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warning/0001/nlin0001031.html | ar5iv | text | # Temporal correlation function in 3-𝐷 Turbulence.
## Abstract
We observe oscillatory decay in the two-point, non-equal time, velocity correlation function of homogeneous, isotropic turbulence. We found this through a direct numerical simulation (DNS) of the three dimensional Navier-Stokes ($`3D`$ NS) equation. We give an approximate analytic theory which explains this oscillatory behaviour. The wave-number and frequency dependent effective viscosity turns out to be complex; the imaginary part gives rise to the temporal oscillation. We find that, at least for the decay at short times, data collapse occur among the inertial range velocity wave-vector modes with the long time dynamic exponent $`z=2/3`$, but the time period of the temporal oscillation is not universal.
In homogeneous, isotropic turbulence the main interest is to understand the long-range spatio-temporal correlations exhibited by the velocity field. Towards this end one studies the scaling behaviour of the velocity structure functions $`S_p(l,t)[v_i(0,0)v_i(l,t)]^p`$ with respect to $`l`$ and $`t`$. Here $`v_i`$ is the velocity field, $`l`$ and $`t`$ are the spatial and temporal separations. Kolmogorov had predicted through his dimensional analysis argument that the equal time structure function $`S_p(l,0)=(ϵl)^{p/3}`$. He had assumed that in the inertial range i.e., $`\eta _dlL`$ ($`\eta _d`$ and $`L`$ are, respectively, dissipation and forcing scales) $`S_p(l,0)`$ is a function of $`ϵ`$ (the mean energy dissipation rate) and $`l`$ only. But after anomalous scaling properties of $`S_p(l)`$ were discovered through experiments and simulations , the importance of $`L`$ (which is also called the integral scale) has been recognised. Now it is known that $`S_p(l)l^{\zeta _p}`$, when $`\zeta _p`$ is a monotonically increasing, convex, nonlinear function of $`p`$. The negative correction to the exponent $`\delta \zeta _p=\zeta _pp/3`$ has to appear as the exponent of a dimensionless quantity $`(l/L)`$ in order to keep the dimension of $`S_p(l)`$ unchanged. In all experiments and simulations of the $`3D`$ NS equation $`L`$ is finite. Infact a grand challenge for analytic theories of turbulence is to show that finite limit for $`\zeta _p`$ exist for $`lL\mathrm{}`$. Such a scheme has been successfully carried out for the passive scalar field advected by a random velocity field in three dimensions (the Kraichnan model ). It has been shown that the anomalous scaling exponents of the passive scalar structure functions are independent of $`L`$ (as $`L\mathrm{}`$), but the amplitudes do depend on $`L`$. Given that so much effort have been made to understand $`S_p(l,0)`$, not much is known about $`S_p(l,t)`$ even for small integer values of $`p`$ (of course $`S_p(l,t)`$ is more complicated than $`S_p(l,0)`$). Literature on $`S_p(l,t)`$ in NS turbulence or related Burgers turbulence is rather sparse . Dynamic renormalisation group (DRG) calculations , one loop self-consistent calculations with the randomly forced $`3D`$ NS equation suggest a long time dynamic exponent $`z=2/3`$. In Ref large wavenumber limit of the velocity correlation function was also explored assuming dynamic scaling hypothesis to be valid. In experiments with a mean flow, because of large scale background velocity, one expects to measure $`z=1`$. But even with zero mean velocity it has been shown , in the context of an $`1D`$ Burgers equation, how $`z=1`$ could arise.
In this work we focus on the simplest two point non-equal-time velocity correlation function in the wave-number ($`k`$) space. We show that in $`3D`$ fluid turbulence, in the large $`L`$ limit, the real part of the non-equal time velocity correlation function $`C_{ij}(𝐤,t)\mathrm{R}v_i(𝐤,t)v_j(𝐤,0)`$, for the inertial scales ($`l=k^1`$), has oscillatory behaviour within a decaying envelope. From the incompressibility ($`.𝐯=0`$) and isotropy assumptions it follows $`C_{ij}(𝐤,t)=c(k,t)P_{ij}(𝐤)`$. Here $`P_{ij}(𝐤)=\delta _{ij}k_ik_j/k^2`$ is the tranverse projector.
An attempt to calculate the $`c(k,t)`$ has been carried out by L’vov et.al. in the context of a turbulent flow with a mean velocity field $`V_0`$. But they treated Navier-Stokes equation at a linear level. They had compensated for the non-linear term to some extent by using a wave number dependent renormalised viscosity instead of the bare viscosity. They predict an oscillatory behaviour for $`c(k,t)`$, but in the absence of the mean velocity (i.e., $`V_0=0`$) the oscillation vanish (i.e., it is a purely kinematic effect). But our data from a numerical simulation of the $`3D`$ NS equation (with large but finite $`L`$ and zero mean velocity) clearly reveals presence of oscillations. The data (see Fig.4,5) looks like the displacement of an under-damped harmonic oscillator. Also simulation of the REWA (reduced wave vector set approximation) model by Eggers shows a non exponential decay at short times and a negative minima.
Incompressible NS equation, forced randomly with a scale dependent variance, has been shown to be a good model for fluid turbulence as far as multiscaling properties are concerned. But there exist many unresolved theoretical problems with the analytic calculations with this model. Our calculation is based on a variant of this model, where instead of a singular forcing spectrum (which goes as $`k^3`$) we use a spectrum which peaks at a small but finite $`L^1`$, goes to zero at $`k=0`$ and behaves as $`k^3`$ for $`kL^1`$. The equation of motion for the velocity field fourier component $`𝐯_𝐢(𝐤)`$ is
$`\dot{v}_i(𝐤)+\nu _0k^2v_i(𝐤)=`$ $`i`$ $`\lambda M_{ijl}(𝐤){\displaystyle \underset{𝐪}{}}v_j(𝐪)v_l(𝐤𝐪)`$ (1)
$`+`$ $`f_i(𝐤,t).`$ (2)
The random force $`f_i(k,t)`$ is a gaussian, white noise with the variancne
$$f_i(𝐤,t)f_j(𝐤^{},t^{})=\frac{(2\pi )^32D_0k}{(k^2+L^2)^2}P_{ij}(𝐤)\delta (𝐤+𝐤^{})\delta (tt^{})$$
(3)
Here $`M_{ijl}=[k_jP_{il}(𝐤)+k_lP_{ij}(𝐤)]/2`$ and $`\lambda `$ is an artificial coupling constant which will be set to $`1`$ later. Henceforth we will denote the variance of the force $`2D_0k/(k^2+L^2)^2`$ by $`D(k)`$.
In the theories with singular forcing spectrum (or equivalently infinite integral scale), infrared divergences appear if one tries to calculate effective viscosity perturbatively. Also the scheme cannot handle the so called sweeping effect i.e., the interaction of the bigger eddies (of size $`q^1`$) with the eddy of size $`k^1`$ (when $`(q<k)`$). Physically the bigger eddies just advect the smaller eddies without distorting them much, so such divergences are basically defect of such a perturbative scheme. But One should remember that this eddy picture is quite heuristic in nature because velocity fourier modes $`v(k)`$ are global features of the velocity field where as the eddies are spatially correlated patches in the velocity field and hence local in nature. One systematic way to get rid of sweeping divergences in the equal time velocity structure functions, is to go to the lagrangian frame. Another way is to do an RG calculation which excludes the effect of the $`(q<k)`$ modes on the $`k`$ mode, so both the infrared divergence and sweeping effect are eliminated. In these calculations because of universality reasons one is mainly interested in the zero frequency limit of the effective viscosity i.e., $`\delta \nu (k,\omega 0)`$ and assumes that for all frequencies $`\omega `$, the effective propagator $`G(k,\omega )(i\omega +k^2\delta \nu (k,\omega 0))^1`$ i.e., remains a Lorentzian in $`\omega `$. This approximation works well for long time properties. But here since we are interested in $`c(k,t)`$ for all $`t`$ (including the short time behaviour) we need the correct behaviour of $`G(k,\omega )`$ for all $`\omega `$s’.
Our procedure to calculate $`c(k,\omega )`$ is a mixture of a self-consistent and a perturbative scheme. We show that if we assume a large but finite integral scale $`L`$, due to nonlinear interaction among modes, the effective viscosity is complex. It has the regular renormalized real part and a $`k,\omega `$ dependent imaginary part as well. The oscillation in $`c(k,t)`$ arises because of this imaginary part. We use an one loop perturbation theory to determine the complex viscosity. The calculation is similar to the standard RG procedure for evaluating zero frequency viscosity. But unlike in the RG procedure, where only modes greater than the external wave-vector $`k`$ are integrated out, we integrate over all $`q`$ modes, including the range $`[0,k]`$. In the zero frequency limit ($`\omega =0`$) our calculation is self-consistent (at one loop level), but for finite $`\omega `$ it is a perturbative calculation. Since our forcing spectrum is not singular, there is no infrared divergence in our integrals.
Treating the nonlinear term perturbatively the effective response function $`G(k,\omega )`$ can be calculated as $`G^1=i\omega +\nu _0k^2+\delta \nu (k,\omega )`$. In the small $`k`$ and $`\nu _00`$ limit, $`\nu _0k^2`$ is negligible compared to $`\delta \nu (k,\omega )`$. Hence $`G^1=i[\omega k^2\mathrm{I}(\mathrm{S})]+k^2\mathrm{R}(\mathrm{S})`$, where we have denoted $`\delta \nu (k,\omega )`$ by $`k^2S`$, and the real, imaginary parts by R,I. Using $`c(k,\omega )=D(k)|G(k,\omega )|^2`$ we get
$$c(k,\omega )=\frac{D(k)}{[\omega k^2\mathrm{I}(\mathrm{S})]^2+k^4\mathrm{R}(\mathrm{S})^2}$$
(4)
$`c(k,t)`$ is the inverse fourier transform of $`c(k,\omega )`$. So our task is to calculate $`S`$. Following Ref.
$`k^2SP_{lj}(𝐤)`$ $`=`$ $`(i\lambda )^2M_{lmn}(𝐤){\displaystyle }{\displaystyle \frac{d^3qd\omega ^{}}{(2\pi )^4}}M_{nij}(𝐤𝐪)\times `$ (6)
$`P_{im}(𝐪)D(q)|G(q,\omega ^{})|^2G(|𝐤𝐪|,\omega \omega ^{})`$
Multiplying both sides by $`P_{jl}(𝐤)`$ and contracting over $`l,j`$ we get
$`S={\displaystyle \frac{M_{jmn}(𝐤)}{2k^2}}`$ $`{\displaystyle }`$ $`{\displaystyle \frac{d^3qd\omega ^{}}{(2\pi )^4}}M_{nij}(𝐤𝐪)P_{im}(𝐪)D(q)\times `$ (9)
$`|G(q,\omega ^{})|^2G(|𝐤𝐪|,\omega \omega ^{})`$
Integrating the r.h.s. over $`\omega ^{}`$ gives
$`S={\displaystyle \frac{d^3q}{(2\pi )^3}b(𝐤,𝐤𝐪,𝐪)\frac{D(q)}{2\nu _1q^z}\frac{1}{i\omega +\nu _1(q^z+|𝐤𝐪|^z)}}`$ (10)
$`k^2b(𝐤,𝐤𝐪,𝐪)`$ is obtained by contracting $`M_{jmn},M_{nij}`$ and $`P_{im}`$. The expression for $`b`$ is $`b(𝐤,𝐤𝐪,𝐪)=\frac{|𝐤𝐪|}{k}(xy+z^3)`$ . The trio $`(𝐤,𝐤𝐪,𝐪)`$ form a triangle and $`x,y,z`$ are the direction cosines of the angles opposite to $`𝐤,𝐤𝐪`$, and $`𝐪`$ respectively. We have calculated this integral numerically. We have used $`G(k,\omega )^1=i\omega +\nu _1k^z`$ anticipating a renormalisation of the viscosity in the $`\omega 0`$ limit. Here $`\nu _1`$ could be a function of $`L,D_0`$ which we determine later. But as far as the $`k,\omega `$ dependence of the integral is concerned, for fixed $`L`$ and $`D_0`$, we can treat $`\nu _1`$ as a constant. From Eq.5 note that $`\mathrm{I}(\mathrm{S})0`$, as $`\omega 0`$. We check numerically that at $`\omega =0`$, the integral in Eq.5 scales as $`k^{2z}`$. Also by expanding the integral in small $`k`$ and retaining only leading powers in $`k`$ (as is done in a DRG procedure ) one can infer this. If self-consistency has to be achieved at $`\omega =0`$, then on the l.h.s. of Eq.5, $`R[S(k,0)]=\nu _1k^{z2}`$. This fixes $`z=2/3`$. Numerical evaluation of the integral with a finite $`L`$ corroborates the approximate analytic prediction for small $`k`$ because $`L^1<<k`$. Now we determine $`\nu _1`$ in a self consistent way .
At $`\omega =0`$, self consistency of Eq.5 requires
$`\nu _1k^{4/3}={\displaystyle \frac{D_0}{(2\pi )^2\nu _1^2}}k^{4/3}g(L),\text{so}\nu _1^3={\displaystyle \frac{D_0}{(2\pi )^2}}g(L)`$ (11)
Where $`g(L)`$ is the $`L`$ dependent integral on the r.h.s. of Eq.5. In Fig.1 we show that $`g(L)`$ converges to a finite value as $`L`$ grows large. Using $`g(L).95`$ from Fig.1, we get $`\nu _1(L)0.28D_0^{1/3}`$. While evaluating the integral for inertial range $`k`$ modes we used $`G^1=i\omega +\nu _1k^{2/3}`$ (neglecting the $`\nu _0k^2`$ term) for all the $`q`$ and $`kq`$ modes in the integrand, though the integral runs over both inertial and dissipation modes. But this is approximately correct because modes very far from the $`k`$ mode do not contribute (locality in $`k`$ space).
Returning to Eq.5, the dominant pole of $`c(k,\omega )`$ will decide the oscillation and the decay of $`c(k,t)`$. If the pole lies at $`\omega =\omega _1+i\omega _2`$ then $`\omega _1=k^2\mathrm{I}(\mathrm{S}[\omega _1])`$ and $`\omega _2=k^2\mathrm{R}(\mathrm{S}[\omega _2])`$. From the shapes of $`\mathrm{R},\mathrm{I}(\mathrm{S})`$ in Fig.2 and the fact that $`\mathrm{R},\mathrm{I}(\mathrm{S})`$ are even,odd functions of $`\omega `$, we can infer that there will be two solutions to these transcendental equations. They are of the form $`\pm \omega _1+i\omega _2`$, when $`\omega _2`$ is negative. This ensures the causality of the effective response function $`G`$. The approximate data collapse for the inertial range $`k`$ modes ($`\eta _D^1>k>L^1`$) in our Fig.3 implies that the oscillation period scales with $`k^{2/3}`$.
We performed a DNS of the $`3D`$ NS equation and calculated $`c(k,t)`$ versus $`t`$ for various $`𝐯(𝐤)`$ modes. The data (see Fig.4,5) clearly shows oscillation in time. In our simulation forcing was present only at large length scales (i.e., in the $`k`$-space the $`𝐯(𝐤)`$ modes in the smallest two shells were forced). Our pseudo-spectral scheme for the DNS is same as in . We used a $`32^3`$ grid with periodic boundary condition. We obtained a short inertial range (shown in Fig.6) with $`Re_\lambda 22`$. We have obtained very long time series ($`68T(L)`$) for averaging $`c(k,t)`$. A preliminary test run on a $`64^3`$ grid also confirms existence of such oscillations.
Our simulation data in Fig.4 indicates that the oscillation time period approximately scales as $`k^{2/3}`$ in the inertial range. But Fig.5 shows that modes closed to the forced modes and the dissipative modes differ widely from the inertial ones in periodicity and decay. We have averaged the data for $`68T(L)`$ (when $`T(L)`$ is the large eddy turnover time). This is sufficienty long averaging time for an equal time, shell averaged, correlation functions to converge; but here for $`c(k,t)`$ since (a) single $`𝐤`$ mode is involved, (b) long time history (few $`T(L)`$’s) of the mode is important, it requires a longer averaging time. In a simulation with small but finite $`\nu _0`$ the $`\nu _1k^{2/3}`$ term will loose its dominance over $`\nu _0k^2`$ as $`k`$ increases towards the dissipation range. That explains the strong damping seen in the $`𝐤=(0,5,6)`$ mode in Fig.5.
In the simulation we calculate $`T(L)`$ using the formula $`T(L)=L_{box}/v_{rms}`$, where $`L_{box}`$ is the simulation box size. From dimensional analysis arguments the scale dependent eddy-turnover time $`T(k)=Aϵ^{1/3}k^{2/3}`$, when $`A`$ is a constant of $`𝒪(1)`$. Equating $`T(L)=Aϵ^{1/3}(2\pi /L)^{2/3}`$ we get the prefactor $`Aϵ^{1/3}`$ and hence can calculate $`T(k)`$. Fig.4 shows that the time period of oscillation $`\lambda _k25T(k)`$. This ratio may be non-universal. To get a clue let us look at our randomly forced model, where $`I(S)`$ is an explicit function of $`D_0`$, $`\nu _1`$ and $`\nu _1D_0^{1/3}`$ (see Eq.5 and 6). Hence $`\lambda _k`$ depends on $`D_0`$ in a complicated way. We explore the dependence of $`T(k)`$ on $`D_0`$ below. In this model $`T(k)`$ cannot be determined in a simple way because here all the $`k`$ shells are being forced and hence the energy flux is not a constant but increses logarithmically with $`k`$ . The scale dependent energy flux $`\mathrm{\Pi }(k)=_0^kD(q)d^3q/(2\pi )^3=\frac{D_0}{2\pi ^2}[\mathrm{ln}(a^2+k^2)k^2(a^2+k^2)^12\mathrm{ln}(a)]`$, when $`a=L^1`$. Again using dimensional analysis we get $`T(k)=A.\mathrm{\Pi }(k)^{1/3}k^{2/3}`$ (hence $`T(k)D_0^{1/3}`$). Evaluating $`\mathrm{\Pi }(k)^{1/3}`$ for $`D_0=0.01`$ and $`L=100`$ (which we used for our theoretical graphs in Fig.2,3) gives $`\mathrm{\Pi }(0.1)^{1/3}=6.17`$, $`\mathrm{\Pi }(1.0)^{1/3}=8.1`$. Allthough this $`k`$ dependence is weak, but the perfect data collapse with respect to $`tk^{2/3}`$ in Fig.3 implies that the above dimensional analysis estimate is not accurate. In Fig.3 the oscillation period $`\lambda _k6k^{2/3}`$ and hence $`\lambda _kT(k)`$ (neglecting the constant $`A`$ of $`𝒪(1)`$).
The imaginary part of the viscosity, generated by an one loop perturbation theory, cannot be interpreted as a background velocity which is slowly varying in time. This mis-interpretation may be provoked by the fact that if we had a mean background flow $`V_0`$ in the problem, then the nonlinear term would generate an extra term $`i(𝐕_\mathrm{𝟎}.𝐤)𝐯(𝐤,t)`$ in the NS equation. With this extra linear term the bare propagator will be $`G(k,\omega )(i\omega i𝐕_\mathrm{𝟎}.𝐤+\nu _0k^2)^1`$ . But the $`k,\omega `$ dependent complex viscosity, which we find in our theory, when transformed to $`(𝐱,t)`$ space, gives an additional memory term $`𝑑𝐱^{}𝑑t^{}\kappa (|𝐱𝐱^{}|,tt^{})𝐯(𝐱^{},t^{})`$ in the equation of motion (e.o.m.). But it is true that the $`\kappa (k,t)`$ field oscillates at a slower time scale than that of the $`𝐯(𝐤,t)`$ mode itself (as $`\mathrm{I}\delta \nu (k,\omega )`$ has a peak at a lower $`\omega `$ than that of $`c(k,\omega )`$). Also a significant contribution to the integral for $`\delta \nu (k,\omega )`$ comes from the $`q<k`$ modes. So it does resemble sweeping by larger (and hence slower) eddies to some extent. But rigorously convection by a background velocity $`𝐕_\mathrm{𝟎}(𝐱,t)`$ should look like $`𝐕_\mathrm{𝟎}(𝐱,t).𝐯(𝐱,t)`$ (which is local in $`(𝐱,t)`$). We note that in the turbulence context, complex effective viscosity has been proposed before by J.K. Bhattacharjee in . The author had assumed dynamic scaling hypothesis (as in dynamic critical phenomena) to be valid for fluid turbulence and had predicted a form for the effective viscosity in the high frequency limit. Then an interpolation scheme had been used to connect the two limits (small $`\omega `$ and large $`\omega `$) of the effective viscosity.
In conclusion, we have given numerical evidence and an approximate theory for the novel oscillatory decay of the two-point, temporal correlation function in $`3D`$ fluid turbulence. This behaviour is similar to viscoelastic effect seen in complex fluids.
We would like to thank J.K. Bhattacharjee (of IACS Calcutta, India) for some critical comments and R. Pandit, S. Ramaswamy (of IISc Bangalore, India) and M. Plischke (of SFU, Canada) for discussions. Also thanks to SERC,IISc Bangalore, (India) for computational facilities and NSERC (Canada) for supporting work in Canada. |
warning/0001/hep-ph0001156.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In theories with non-trivial structure of vacua a number of interesting physical effects, induced by instanton solutions, appear. In the present article we will study shadow processes . These are non-perturbative processes in which both the initial and the final state are in the false vacuum. Apart from standard perturbative contributions, the processes which start and end in the false vacuum acquire additional contributions due to the underbarrier tunelling of the system to another vacuum and its return to the initial one. This transition is obviously induced by an instanton solution and goes through the intermediate state containing a bubble of the true vacuum. We would like to mention that other examples of instanton induced processes are transitions with baryon number violation between the vacua in the electroweak theory and the decay of a metastable (false) vacuum due to underbarrier tunelling from a false vacuum to the true one .
Much work has been done to study the instanton induced transitions, and quite effective techniques for the calculation of the probabilities of such trunsitions have been developed (see Refs. for a review). We will study the instanton contribution to the total cross section $`\sigma _2(E)`$ of a process ($`2`$ any) with two initial particles of the total energy $`E`$ in the $`(\lambda \varphi ^4)`$-theory. There is a number of arguments showing that $`\sigma _2(E)`$ can be presented in the following exponential form:
$$\sigma _2(E)e^{\frac{1}{\lambda }F(ϵ)+𝒪(1)},$$
(1)
where $`\lambda `$ is the coupling constant in the model, $`ϵ=E/E_{sph}`$ and $`E_{sph}`$ is the energy of the sphaleron configuration which characterizes the height of the barrier separating the vacua. The leading order approximation of the function $`F(ϵ)`$ for small $`ϵ`$ was studied in Refs. . The next-to-leading term is a propagator correction for it includes contributions from the propagator in the instanton background. Hence, calculation of the next-to-leading correction requires knowledge of the instanton propagator. It turns out that in the $`(\lambda \varphi ^4)`$-theory an exact expression for the instanton propagator can be obtained. Calculation and discussion of the propagator correction to the function $`F(ϵ)`$ is one of the purposes of this article.
An important issue is that of the validity of formula (1). In the electroweak theory a proof based on the properties of the propagator in the instanton background was given in Ref.. We apply the arguments of Ref. in the $`(\lambda \varphi ^4)`$-theory making use of the explicit expression for the propagator.
According to arguments of Refs. for the multiparticle initial state the total cross section is semiclassical and has the form
$$\sigma _N(E)e^{\frac{1}{\lambda }F(ϵ,\nu )+𝒪(1)},$$
(2)
where $`N`$ is the number of initial particles, $`\nu =N/N_{sph}`$, and $`N_{sph}1/\lambda `$ is a characteristic number of particles contained in the sphaleron. Note that in the regime $`\lambda 0`$ and $`\nu `$ fixed $`N\nu /\lambda `$ is a large number. The function $`F(ϵ,\nu )`$ for the $`(\lambda \varphi ^4)`$-theory was calculated numerically for a certain range of $`ϵ`$ and $`\nu `$ in Ref. . In Refs. it was argued that the leading exponential term of the two-particle cross section can be calculated from the following formula:
$$\underset{\lambda 0}{lim}\lambda \mathrm{ln}\sigma _2=\underset{\nu 0}{lim}F(\frac{E}{E_{sph}},\nu ).$$
(3)
In this conjecture it is assumed that the limit $`\nu 0`$ exists. The problem is that the function $`F(E/E_{sph},\nu )`$ is known to contain contributions singular in $`\nu `$. In particular, in the $`(\lambda \varphi ^4)`$-theory such contributions already appear in the propagator correction. The conjecture basically claims that terms singular in $`\nu `$ cancel each other in the final answer. Its validity, of course, means that the semiclassical form of the two-particle cross section is indeed given by Eq. (1) with $`F(E/E_{sph})=F(E/E_{sph},0)`$. Verification of conjecture (3) in the next-to-leading order is another purpose of this paper. Note that different arguments in favor of this conjecture were given in Refs..
The plan of the article is the following. In Sect. 2 we describe the model and discuss the propagator in the instanton background. Namely, we present the high energy asymptotics of the propagator and discuss the implementation of Mueller’s idea in the scalar model. We also discuss the exact expression of the double residue of the instanton propagator. In Sect. 3 we apply it for the evaluation of the next-to-leading order (propagator correction) of the function $`F(ϵ,\nu )`$. There we explicitly demonstrate the appearance of terms singular in $`\nu `$ for $`\nu 0`$ and their cancellation in the final result. Sect. 4 contains some discussion of the results. In particular, the range of validity of the next-to-leading order approximation is estimated.
## 2 Instanton propagator in the scalar model
We consider the model of one component real scalar field, defined by the Minkowskian action
$$S=d^4x\left[\frac{1}{2}\left(_\mu \varphi \right)^2\frac{1}{2}m^2\varphi ^2+\frac{\lambda }{4!}\varphi ^4\right],$$
(4)
where $`\lambda >0`$. The potential of the model is unbounded from below, hence the minimum $`\varphi =0`$ is metastable. Underbarrier tunelling from this vacuum to the instability region and its return to the trivial vacuum is the transition which gives rise to the shadow process we are going to study here.
Let us consider first the case $`m=0`$. There is a well known instanton solution in the massless theory given by the formula
$$\varphi _{inst}(x;x_0,\rho )=\frac{4\sqrt{3}}{\sqrt{\lambda }}\frac{\rho }{(xx_0)^2+\rho ^2}.$$
(5)
Here $`x_{0\mu }`$ is the center of the instanton and $`\rho `$ is its size. Due to the conformal invariance of the massless theory the action of the instanton does not depend on $`\rho `$,
$$S_{inst}^{(0)}S(\varphi _{inst})=\frac{16\pi ^2}{\lambda }.$$
(6)
In the case $`m0`$ the mass term breaks the conformal invariance. Using standard scaling arguments it can be shown that there are no regular solutions of the Euclidean equations of motion with finite action. The decay of the vacuum $`\varphi =0`$ is dominated by the constrained instanton, a configuration which can be regarded as an approximate solution of the equations of motion. It minimizes the action under the constraint that the size of the configuration is $`\rho `$. A formalism for construction of such configurations and evaluation of the functional integral was developed in .
When $`\rho m1`$ the constrained instanton configuration behaves like the instanton (5) of the massless theory at $`x\rho `$ and as a solution of the free massive theory for $`x>m^1`$. The action of such configuration is
$$S_{inst}(\rho )=\frac{16\pi ^2}{\lambda }\frac{24\pi ^2}{\lambda }(\rho m)^2\left[\mathrm{ln}\frac{\rho ^2m^2}{4}+2C_E+1\right]+𝒪(\rho ^4m^4),$$
(7)
where $`C_E=0.577\mathrm{}`$ is the Euler constant. For the class of constraints mentioned above the terms given in (7) do not depend on the explicit form of the constraint, whereas the corection $`𝒪(\rho ^4m^4)`$ does. In our analysis we limit ourselves to the constraint independent order of the approximation.
For $`m0`$ the potential barrier separating the trivial vacuum $`\varphi =0`$ from the instability region is finite. Its height is characterized by a sphaleron solution, a static $`SO(3)`$-symmetric configuration satisfying the equation of motion. In Ref. it was found that the sphaleron energy and the sphaleron number of particles are
$$E_{sph}=\kappa \frac{m}{\lambda },\kappa =113.4\text{and}N_{sph}=63\frac{1}{\lambda }$$
(8)
respectively.
Now let us study the propagator in the instanton background. It is defined by the operator $`\widehat{D}_x`$ of quadratic fluctuations appearing in the expansion of the action around the instanton solution. In the massless case this operator is equal to
$$\widehat{D}_x=\frac{^2}{x_\mu ^2}+\frac{\lambda }{2}\varphi _{inst}^2(x;0,\rho )=\frac{^2}{x_\mu ^2}+\frac{24\rho ^2}{(\rho ^2+x^2)^2}.$$
It can be easily seen that it possesses five zero modes $`\psi _A(x)`$ $`(A=1,2,3,4,5)`$ corresponding to the translational invariance and the scale invariance of the massless theory. The zero modes can be obtained by differentiation of the instanton solution with respect to the parameters $`x_0`$ and $`\rho `$:
$$\psi _A\frac{}{\zeta _A}\varphi _{inst}(x;x_0,\rho )|_{x_0=0},\zeta _\mu =(x_0)_\mu ,\zeta _5=\rho .$$
(9)
Because of the existence of the zero modes there is an ambiguity in the definition of the propagator that can be fixed by imposing additional constraints. Let $`G_f(x,y)`$ be the inverse of $`\widehat{D}_x`$ on the subspace of functions orthogonal to some functions $`f_A(x)`$. The latter satisfy the only condition that the matrix
$$\mathrm{\Omega }_{AB}=𝑑x\psi _A(x)f_B(x)$$
is invertible . According to this definition the instanton propagator satisfies the equation
$$\widehat{D}_xG_f(x,y)=\delta (xy)\underset{A}{}f_A(x)\mathrm{\Omega }_{AB}^1\psi _B(x)$$
(10)
and the orthogonality constraints
$$𝑑xf_A(x)G_f(x,y)=0=𝑑yG_f(x,y)f_B(y).$$
(11)
The r.h.s. of Eq. (10) is the projector onto the subspace orthogonal to the functions $`f_A(x)`$. Physical results, of course, do not depend on particular choice of the functions $`f_A(x)`$. Below, following the ideas of Refs. , we will use the freedom of choosing constraints (11) to eliminate the leading asymptotics of the instanton propagator and simplify the analysis of the initial-state corrections.
For a particularly simple and natural choice of the functions $`f_A(x)`$, namely for
$$f_A(x)=w(x)\psi _A(x),$$
(12)
where the weight function
$$w(x)=\frac{4\rho ^2}{(\rho ^2+x^2)^2},$$
the propagator in the instanton background was calculated explicitly in Ref. (see also ). Allowing some abuse of notation we denote this propagator by $`G_\psi (x,y)`$. It is equal to
$`G_\psi (x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \frac{\rho ^2}{(\rho ^2+x^2)(\rho ^2+y^2)}}\{{\displaystyle \frac{1}{2d(x,y)}}3\mathrm{ln}d(x,y)`$ (13)
$``$ $`{\displaystyle \frac{43}{5}}+6d(x,y)\mathrm{ln}d(x,y)+{\displaystyle \frac{56}{5}}d(x,y)\},`$
where
$$d(x,y)=\frac{\rho ^2(xy)^2}{(\rho ^2+x^2)(\rho ^2+y^2)}.$$
The relation between $`G_\psi (x,y)`$ and the propagator $`G_f(x,y)`$ for an arbitrary constraint (11) is given by the following formula:
$`G_f(x,y)`$ $`=`$ $`G_\psi (x,y)\left({\displaystyle 𝑑zG_\psi (x,z)f_A(z)}\right)\mathrm{\Omega }_{AB}^1\psi _B(y)`$ (14)
$``$ $`\psi _A(x)\left(\mathrm{\Omega }_{AB}^T\right)^1{\displaystyle 𝑑zf_B(z)G\psi (z,y)}`$
$`+`$ $`\psi _A(x)\left(\mathrm{\Omega }_{AB}^T\right)^1\left({\displaystyle 𝑑z𝑑z^{}f_B(z)G\psi (z,z^{})f_C(z^{})}\right)\mathrm{\Omega }_{CD}^1\psi _D(y).`$
The Fourier transform of the instanton propagator is defined in the standard way:
$$G_f(p,q)=𝑑x𝑑ye^{ipx+iqy}G_f(x,y).$$
In principle, using the exact result, Eq. (13), the function $`G_\psi (p,q)`$ can be obtained by direct calculation. We did not find the complete expression. Instead we derived the asymptotic formula for the Fourier transform of the instanton propagator in the regime when $`p^2`$, $`q^2`$ are fixed and $`s(p+q)^2\mathrm{}`$. The growing terms of the asymptotics are given by
$$G_\psi (p,q)=\frac{16\pi ^2}{p^2q^2}\left[s\rho ^2\mathrm{ln}(s\rho ^2)\mathrm{\Pi }_1(p,q)+(s\rho ^2)\mathrm{\Pi }_2(p,q)+\mathrm{ln}(s\rho ^2)\mathrm{\Pi }_3(p,q)+\mathrm{}\right],$$
(15)
where
$`\mathrm{\Pi }_1(p,q)`$ $`=`$ $`{\displaystyle \frac{3}{4}}𝒮_1(p\rho )𝒮_1(q\rho ),`$ (16)
$`\mathrm{\Pi }_2(p,q)`$ $`=`$ $`{\displaystyle \frac{3}{2}}\left(C_E{\displaystyle \frac{1}{15}}\mathrm{ln}2\right)𝒮_1(p\rho )𝒮_1(q\rho ),`$ (17)
$`\mathrm{\Pi }_3(p,q)`$ $`=`$ $`\{𝒮_1(p\rho )[{\displaystyle \frac{9}{2}}𝒮_2(q\rho )({\displaystyle \frac{27}{4}}+{\displaystyle \frac{3}{4}}q^2\rho ^2)𝒮_1(q\rho )]`$ (18)
$`+`$ $`[{\displaystyle \frac{9}{2}}𝒮_2(p\rho )({\displaystyle \frac{27}{4}}+{\displaystyle \frac{3}{4}}p^2\rho ^2)𝒮_1(p\rho )]𝒮_1(q\rho ){\displaystyle \frac{3}{2}}𝒮_2(p\rho )𝒮_2(q\rho )\}.`$
Here $`𝒮_n(z)`$ is defined by $`𝒮_n(z)=z^nK_n(z)`$, where $`K_n(z)`$ is the modified Bessel function. Using the explicit expressions for the translational zero modes (see Eq. (9)) normalized with respect to the weight function $`w(x)`$, the first two terms of the asymptotic expansion (15) can be written as
$$G_\psi (p,q)=\frac{1}{5\rho ^2}\mathrm{ln}(\rho ^2s)\psi _\mu (p)\psi _\mu (q)\frac{2}{5\rho ^2}\left(C_E\frac{1}{15}\mathrm{ln}2\right)\psi _\mu (p)\psi _\mu (q)+\mathrm{}$$
(19)
The leading term of the asymptotics of the propagator in the instanton background was calculated in Ref. . This result is in complete agreement with the first term in Eq. (19).
In Ref. Mueller proposed an idea to use the ambiguity in the choice of the function $`f_A(x)`$ in order to cancel the two leading terms in the asymptotics of the propagator $`G_\psi (p,q)`$. Then the propagator contribution, as well as loop contributions of the initial state corrections disappear. As a consequence, such corrections do not exponentiate, i.e., do not give contributions to the function $`F(ϵ)`$. In addition, in this case the initial-final state corrections can be described semiclassically. Namely, the effect of the initial state lines can be taken into account by substituting the instanton by a new field configuration which is a solution to the classical equation of motion with an external source (see Ref. for details).
Now we explain how the functions $`f_A(x)`$ can be chosen to provide vanishing of the two leading terms of the asymptotics of $`G_f(p,q)`$. For this we repeat the arguments of Ref. . It turns out that for such functions the corresponding propagator constraint (11) is not relativistically covariant. Let $`p_1`$ and $`p_2`$ be the arguments of the Fourier transform of the propagator. We choose a coordinate system such that $`p_{1j}=p_{2j}=0`$ for $`j=2,3`$, whereas $`p_{1+}\rho =p_2\rho 1`$ and $`p_1^2`$ and $`p_2^2`$ are fixed. Here the $`\pm `$ components of the momenta are defined by
$$p_{j\pm }=\frac{(p_j)_0\pm (p_j)_1}{\sqrt{2}}.$$
Then $`(p_1,p_2)\rho ^2p_{1+}p_2\rho ^21`$. Only the components $`f_\mu (x)`$, corresponding to translations, modify the leading asymptotics of the propagator. The Fourier transforms $`\stackrel{~}{f}_\mu (p)`$ of the functions $`f_\mu (x)`$, defining the required propagator constraint, are chosen in the following way:
$$\stackrel{~}{f}_\mu (p)=\delta (p_+M)\delta (p_{}+M)\overline{f}_\mu (p_2,p_3),$$
(20)
where $`M`$ is an arbitrary parameter of the dimension of mass. Substituting these functions into Eq. (14) one finds after some calculations that
$$G_f(p_1,p_2)=\psi _\mu (p_1)𝒢_{\mu \nu }\psi _\nu (p_2),$$
(21)
where the $`4\times 4`$ constant matrix $`𝒢_{\mu \nu }`$ is equal to
$$𝒢_{\mu \nu }=\frac{1}{5\rho ^2}\mathrm{\Omega }_{\mu \sigma }^1\frac{1}{(2\pi )^8}d^4q_1d^4q_2f_\sigma (q_1)\psi _\rho (q_1)\mathrm{ln}\frac{(q_1,q_2)}{M^2}\psi _\rho (q_2)f_\tau (q_2)\left(\mathrm{\Omega }^T\right)_{\tau \nu }^1.$$
Using the freedom of choosing the functions $`\overline{f}_\mu `$ one can make the constant real symmetric matrix $`𝒢_{\mu \nu }`$ equal to zero. We would like to stress that the knowledge of the exact formulas for the leading terms of the asymptotics of $`G_\psi (p,q)`$, Eqs. (15) - (19), allows us to get the explicit expression of the matrix $`𝒢_{\mu \nu }`$. This is in contrast with the case of the electroweak theory where only a general structure of the analogous matrix can be derived .
For the perturbative calculations of the function $`F(ϵ,\nu )`$ the on-mass-shell residue of the instanton solution will be needed. By definition it is equal to
$$R_{inst}(𝐤)=(k^2+m^2)\stackrel{~}{\varphi }_{inst}(k;0,\rho )|_{k_0=i\omega _𝐤},$$
(22)
where $`\stackrel{~}{\varphi }_{inst}(k;x_0,\rho )`$ is the Fourier transform of the instanton,
$$\stackrel{~}{\varphi }_{inst}(k;x_0,\rho )=d^4xe^{ikx}\varphi _{inst}(x;x_0,\rho )$$
and $`\omega _𝐤=\sqrt{𝐤^2+m^2}`$. For the instanton solution (5) in the massless theory
$$R_{inst}=\frac{1}{\sqrt{\lambda }}16\sqrt{3}\pi ^2\rho .$$
(23)
Correspondingly, to calculate of the next-to-leading correction to the function $`F(ϵ,\nu )`$ we need the expressions for the double on-mass-shell residues of the instanton propagator. We will consider the propagator orthogonal to functions (12). Let us introduce the following notations:
$`R_{aa}(𝐤,𝐪)`$ $`=`$ $`(k^2+m^2)(q^2+m^2)G_\psi (k_0,𝐤;q_0,𝐪)|_{k_0=i\omega _𝐤,q_0=i\omega _𝐪},`$ (24)
$`R_{ab}(𝐤,𝐪)`$ $`=`$ $`(k^2+m^2)(q^2+m^2)G_\psi (k_0,𝐤;q_0,𝐪)|_{k_0=i\omega _𝐤,q_0=i\omega _𝐪},`$ (25)
$`R_{bb}(𝐤,𝐪)`$ $`=`$ $`(k^2+m^2)(q^2+m^2)G_\psi (k_0,𝐤;q_0,𝐪)|_{k_0=i\omega _𝐤,q_0=i\omega _𝐪},`$ (26)
The indices $`a`$ and $`b`$ correspond to initial and final particles, respectively (in the notations of Ref. ). For the scalar massless field $`\omega _𝐤=|𝐤|`$ and all three residues (24)-(26) can be expressed in terms of one function:
$$R_\mathrm{\#}(𝐤,𝐪)=\rho ^2R\left(\rho ^2s_\mathrm{\#}^{(0)}(𝐤,𝐪)\right),$$
(27)
where $`\mathrm{\#}=aa`$, $`ab`$, $`bb`$ and the function $`s_\mathrm{\#}^{(0)}(𝐤,𝐪)`$ is the $`s`$-variable for the corresponding particles on the mass shell,
$$s_{aa}^{(0)}(𝐤,𝐪)=s_{bb}^{(0)}(𝐤,𝐪)=s_{ab}^{(0)}(𝐤,𝐪)=2(|𝐤||𝐪|\mathrm{𝐤𝐪}).$$
However, in the calculation of the next-to-leading order corrections due to non-zero mass must be taken into account. It turns out that within the accuracy set by Eq. (7) it is enough to consider the residues defined through the relation
$$R_\mathrm{\#}(𝐤,𝐪)=\rho ^2R\left(\rho ^2s_\mathrm{\#}(𝐤,𝐪)\right),$$
(28)
(cf. (27)), where the function is calculated for the instanton propagator of the massless theory, whereas the $`s`$-variable is taken for the massive one:
$`s_{aa}(𝐤,𝐪)`$ $`=`$ $`s_{bb}(𝐤,𝐪)=2m^22(\omega _𝐤\omega _𝐪\mathrm{𝐤𝐪}),`$ (29)
$`s_{ab}(𝐤,𝐪)`$ $`=`$ $`2m^2+2(\omega _𝐤\omega _𝐪\mathrm{𝐤𝐪}).`$ (30)
The consistency of this procedure is discussed in Sect. 4.
The exact expression for the function $`R(s)`$ was obtained in Ref. and is given by
$`R(s)`$ $`=`$ $`16\pi ^2\left\{\alpha _1\left[s\mathrm{ln}{\displaystyle \frac{s}{4}}+2\left(C_E{\displaystyle \frac{1}{15}}\right)s\right]\alpha _2\left[\mathrm{ln}{\displaystyle \frac{s}{4}}+2\left(C_E+{\displaystyle \frac{43}{30}}\right)\right]\right\},`$ (31)
$`\alpha _1`$ $`=`$ $`3/4,\alpha _2=3/2.`$
In the next section this result will be used for the calculation of the next-to-leading correction to the function $`F(ϵ,\nu )`$.
## 3 Multiparticle cross section
Formula (2) for the multiparticle cross-section of shadow processes comes from the following expression derived in Ref. ,
$`\sigma _N(E)`$ $``$ $`{\displaystyle }d^4x_0d\rho d^4\xi d\theta \mathrm{exp}[2S_{inst}(\rho )+{\displaystyle \frac{1}{\lambda }}W^{(1)}(x_0,\rho ,\xi ,\theta )`$ (32)
$`+`$ $`{\displaystyle \frac{1}{\lambda }}W^{(2)}(x_0,\rho ,\xi ,\theta )+\mathrm{}],`$
where we integrate over the position $`x_0`$ and the size $`\rho `$ of the instanton, as well as over auxiliary variables $`\xi _\mu `$ and $`\theta `$. We also indicated explicitly the dependence of the action on the size of the instanton (see Eq. (7)). The terms $`W^{(i)}`$ account for fluctuations in the instanton background: $`W^{(1)}`$ corresponds to leading diagrams without propagator lines, $`W^{(2)}`$ corresponds to diagrams with one internal propagator in the instanton background, etc. Diagrams with loops do not appear in the $`(1/\lambda )`$ order of the semi-classical approximation, they contribute to $`𝒪(1)`$ terms in (2).
General expressions for the functions $`W^{(1)}`$ and $`W^{(2)}`$ were derived in . The integrals in Eq. (32) are evaluated by the saddle point method. It can be checked that up to the next-to-leading order the saddle point values of $`x_0`$, $`\rho `$, $`\xi `$ and $`\theta `$ are determined by the leading-order equations. These equations are obtained by differentiation of the expression $`(2S_{inst}(\rho )+W^{(1)}/\lambda )`$ with respect to $`x_0`$, $`\rho `$, $`\xi `$ and $`\theta `$. The physically relevant saddle point has $`(x_0)_i=0`$, $`\xi _i=0`$ ($`i=1,2,3`$), while $`(x_0)_0`$, $`\xi _0`$ and $`\theta `$ are purely imaginary. It is convenient to introduce the following notations: $`x_0=i\tau `$, $`\xi _0=i\chi `$, and $`\theta =i\mathrm{ln}\gamma `$. In accordance with Eq. (32) the function $`F(ϵ,\nu )`$ is represented as
$$F(ϵ,\nu )=32\pi ^2+F^{(1)}(ϵ,\nu )+F^{(2)}(ϵ,\nu )+\mathrm{}.$$
(33)
The first term in the r.h.s. is just $`(2\lambda S_{inst}^{(0)})`$, where $`S_{inst}^{(0)}`$ is the instanton action in the massless theory, Eq. (6). The non-trivial leading order correction $`F^{(1)}(ϵ,\nu )`$ corresponds to the contribution of
$$2\lambda (S_{inst}^{(0)}S_{inst}(\rho ))+W^{(1)}$$
in Eq. (32). The next-to-leading (propagator) correction $`F^{(2)}(ϵ,\nu )`$ is given by $`W^{(2)}`$ evaluated at the saddle point solution.
### 3.1 Leading order correction
For general values of $`ϵ`$ and $`\nu `$ the system of saddle point equations is too complicated and we studied it numerically. The results are described at the end of the section.
In the limit of small $`\nu `$ the calculations simplify considerably. Keeping only relevant terms we obtain that the function $`W^{(1)}`$ of Eq. (32) reads
$`{\displaystyle \frac{1}{\lambda }}W^{(1)}(\tau ,\rho ,\chi ,\gamma )=E\chi N\mathrm{ln}\gamma +{\displaystyle \frac{1}{(2\pi )^3}}{\displaystyle \frac{d𝐩}{2\omega _𝐩}R_{inst}(𝐩)e^{\omega _𝐩\tau }R_{inst}(𝐩)}`$
$`+{\displaystyle \frac{\gamma }{(2\pi )^3}}{\displaystyle \frac{d𝐩}{2\omega _𝐩}R_{inst}(𝐩)e^{\omega _𝐩(\chi \tau )}R_{inst}(𝐩)}+\mathrm{}`$
$`=E\chi N\mathrm{ln}\gamma +192\pi ^2{\displaystyle \frac{\rho ^2m^2}{\lambda }}\left[\mathrm{\Phi }(m\tau )+{\displaystyle \frac{\gamma }{m^2(\chi \tau )^2}}\right]+\mathrm{},`$ (34)
where $`R_{inst}`$ is given by expression (23), $`\omega _𝐩=\sqrt{𝐩^2+m^2}`$, and
$$\mathrm{\Phi }(z)\frac{1}{z}K_1(z).$$
We would like to stress that in these calculations the expression for the energy $`\omega _𝐩`$ of the massive theory is used, whereas it is enough to substitute the residue $`R_{inst}`$ of the instanton solution of the massless theory. This is consistent with the approximation we are considering in the present paper. The question of validity of this procedure is discussed in Sect. 4.
For further calculations it is convenient to introduce the variables $`\stackrel{~}{ϵ}=E\lambda /m`$ and $`\stackrel{~}{\nu }=N\lambda `$. From Eqs. (8) it follows that $`\stackrel{~}{ϵ}=\kappa ϵ`$ and $`\stackrel{~}{\nu }=63\nu `$. To the leading order in $`\nu `$ the saddle point solution can be written in the form
$$\stackrel{~}{\rho }^2=\frac{1}{192\pi ^2m^2}\frac{\stackrel{~}{ϵ}}{\mathrm{\Phi }^{}(m\stackrel{~}{\tau })};\stackrel{~}{\gamma }=4\left(\frac{\stackrel{~}{\nu }}{\stackrel{~}{ϵ}}\right)^3\mathrm{\Phi }^{}(m\stackrel{~}{\tau });\stackrel{~}{\chi }=\stackrel{~}{\tau }+\frac{2}{m}\frac{\stackrel{~}{\nu }}{\stackrel{~}{ϵ}}.$$
(35)
Here the prime denotes the derivative, $`\mathrm{ln}C=\mathrm{ln}4+2C_E+1`$, and $`\stackrel{~}{\tau }=\stackrel{~}{\tau }(ϵ)`$ is determined by the equation
$$\mathrm{ln}\left(\frac{\stackrel{~}{ϵ}Ce}{192\pi ^2\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}\right)+4\mathrm{\Phi }(\stackrel{~}{m\tau })=0.$$
(36)
Note that the saddle point solution satisfies the relation
$$\frac{2\stackrel{~}{\gamma }}{m^3(\stackrel{~}{\chi }\stackrel{~}{\tau })^3}+\mathrm{\Phi }^{}(m\stackrel{~}{\tau })=0,$$
(37)
which will be used later.
Substituting the saddle point solution into Eq. (34) we obtain the leading order contribution $`F^{(1)}`$:
$$F^{(1)}(ϵ,\nu )=\kappa ϵ\left[m\stackrel{~}{\tau }_0(ϵ)+\frac{1}{4\mathrm{\Phi }^{}(m\stackrel{~}{\tau }_0(ϵ))}\right]+𝒪(\nu ).$$
(38)
In the limit $`ϵ0`$ Eq. (36) can be solved iteratively. One gets
$$m\stackrel{~}{\tau }(ϵ)=\frac{2}{\sqrt{\mathrm{ln}\frac{1}{ϵ}}}+\frac{\mathrm{ln}\mathrm{ln}\frac{1}{ϵ}}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^{3/2}}+\mathrm{}$$
(39)
Then in the leading order in energy solutions (35) become
$$(m\stackrel{~}{\rho })^2=\frac{1}{48\pi ^2}\frac{\stackrel{~}{ϵ}}{\left(\mathrm{ln}\frac{1}{\stackrel{~}{ϵ}}\right)^{3/2}};\stackrel{~}{\gamma }=\left(\frac{\stackrel{~}{\nu }}{\stackrel{~}{ϵ}}\right)^3\left(\mathrm{ln}\frac{1}{\stackrel{~}{ϵ}}\right)^{3/2};m(\stackrel{~}{\chi }\stackrel{~}{\tau })=2\frac{\stackrel{~}{\nu }}{\stackrel{~}{ϵ}}.$$
(40)
In this regime the function $`F^{(1)}(ϵ,\nu )`$ is equal to
$$F^{(1)}(ϵ,\nu )=2\frac{\kappa ϵ}{\sqrt{\mathrm{ln}\frac{1}{ϵ}}}\left[1+𝒪\left(\frac{\mathrm{ln}\mathrm{ln}\frac{1}{ϵ}}{\mathrm{ln}\frac{1}{ϵ}}\right)\right]+𝒪(\nu ).$$
(41)
### 3.2 Propagator correction
The next-to-leading order function $`W^{(2)}`$ can be written as the sum of contributions involving the propagator between final states, between initial and final states and between initial states, respectively:
$$W^{(2)}=W_{(ff)}^{(2)}+W_{(if)}^{(2)}+W_{(ii)}^{(2)}.$$
(42)
As we have already mentioned the expressions for these terms are given in Ref. . The complete propagator correction was calculated numerically, the results are discussed in Sect. 3.3. Here we study the propagator correction analytically in the limit of small $`\nu `$. Keeping only relevant contributions we obtain that
$`{\displaystyle \frac{1}{\lambda }}W_{(ff)}^{(2)}`$ $`=`$ $`I_{bb}(\tau ,\tau )+\mathrm{},`$ (43)
$`{\displaystyle \frac{1}{\lambda }}W_{(if)}^{(2)}`$ $`=`$ $`2\gamma I_{ab}(\tau ,\chi \tau )+\mathrm{},`$ (44)
$`{\displaystyle \frac{1}{\lambda }}W_{(ii)}^{(2)}`$ $`=`$ $`\gamma ^2I_{aa}(\chi \tau ,\chi \tau )+\mathrm{},`$ (45)
where
$`I_\mathrm{\#}(\tau _1,\tau _2)`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^6m^2}}{\displaystyle \frac{d𝐤}{2\omega _𝐤}\frac{d𝐪}{2\omega _𝐪}R_{inst}(𝐤)e^{\omega _𝐤\tau _1}R_\mathrm{\#}(𝐤,𝐪)R_{inst}(𝐪)e^{\omega _𝐪\tau _2}}`$
$`=`$ $`48{\displaystyle \frac{\rho ^4}{\lambda }}{\displaystyle \frac{d𝐤}{2\omega _𝐤}\frac{d𝐪}{2\omega _𝐪}e^{\omega _𝐤\tau _1}\frac{R\left(\rho ^2s_\mathrm{\#}(𝐤,𝐪)\right)}{16\pi ^2}e^{\omega _𝐪\tau _2}}.`$
The functions $`R_\mathrm{\#}(𝐤,𝐪)`$ and $`R(\rho ^2s)`$ are given by Eqs. (28) - (31), all necessary notations were introduced in Sect. 2.
In the limit of small $`\nu `$ the expression for the propagator correction in terms of simple integrals can be obtained. However, it is quite cumbersome and we do not present this result here. Instead we calculate and analyze groups of terms which are singular in $`\nu `$. From Eqs. (35) it follows that for the saddle point solution in the limit $`\nu 0`$ we have
$$m(\stackrel{~}{\chi }\stackrel{~}{\tau })\nu 0,\stackrel{~}{\gamma }\nu ^3.$$
(47)
Using these properties it is easy to select and calculate the terms in Eqs. (44) and (45) which are singular in $`\nu `$. Evaluating these terms at the saddle point solution (35), we obtain that
$`F_{(ii)}^{(2)}`$ $`=`$ $`32(192\pi ^2)\alpha _1\stackrel{~}{\rho }^6{\displaystyle \frac{\stackrel{~}{\gamma }^2}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^6}}\left[\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\rho }^2}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^2}}+\mathrm{}\right]`$ (48)
$`=`$ $`{\displaystyle \frac{8\alpha _1}{(192\pi )^2}}{\displaystyle \frac{\stackrel{~}{ϵ}^3}{\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\left[2\mathrm{ln}{\displaystyle \frac{1}{\stackrel{~}{\nu }}}+\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{ϵ}^3}{768\pi ^2\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\right)+\mathrm{}\right],`$
$`F_{(if)}^{(2)}`$ $`=`$ $`16(192\pi ^2)\alpha _1\stackrel{~}{\rho }^6{\displaystyle \frac{2\stackrel{~}{\gamma }}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^3}}\left[\mathrm{\Phi }^{}(m\stackrel{~}{\tau }_0)\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\rho }^2}{\stackrel{~}{\tau }(\stackrel{~}{\chi }\stackrel{~}{\tau })}}+\mathrm{}\right]`$ (49)
$`=`$ $`{\displaystyle \frac{16\alpha _1}{(192\pi )^2}}{\displaystyle \frac{\stackrel{~}{ϵ}^3}{\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\left[\mathrm{ln}{\displaystyle \frac{1}{\stackrel{~}{\nu }}}+\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{ϵ}^2}{384\pi ^2m\stackrel{~}{\tau }\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\right)+\mathrm{}\right],`$
where the dots stand for non-singular terms. Summing contributions (48) and (49) one gets
$`F_{(ii)}^{(2)}`$ $`+`$ $`F_{(if)}^{(2)}`$ (50)
$`=`$ $`16(192\pi ^2)\alpha _1\stackrel{~}{\rho }^6[{\displaystyle \frac{\stackrel{~}{\gamma }}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^3}}({\displaystyle \frac{2\stackrel{~}{\gamma }}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^3}}+\mathrm{\Phi }^{}(m\stackrel{~}{\tau }))\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\rho }^2}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^2}}`$
$`+`$ $`{\displaystyle \frac{\stackrel{~}{\gamma }}{(\stackrel{~}{\chi }\stackrel{~}{\tau })^3}}\mathrm{\Phi }^{}(m\stackrel{~}{\tau })\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\rho }^2}{\stackrel{~}{\tau }^2}}+\mathrm{}]`$
$`=`$ $`{\displaystyle \frac{8\alpha _1}{(192\pi )^2}}{\displaystyle \frac{\stackrel{~}{ϵ}^3}{\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\left[\mathrm{ln}\left({\displaystyle \frac{\stackrel{~}{ϵ}}{192\pi ^2m\stackrel{~}{\tau }\mathrm{\Phi }^{}(m\stackrel{~}{\tau })}}\right)+\mathrm{}\right].`$ (51)
As we see from the last line, Eq. (51), the singular terms $`\mathrm{ln}(1/\nu )`$ cancel each other. Eq. (50), reveals the reason of this cancellation: due to relation (37) the coefficient of the term $`\mathrm{ln}(\stackrel{~}{\rho }^2/(\stackrel{~}{\chi }\stackrel{~}{\tau })^2)`$, which gives rise to the singularity $`\mathrm{ln}(1/\nu )`$, is equal to zero exactly. This result is general and does not depend on any approximation.
We would like to remark that the terms singular in $`\nu `$ are proportional to $`\alpha _1`$. From Eq. (31) it follows that they originate from the terms proportional to $`s\mathrm{ln}s`$ and $`s`$ in the residue of the instanton propagator. If one uses the instanton propagator $`G_f(p,q)`$ satisfying constraint (11) with the functions $`f_A`$ such that two leading terms in the asymptotics (15) vanish, then the leading asymptotics $`s\mathrm{ln}s`$ and $`s`$ of the propagator for large $`s`$ are absent. As a consequence, the singular terms $`\mathrm{ln}(1/\nu )`$ do not appear.
For energies small enough, such that $`m\stackrel{~}{\tau }1`$, the expressions simplify further and the result for the next-to-leading correction can be written in a simple form. We obtain that
$$F^{(2)}(ϵ,\nu )=\frac{\alpha _2}{192\pi ^2}(\stackrel{~}{ϵ}m\stackrel{~}{\tau })^2\left[\mathrm{ln}\frac{\stackrel{~}{ϵ}m\stackrel{~}{\tau }}{384\pi ^2}+\frac{58}{15}+𝒪(m^2\stackrel{~}{\tau }^2)\right].$$
(52)
In the limit $`ϵ0`$ we use solution (39) and obtain
$$F^{(2)}(ϵ,\nu )=\frac{4\alpha _2\kappa ^2ϵ^2}{192\pi ^2}\left(1+\frac{1}{2}\frac{\mathrm{ln}\mathrm{ln}\frac{1}{ϵ}}{\mathrm{ln}\frac{1}{ϵ}}+\mathrm{}\right)=\frac{\kappa ^2ϵ^2}{32\pi ^2}\left(1+\frac{1}{2}\frac{\mathrm{ln}\mathrm{ln}\frac{1}{ϵ}}{\mathrm{ln}\frac{1}{ϵ}}+\mathrm{}\right)$$
(53)
We see that at low energies the main contribution is proportional to $`\alpha _2`$, i.e. comes from the $`\mathrm{ln}s`$ and constant terms in the residue (31) of the instanton propagator. In fact it is easy to check that it is precisely the term $`\mathrm{ln}s`$ in Eq. (31) which gives the contribution (53).
### 3.3 Numerical results
For arbitrary $`ϵ`$ and $`\nu `$ the functions $`F^{(1)}(ϵ,\nu )`$ and $`F^{(2)}(ϵ,\nu )`$ were studied numerically. It turned out that the saddle point solution exists only for a certain region in the $`(ϵ,\nu )`$-plane. It lies inside the rectangle $`0<ϵ<ϵ_{max}=0.55`$ and $`0<\nu <\nu _{max}=0.25`$. We performed the numerical analysis for this whole region.
To present the results it is convenient to introduce the following functions:
$$_1(ϵ,\nu )=1\frac{F^{(1)}(ϵ,\nu )}{32\pi ^2},_2(ϵ,\nu )=1\frac{F^{(1)}(ϵ,\nu )+F^{(2)}(ϵ,\nu )}{32\pi ^2}.$$
They are normalized by the conditions $`_1(0,\nu )=_2(0,\nu )=1`$.
Lines of constant $`_2(ϵ,\nu )`$ are plotted in Fig. 1. We want to study the cross section for shadow processes with a few initial particles. Then according to conjecture (3) points where the lines cross the $`\nu =0`$ axis are of particular interest. For example, $`_2(ϵ,0)=0.95`$ at $`ϵ=0.180`$, $`_2(ϵ,0)=0.85`$ at $`ϵ=0.492`$. We would like to mention that, in fact, in the studied region of $`(ϵ,\nu )`$ the propagator correction is quite small comparing to the leading order. Thus, the difference between $`_1(ϵ,\nu )`$ and $`_2(ϵ,\nu )`$ does not exceed $`10^2`$.
The curves in Fig. 1 end at the line formed by saddle points corresponding to the periodic instanton solutions. For them $`\stackrel{~}{\tau }(ϵ,\nu )=\stackrel{~}{\chi }(ϵ,\nu )/2`$. This line is directed from the zero energy instanton ($`ϵ=\nu =0`$) to the sphaleron ($`ϵ=\nu =1`$).
As it has been already mentioned, the complete function $`F(ϵ,\nu )`$ was calculated numerically in the range $`0.4<ϵ<3.5`$ and $`0.25<\nu <1`$ in Ref. . The computation was performed by solving a certain classical boundary value problem on the lattice. With the size of the lattice used in the numerical calculation in Ref. , the authors did not obtain data for smaller $`ϵ`$ and $`\nu `$ except for the line of the periodic instantons. The comparison shows that our perturbative results do not differ significantly from the exact ones of Ref. for $`ϵ<0.25`$ and $`\nu <0.2`$. These values can be regarded as a rough estimate of the range of validity of the leading and next-to-leading approximations.
## 4 Discussion and conclusions
In the present paper we have analyzed the multiparticle cross section of the shadow processes induced by instanton transitions in the simple scalar model (4). Using the exact analytical expression for the on-shell residue of the propagator of quantum fluctuations in the instanton background we calculated the suppression factor in the next-to-leading order.
The calculation of the leading and next-to-leading orders of $`F(ϵ,\nu )`$ was performed assuming that the size of the instanton solution is small enough, namely $`(\stackrel{~}{\rho }m)1`$. Neglecting $`𝒪(\rho ^4m^4)`$ terms in the action (7) and using the instanton and the residue of the instanton propagator of the massless theory in Eqs. (22) and (28) amount to omission of corrections of the type
$$\frac{\rho ^2}{\tau ^2}(\rho ^2m^2)$$
(54)
in $`F^{(1)}`$ and
$$(\rho m)^4,\frac{\rho ^4}{\tau ^4}(\rho ^2m^2),\frac{\rho ^4}{\tau ^4}(\rho ^2m^2)\mathrm{ln}\frac{\rho ^2}{\tau ^2}\text{and}\frac{\rho ^4}{\tau ^4}(m^2\tau ^2)^k$$
(55)
in $`F^{(2)}`$. We checked numerically that in the region of $`ϵ`$ and $`\nu `$, where the saddle point solution exists, the terms in Eqs.(54), (55) are really small. As an illustration let us consider the case of very small $`ϵ`$ and use the saddle point solution (39), (40). We obtain that
$`{\displaystyle \frac{\stackrel{~}{\rho }^2}{\stackrel{~}{\tau }^2}}(\stackrel{~}{\rho }^2m^2)`$ $``$ $`{\displaystyle \frac{ϵ^2}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^{5/2}}},`$
$`(\stackrel{~}{\rho }m)^4`$ $``$ $`{\displaystyle \frac{ϵ^2}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^3}},{\displaystyle \frac{\stackrel{~}{\rho }^4}{\stackrel{~}{\tau }^4}}(\stackrel{~}{\rho }^2m^2){\displaystyle \frac{ϵ^3}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^{5/2}}},`$
$`{\displaystyle \frac{\stackrel{~}{\rho }^4}{\stackrel{~}{\tau }^4}}(\stackrel{~}{\rho }^2m^2)\mathrm{ln}{\displaystyle \frac{\stackrel{~}{\rho }^2}{\stackrel{~}{\tau }^2}}`$ $``$ $`{\displaystyle \frac{ϵ^3}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^{3/2}}},{\displaystyle \frac{\stackrel{~}{\rho }^4}{\stackrel{~}{\tau }^4}}(m^2\stackrel{~}{\tau }^2)^k{\displaystyle \frac{ϵ^2}{\left(\mathrm{ln}\frac{1}{ϵ}\right)^{1+k/2}}}.`$
All these corrections are subleading compared to the terms retained in the function $`F^{(2)}(ϵ,\nu )`$, Eq. (53). Contributions due to non-zero mass amount to corrections in powers of $`(m\stackrel{~}{\tau })`$, where $`\stackrel{~}{\tau }(ϵ,\nu )`$ is the saddle point solution for $`\tau `$. In general, these corrections are not small, and all of them were taken into account by using the $`s`$-variable and the energy $`\omega _𝐤`$ of massive particles in Eqs. (28) - (30), (34), (3.2), etc. Our numerical analysis shows that the inequality $`m\stackrel{~}{\tau }<1`$ is satisfied, for example, for $`ϵ<0.4`$ if $`\nu `$ is close to $`\nu =0`$ and for $`ϵ<0.02`$ for the periodic instanton solutions. Comparing this to the region in the $`(ϵ,\nu )`$-plane in Fig. 1, for which we carried out the calculation in this article, one can see that our formalism, accounting for arbitrary $`m\stackrel{~}{\tau }`$, allows to enlarge considerably the range of validity of the next-to-leading approximation.
The range of validity of the next-to-leading order approximation of the function $`E(ϵ,\nu )`$ was estimated by comparing our results with numerical computations in Ref. for the values of $`ϵ`$ and $`\nu `$ for which the latter can be translated to the case of shadow processes, i.e., for periodic instantons. The comparison shows that the perturbative results do not differ significantly from the exact ones for $`ϵ0.25`$ and $`\nu 0.2`$.
For this range of values of $`ϵ`$ and $`\nu `$ and away from the line of periodic instantons, methods of Ref. do not allow to obtain exact results. Therefore, at the moment our perturbative calculations are the only ones which give quantitative behaviour of the supression factor in this range.
From Eqs. (39) we see that approximate formulas (41) and (53) are valid as long as
$$\frac{\mathrm{ln}\mathrm{ln}\frac{1}{ϵ}}{\mathrm{ln}\frac{1}{ϵ}}1.$$
For this range of energies we obtained the analytical expressions for the suppression factor and values of the saddle point parameters $`\stackrel{~}{\rho }`$, $`\stackrel{~}{\chi }`$, $`\stackrel{~}{\tau }`$ and $`\stackrel{~}{\gamma }`$. Formula (52) for the propagator correction for small $`\nu `$ is valid when $`m\stackrel{~}{\tau }1`$. According to the estimate, mentioned above, this condition is satisfied if $`ϵ0.4`$. This can be also verified by analyzing Eq. (36).
We also checked the cancellation of terms singular in the limit $`\nu 0`$ in the propagator correction $`F^{(2)}`$. As we have explained, this is closely related to the problem of quasiclassical evaluation of contributions of initial states and initial-final states. In the article we also discussed this problem within the approach proposed by Mueller. Namely, we calculated the leading asymptotics of the instanton propagator at large $`s`$ and showed that it can be cancelled by an appropriate choice of the propagator constraint. According to Ref. , with such propagator the problem of semiclassical calculation of contributions due to initial states and initial-final states can be tackled properly.
## Acknowledgments
We would like to thank A. Ringwald and V. Rubakov for discussions and valuable comments. Y.K. acknowledges financial support from the Russian Foundation for Basic Research (grant 98-02-16769-a) and grant CERN/P/FIS/1203/98. |
warning/0001/cond-mat0001350.html | ar5iv | text | # Heisenberg frustrated magnets: a nonperturbative approach
Understanding the effect of competing interactions in three dimensional classical spin systems is one of the great challenges of condensed matter physics. However, after twenty five years of investigations, the nature of the universality class for the phase transition of the simplest frustrated model, the antiferromagnetic Heisenberg model on a triangular lattice (AFHT model), is still a strongly debated question$`^{\text{[1]}}`$. Due to frustration, the ground state of the AFHT model is given by a canted configuration – the famous 120 structure – that implies a matrix-like order parameter$`^{\text{[2]}}`$ and thus, the possibility of a new universality class. Experiments performed on materials supposed to belong to the AFHT universality class display indeed exponents different from those of the standard $`O(N)`$ universality class: for VCl<sub>2</sub>$`^{\text{[3]}}`$: $`\beta =0.20(2),\gamma =1.05(3),\nu =0.62(5)`$, for VBr<sub>2</sub>$`^{\text{[4]}}`$: $`\alpha =0.30(5)`$, for CuFUD$`^{\text{[5]}}`$: $`\beta =0.22(2)`$ and for Fe\[S<sub>2</sub>CN(C<sub>2</sub>H<sub>5</sub>)<sub>2</sub>\]<sub>2</sub>Cl$`^{\text{[6, 7, 8]}}`$: $`\beta =0.24(1),\gamma =1.16(3)`$. These results however call for several comments. First, the exponents violate the scaling relations, at least by two standard deviations. Second, they differ significantly from those obtained by Monte Carlo (MC) simulations performed either directly on the AFHT model ($`\nu 0.59(1),\gamma 1.17(2),\beta 0.29(1),\alpha 0.24(2)`$), and on models supposed to belong to the same universality class: AFHT with rigid constraints ($`\nu =0.504(10),\gamma =1.074(29),\beta =0.221(9),\alpha =0.488(30)`$), dihedral (i.e. $`V_{3,2}`$ Stiefel) models ($`\nu 0.51(1),\gamma 1.13(2),\beta 0.193(4),\alpha 0.47(3)`$). See Ref. for a review, and references therein. Finally, the anomalous dimensions $`\eta `$ obtained by means of scaling relations is found to be negative in experiments as well as in MC simulations, a result forbidden by first principles for second order phase transitions$`^{\text{[10]}}`$. All these results are hardly compatible with the assumption of universality. It has been proposed that the exponents are, in fact, effective exponents characterizing a very weakly first order transition, the so-called “almost second order phase transition$`^{\text{[11, 12, 13]}}`$”.
From the theoretical point of view the situation is also very unsatisfactory since one does not have a coherent picture of the expected critical behaviour of the AFHT model between two and four dimensions. On the one hand, the weak coupling expansion performed on the suitable Landau-Ginzburg-Wilson (LGW) model in the vicinity of $`d=4`$ leads to a first order phase transition due to the lack of a stable fixed point$`^{\text{[14, 15, 16]}}`$. On the other hand, the low temperature expansion performed around two dimensions on the Non-Linear Sigma (NL$`\sigma `$) model predicts a second order phase transition of the standard $`O(4)/O(3)`$ universality class$`^{\text{[17]}}`$. Since there is no indication that these perturbative results should fail in their respective domain of validity – i.e. for small $`ϵ=4d`$ and small $`ϵ=d2`$ – this situation raises two problems. First, and contrary to what happens in the non-frustrated case, one cannot safely predict the three dimensional behaviour from naïve extrapolations of the perturbative results. Although a direct computation in three dimensions, possible on the LGW model$`^{\text{[18, 19]}}`$, can circumvent this difficulty, such an approach misses a second fondamental problem: the incompatibility between the symmetries of the NL$`\sigma `$ and LGW models. Indeed, the renormalization group flow drives the NL$`\sigma `$ model action towards an $`O(4)`$ symmetric regime, more symmetric than the microscopical system, a phenomenon that cannot occur within all previous treatments of the LGW model (see ref. and below). The LGW model is therefore unable to find the $`O(4)`$ behaviour which has been nevertheless observed numerically in $`d=2`$$`^{\text{[20]}}`$. This raises serious doubts on the perturbative analysis of the LGW model away from $`d=4`$. Reciprocally, the perturbative analysis of the NL$`\sigma `$ model, based on a Goldstone mode expansion, predicts an $`O(4)/O(3)`$ fixed point everywhere between $`d=2`$ and $`d=4`$, as for the $`N=4`$ ferromagnetic model, in contradiction with the perturbative LGW results and the experimental and numerical situation in $`d=3`$. All this suggests that non perturbative features could play a major role and thus imposes to go beyond the standard perturbative approaches.
In this letter we realize this program by using the Wilson renormalization group framework$`^{\text{[21]}}`$. We obtain a coherent picture of the physics of the AFHT model everywhere between $`d=2`$ and $`d=4`$. We find that the fixed point expected from the NL$`\sigma `$ model approach exists indeed in the vicinity of $`d=2`$ but disappears below – and close to – three dimensions. The transition for AFHT in $`d=3`$ is thus weakly first order contrary to the different predictions of both a new universality class$`^{\text{[2]}}`$ and an $`O(4)/O(3)`$ second order behaviour$`^{\text{[17]}}`$. We get effective exponents compatible with the numerical and experimental data quoted above. For generalization to $`N>4`$-component spins, we find the transition in $`d=3`$ to be second order with exponents in good agreement with recent extensive MC simulations – contrary to those found from three loop Padé-Borel resummed series$`^{\text{[19]}}`$.
Our approach relies on the concept of effective average action$`^{\text{[22, 23]}}`$, $`\mathrm{\Gamma }_k[\varphi ]`$, which is a coarse grained free energy where only fluctuations with momenta $`qk`$ have been integrated out. The field $`\varphi `$ corresponds to an average order parameter at scale $`k`$, the analog of a magnetization at this scale. At the scale of the inverse lattice spacing $`\mathrm{\Lambda }`$, $`\mathrm{\Gamma }_{k=\mathrm{\Lambda }}`$ is the continuum limit of the lattice hamiltonian obtained, for example, by means of an Hubbard-Stratonovich transformation. On the other hand, the usual free energy $`\mathrm{\Gamma }`$, generating one particle-irreducible correlation functions, is recovered in the limit $`k0`$. The $`k`$-dependence of $`\mathrm{\Gamma }_k`$ is controlled by an exact evolution equation$`^{\text{[24, 25]}}`$:
$$\frac{\mathrm{\Gamma }_k}{t}=\frac{1}{2}\text{Tr}\left\{(\mathrm{\Gamma }_k^{(2)}+R_k)^1\frac{R_k}{t}\right\}$$
(1)
where $`t=\mathrm{ln}k/\mathrm{\Lambda }`$. The trace has to be understood as a momenta integral as well as a summation over internal indices. In Eq.(1), $`R_k`$ is the effective infrared cut-off which suppresses the propagation of modes with momenta $`q<k`$. A convenient cut-off is provided by$`^{\text{[24, 26]}}`$: $`R_k(q)=Zq^2/(\mathrm{exp}(q^2/k^2)1)`$, where $`Z`$ is the field renormalization. In Eq.(1), $`\mathrm{\Gamma }_k^{(2)}`$ is the exact field-dependent inverse propagator – i.e. the second derivative of $`\mathrm{\Gamma }_k`$.
The effective average action $`\mathrm{\Gamma }_k`$ is a functional invariant under the symmetry group of the system and thus depends on all the invariants built from the average order parameter. In our case, it is well known that the order parameter is a set of two vectors $`\stackrel{}{\varphi _1}`$ and $`\stackrel{}{\varphi _2}`$ that can be gathered in a real $`N\times 2`$ matrix $`\varphi _{ab}`$ for $`N`$-component spins$`^{\text{[2]}}`$. The symmetry of the system is the usual spatial rotation group $`O(N)`$ times a $`O(2)`$ corresponding to the symmetry of the underlying triangular lattice$`^{\text{[17]}}`$. This $`O(2)`$ is realized on $`\varphi _{ab}`$ as a right $`O(2)`$ “rotation” that turns the $`\stackrel{}{\varphi _i}`$ into each other. There are two independent $`O(N)O(2)`$ invariants built out of $`\varphi _{ab}`$: $`\rho =\text{Tr}^t\varphi \varphi `$ and $`\tau =\frac{1}{2}\text{Tr}(^t\varphi \varphi )^2\frac{1}{4}(\text{Tr}^t\varphi \varphi )^2`$.
The exact effective average action involves all the powers of $`\rho ,\tau `$ and of derivative terms, and so Eq.(1) is a nonlinear functional equation, too difficult to be solved exactly in general. We therefore need to truncate it. One possibility is to keep in $`\mathrm{\Gamma }_k`$ only the momentum (i.e. derivative)-independent part, an approximation called the Local Potential Approximation (LPA). In the case of frustrated magnets, this has been considered by Zumbach$`^{\text{[11, 12, 13]}}`$. This approximation however misses the field-renormalization and worse, as described below, the phenomenon of enlarged symmetry around $`d=2`$ found perturbatively in the NL$`\sigma `$ model$`^{\text{[17]}}`$. This does not mean that this approximation is not useful: it is simply, in essence, unable to answer the question of the matching of the different perturbative approaches. Another truncation is however possible which preserves this possibility: it consists in an expansion of $`\mathrm{\Gamma }_k`$ around its minimum in order to keep a finite number of monomials in the invariants $`\rho `$ and $`\tau `$ while including the derivative terms allowing to recover the different perturbative results. We choose the simplest such truncation:
$$\begin{array}{c}\mathrm{\Gamma }_k=d^dx\{\frac{Z}{2}\varphi _{ab}\varphi _{ab}+\frac{\omega }{4}(ϵ_{ab}\varphi _{ca}\varphi _{cb})^2\hfill \\ \\ +\frac{\lambda }{4}(\frac{\rho }{2}\kappa )^2+\frac{\mu }{4}\tau \}\hfill \end{array}$$
(2)
where $`\{\omega ,\lambda ,\kappa ,\mu ,Z\}`$ are the coupling constants which parametrize the model. All terms but one - the “current term” $`(ϵ_{ab}\varphi _{ca}\varphi _{cb})^2`$ \- are very natural and correspond to those appearing in the usual LGW action that realizes the symmetry breaking scheme of frustrated magnets. Indeed for $`\lambda `$ and $`\mu 0`$, the minimum of the action is realized by a configuration of the form $`\varphi _{ab}^{min}=\sqrt{\kappa }R_{ab}`$, where $`R_{ab}`$ is a matrix built with two orthonormal $`N`$-component vectors. The symmetry of this minimum is a product of a diagonal $`O(2)`$ group and a residual $`O(N2)`$ group. The symmetry breaking scheme is thus $`O(N)O(2)O(N2)O(2)_{diag}`$$`^{\text{[17]}}`$. Note that for $`\varphi _{ab}=\varphi _{ab}^{min}`$ one has: $`\rho =2\kappa `$ and $`\tau =0`$ so that Eq.(2) corresponds indeed to a quartic expansion around the minimum. The spectrum in the low temperature phase consists in $`2N3`$ Goldstone modes and three massive modes: one singlet of mass $`m_1=\kappa \lambda `$ and one doublet of mass $`m_2=\kappa \mu `$ which correspond to fluctuations of the relative angle and of the norms of the two vectors $`\stackrel{}{\varphi _1}`$ and $`\stackrel{}{\varphi _2}`$.
Without the current term, the truncation Eq.(2) is however not sufficient in our case. This term plays a crucial role since, for $`N=3`$, it allows the model to enlarge its symmetry from $`O(3)O(2)`$ to $`O(3)O(3)O(4)`$ at the fixed point around $`d=2`$, leading to the well known $`O(4)/O(3)`$ behaviour$`^{\text{[17]}}`$. The current term is systematically discarded in the perturbative treatment of the LGW model around four dimensions, for the - correct - reason that it is power-counting irrelevant. Here we can include it in our ansatz since it is anyway present in the full effective action $`\mathrm{\Gamma }_k`$ and, in fact, must include it since it becomes relevant somewhere between two and four dimensions. The formalism we use is in charge to decide where it is important.
Let us emphasize that the effective average action method leads to non trivial and/or new results even within a quartic truncation of $`\mathrm{\Gamma }_k`$. One can mention the Kosterlitz-Thouless phase transition$`^{\text{[27]}}`$, low energy QCD$`^{\text{[28]}}`$, the abelian Higgs model and superconductivity$`^{\text{[29, 30]}}`$, etc. The accuracy of the results thus obtained depends on two main features: i) the smallness of the anomalous dimension $`\eta `$ and ii) the fact that the thermodynamics of the system is controlled by a unique minimum of $`\mathrm{\Gamma }_k`$. Note finally that this technique has been successfully employed in the case of the principal chiral model to solve a conflict between perturbative approaches$`^{\text{[31]}}`$, similar to what is studied here. However we stress that in the principal chiral case, there was no conflict between the symmetries of the LGW and NL$`\sigma `$ models.
The flow equations for the different coupling constants $`\kappa `$, $`\lambda `$, $`\mu `$, $`\omega `$ and $`Z`$ are derived by using Eq.(1) and Eq.(2) along the same lines as in . The explicit recursion equations are too long to display and not particularly illuminating (see ). Moreover, they require a numerical analysis, apart in $`d=2+ϵ`$ and in $`d=4ϵ`$ where, as we now see, they get analytically tractable.
The physics around two dimensions. Around two dimensions, one expects that the perturbative “Goldstone mode” expansion of the NL$`\sigma `$ model works well. In the Goldstone regime, the fluctuations of the modulus of $`\stackrel{}{\varphi }_1`$ and $`\stackrel{}{\varphi }_2`$ and of their relative angle are frozen. This corresponds to the large mass limit $`m_{1r}`$, $`m_{2r}\mathrm{}`$. In this limit, our equations greatly simplify since the coupling constants divide in two sets $`\{\kappa ,\omega ,Z\}`$ and $`\{\lambda ,\mu \}`$ that do not mix. We only quote here the flow equations for the renormalized coupling constants of the first set:
$$\{\begin{array}{c}\frac{d\kappa _r}{dt}=(d2+\eta )\kappa _r+\frac{N2}{2\pi }+\frac{1}{4\pi (1+\kappa _r\omega _r)}\hfill \\ \frac{d\omega _r}{dt}=(2+d+2\eta )\omega _r+\hfill \\ \frac{1+\kappa _r\omega _r+(N1)\kappa _r^2\omega _r^2+(N2)\kappa _r^3\omega _r^3}{2\kappa _r^2\pi (1+\kappa _r\omega _r)}\hfill \\ \\ \eta =\frac{d\mathrm{ln}Z}{dt}=\frac{3+4\kappa _r\omega _r+2\kappa _r^2\omega _r^2}{4\kappa _r\pi (1+\kappa _r\omega _r)}\hfill \end{array}$$
(3)
These equations admit a fixed point for any $`N>2`$ of coordinates $`\kappa _r1/ϵ`$, $`\omega _rϵ`$, while $`\lambda _r,\mu _r`$ cst. The masses $`m_{1r}^{}`$, $`m_{2r}^{}`$ are thus very large, proving the consistency of the limit. In fact, modulo the change of variables: $`\eta _1=\kappa _r`$ and $`\eta _2=2\kappa _r(1+\kappa _r\omega _r)`$ the equations for $`\kappa _r`$ and $`\omega _r`$ are exactly those obtained at one-loop in the perturbative analysis of the NL$`\sigma `$ model$`^{\text{[17]}}`$. For $`N=3`$, they admit a fixed point for which the model is $`O(4)`$-symmetric.
Let us now recall how this phenomenon of enlarged symmetry for $`N=3`$ can be understood directly on the partition function. At the fixed point, the potential gets infinitely deep so that one recovers the hard constraints of the NL$`\sigma `$ model: $`\stackrel{}{\varphi }_1\stackrel{}{\varphi }_2`$, and $`\stackrel{}{\varphi }_{1}^{}{}_{}{}^{2}=\stackrel{}{\varphi }_{2}^{}{}_{}{}^{2}=\kappa _r^{}`$. For $`N=3`$, this allows us to rewrite the current term as the kinetic term of a third vector, the cross product of the two others: $`(ϵ_{ab}\varphi _{ca}\varphi _{cb})^2(\stackrel{}{\varphi }_3)^2`$ with $`\stackrel{}{\varphi }_3=\stackrel{}{\varphi }_1\stackrel{}{\varphi }_2`$. The order parameter of the system is then a trihedral of orthogonal vectors $`(\stackrel{}{\varphi }_1,\stackrel{}{\varphi }_2,\stackrel{}{\varphi }_3)`$. Thus contrary to what could be expected from a naïve expansion in powers of the fields, the current term plays a role as important as the usual kinetic terms. At the fixed point, $`\omega _r`$ takes a value such that the three vectors play a symmetric role and the symmetry breaking scheme is $`O(3)O(3)/O(3)O(4)/O(3)`$ instead of $`O(3)O(2)/O(2)`$. Such a result is of course missed within the LPA$`^{\text{[11, 12, 13]}}`$. Therefore, the presence of the current term does not only improve the accuracy of the calculation, it is necessary for its consistency.
The physics around four dimensions. Around four dimensions, we have expanded our equations at leading order in the coupling constants $`\lambda _r`$ and $`\mu _r`$. At this order the current term decouples and we are left with the following equations for the quartic coupling constants:
$$\{\begin{array}{c}\frac{d\lambda _r}{dt}=(4+d)\lambda _r+\frac{1}{16\pi ^2}(4\lambda _r\mu _r+4\mu _r^2+\lambda _r^2(4+N))\hfill \\ \\ \frac{d\mu _r}{dt}=(4+d)\mu _r+\frac{1}{16\pi ^2}(6\lambda _r\mu _r+N\mu _r^2).\hfill \end{array}$$
(4)
They are those obtained at one loop from the LGW approach$`^{\text{[14]}}`$. These flow equations admit a stable fixed point for $`N>N_c21.8`$, attesting that the phase transition is second order. For $`N<N_c`$ the transition is first order since no fixed point is found.
To higher orders, $`N_c`$ depends on the dimension. In $`d=3`$, three loop calculations resummed à la Padé-Borel predict $`N_c(d=3)=3.91`$$`^{\text{[18]}}`$. Note however that this calculation exhibits unusual behaviours compared to the $`O(N)`$ case: the coefficients of the series do not decrease monotonically and the series themselves are not alternate$`^{\text{[19]}}`$. These features reveal the poor summability of the series. Finally, in the $`N=6`$ case, for which the transition is second order, the predictions based on a Padé-Borel resummation, which provides $`\nu =0.575`$ and $`\gamma =1.121`$$`^{\text{[19]}}`$, are in clear disagreement with recent numerical simulations, for which $`\nu =0.700(11)`$ and $`\gamma =1.383(36)`$$`^{\text{[19]}}`$.
From this point of view our approach has several avantages: first, since it matches with the one loop perturbative results in $`d=2`$ and $`d=4`$ it is likely that the error does not vary much with the dimension – a fact that has been confirmed in the $`O(N)`$ case for which the precision for a given truncation is almost uniform with $`d`$. Second, it does not rely on a Padé-Borel resummation and therefore is free of the above mentionned problems of convergence. Of course, our results will change while improving the ansatz Eq.(2) by incorporating terms of higher order in fields and derivatives. However, all cases already treated within the average action method suggest that the lowest order approximation gives fairly good results, even with this crude approximation. For example, in the ferromagnetic $`O(3)`$ model, one finds $`\nu =0.703`$$`^{\text{[33]}}`$ which has to be compared to the six loop resummed perturbation series in three dimensions which provide $`\nu =0.705`$$`^{\text{[10]}}`$.
The physics between two and four dimensions. Let us first study the fate of the fixed point found analytically in $`d=2+ϵ`$ for $`N=3`$. By numerically integrating the flow equations, we find that this stable $`O(4)/O(3)`$ fixed point describes a smooth trajectory in the coupling constant space while $`d`$ is increased. Our flow equations actually admit another – but unstable – fixed point, which moves toward the stable fixed point as the dimension is increased. At a critical dimension $`d_c2.87`$, the two fixed points collapse and disappear. Above $`d_c`$, no other stable fixed point is found and we conclude that the transition is first order in $`d=3`$. We thus show that the $`O(4)/O(3)`$ fixed point obtained from the NL$`\sigma `$ model plays no role in the three dimensional physics of frustrated magnets, as conjectured for example by Jolicœur and David$`^{\text{[34]}}`$ and Dobry and Diep$`^{\text{[35]}}`$. We also discard the possibility of a new universality class conjectured on the basis of a naïve extrapolation of the $`ϵ=4d`$ calculation$`^{\text{[1, 2]}}`$. The proximity of $`d_c`$ with $`d=3`$ however let open the possibility of a very weakly first order phase transition with effective critical exponents. This behaviour manifests itself in our equations by the existence of a minimum around which the RG flow slows down. This characterizes a very large, although finite correlation length $`\xi `$. A rough estimate of this correlation length – a few hundred lattice spacings – indicates that a pseudo-scaling behaviour can be observed although $`\xi `$ is not large enough to ensure a true universality. This could explain the broad spectrum of effective critical exponents found in experiments and numerical simulations. Although the flow equations do not have a fixed point, we are able to compute effective exponents by linearizing the flow equations around the minimum. We recover here the phenomenon of “almost second order phase transition” first introduced by Zumbach$`^{\text{[11, 12, 13]}}`$ within the LPA. To get accurate results we have to take into account the $`\varphi ^6`$-like terms in our ansatz. We find: $`\nu =0.53`$, $`\gamma =1.03`$ and $`\beta =0.28`$, which lie in between the various sets of exponents found experimentally and numerically (see above). For comparison Zumbach found $`\nu 0.63`$ in the LPA$`^{\text{[11, 12, 13]}}`$, the difference being mainly due to the anomalous dimension.
Finally, we find a true fixed point in $`d=3`$ for $`N`$ larger than a critical value $`N_c(d=3)4`$. For $`N=6`$, we get $`\nu =0.74`$ and $`\gamma =1.45`$ which compare well with the Monte Carlo data $`\nu =0.700(11)`$ and $`\gamma =1.383(36)`$$`^{\text{[19]}}`$. They are close to the LPA results, where $`\nu =0.76`$$`^{\text{[12]}}`$, and much better than those obtained by a three-loop calculation in $`d=3`$$`^{\text{[19]}}`$ (see above). We have checked that our exponents do not vary significantly when monomials of order six in the fields are included in the ansatz Eq.(2).
To conclude, using a non perturbative method, we have reached a global understanding of frustrated Heisenberg magnets including a matching between previous perturbative predictions and a good agreement with experimental and numerical data. It remains to understand the very origin of the disappearance of the NL$`\sigma `$ model fixed point. The role of non trivial topological configurations can be invoked. One can hope a complete understanding of this point through the average action method which successfully describes the Kosterlitz-Thouless phase transition$`^{\text{[27]}}`$.
We thank J. Vidal for a careful reading of the manuscript
LPTHE is a laboratoire associé au CNRS UMR 7589. e-mail: tissier,delamotte,mouhanna@lpthe.jussieu.fr |
warning/0001/astro-ph0001473.html | ar5iv | text | # Radio-optical identification of very-steep spectrum radio sources from the UTR-2 catalogue
## 1 Introduction
A radio survey obtained with the UTR telescope (Kharkov, Ukraine) at frequencies 10–25 MHz has resulted in a catalog of 1822 sources (Braude et al. 1978–1994; www.ira.kharkov.ua/UTR2). Covering about 30% of the sky north of $``$13 declination, this survey is presently the lowest-frequency source catalog of its size, and thus provides an ideal basis to study the little known optical identification content of sources selected at decametric frequencies. In the original version of the UTR-2 catalog (UTR in what follows) there is no radio identification at other frequencies for 7% of the sources, and for 81% there is no optical identification. Our goal is to identify all UTR sources with known radio sources and to search for optical counterparts on the Digitzed Sky Surveys.
Here we present our first results on a subsample of ultra-steep spectrum (USS) sources (spectral index $`\alpha `$1.2, S $`\nu ^\alpha `$). This class of sources is being actively studied by various groups (Parijskij et al. 1996; Röttgering et al. 1997; McCarthy et al. 1997), mainly because they are often identified with very distant radio galaxies, which are probes of the early Universe and thought to be indicators of proto-clusters (e.g. Djorgovski 1987).
## 2 Radio identification
The rather large uncertainties of UTR positions ($``$0.7) require an iterative process for finding radio counterparts at successively higher frequencies (and thus higher positional accuracy). In this we aided ourselves by selecting previously cataloged sources from the CATS database (Verkhodanov et al. 1997a) in a box of RA$`\times `$DEC = 40$`{}_{}{}^{}\times `$40 centred on the nominal UTR position. The “raw” spectra given by these fluxes were refined using computer charts of source locations around UTR positions. All counterparts from TXS, GB6 and PMN within circles of 1 radius were considered one source. Groups of sources lying further apart were assigned separate spectra, each with the UTR flux as upper limit.
We were able to fit spectra for all but 7 of the 2314 radio counterparts to UTR sources. Fits were either straight (S), convex (C<sup>-</sup>), or concave (C<sup>+</sup>) curves in the lg $`\nu `$–lg S plot (see also Verkhodanov et al. 1997b, 1998). The resulting catalog includes information from a large number of electronically available catalogs of radio, infrared, optical and X–ray sources. The distribution of radio source spectra among the various spectral types is given in Table 1, and Fig. 1 shows the spectral index distribution of sources at high and low Galactic latitudes.
## 3 Subsample of ultra-steep spectrum sources
In our catalog of 2314 radio counterparts (the full list will be published in a forthcoming paper) there are 422 S-type sources with “very-steep spectrum” (VSS), and for the present work we selected from these a subsample of 102 “ultra-steep spectrum” (USS) objects ($`\alpha `$1.2). To further increase the radio-positional accuracy, we searched for radio counterparts of USS sources in the February (1998) version of the FIRST catalog (White et al. 1997), resulting in 38 FIRST counterparts for 23 UTR sources (see Table 2). If a UTR source has more than one acceptable counterpart in FIRST, we label these components with letters a, b, c, etc. Only one of the FIRST components (labeled GR1527+51 b) is truly unresolved by the FIRST beam of $``$ 5<sup>′′</sup> (i.e. has a major and minor axis of $`<`$ 2<sup>′′</sup>), while all other objects have a multi-component or extended structure. We checked the sources also in the lower resolution NVSS at 1.4 GHz (Condon et al. 1998). Usually, the larger the source complex, the larger the NVSS/FIRST flux ratio. Radio spectra of some of the source complexes are shown in Fig. 2. Examples of two FIRST maps of multi-component objects GR0910+48 and GR0942+54 overlaid on DSS-2 images are shown in Fig. 3. According to the NASA Extragalactic Database (NED), the complex source GR0135$``$08 is identified with the $`z=0.041`$ galaxy MCG-02-05-020, GR1214$``$03 is an LCRS QSO at z=0.184, and GR1243+04 is a radio galaxy (4C+03.24) at z=3.57. Our investigation of DSS-2 images for the unidentified sources is in progress.
Acknowledgements We wish to thank A.P. Miroshnichenko and D. Krivitskij of Institute of Radio Astronomy (Kharkov) for providing data and for useful discussion. The co-creators of the CATS database Sergei Trushkin and Vladimir Chernenkov (SAO RAS) provided useful comments. |
warning/0001/astro-ph0001199.html | ar5iv | text | # Superluminal Caustics of Close, Rapidly-Rotating Binary Microlenses
## 1 Introduction
Microlensing is gravitational lensing where the separations between the images are too small to be resolved. The directly detectable consequence of microlensing is that the brightness of the source varies in a way determined by the lens properties and the projected lens-source trajectory. Paczyński (1986) pointed out that microlensing would be a useful tool to detect massive compact halo objects. Microlensing surveys have since been carried out towards the Galactic bulge, the Magellanic Clouds, and M31. About 500 microlensing events have been detected to date (see Mao 2000 for a review).
Gravitational lensing by two point masses was carefully studied by Schneider & Weiss (1986). The most striking feature of such binary lensing is its caustics - one or several closed curves in the source (objects to be imaged) plane where a point source is infinitely magnified by the lens. As a reflection of the caustic structure of the magnification, the light curves of such binary lensing events may have multiple peaks. Microlensing surveys have detected about 30 such events (e.g., Udalski et al. 1994; Alcock et al. 2000). Star-planetary systems are an extreme form of binary. Mao & Paczyński (1991) first suggested that extrasolar planetary systems could be discovered by microlensing surveys.
Dominik (2000) conducted the first systematic investigation of the effect of binary rotation on microlensing light curves. Although all physical binaries rotate, static models suffice to reproduce the light curves of the great majority of observed microlensing events, even those with superb data such as MACHO 97-BLG-28 (Albrow et al. 1999a) and MACHO 98-SMC-1 (Afonso et al. 2000). The only event observed to date for which a rotating model is required is MACHO 97-BLG-41 (Albrow et al. 2000), and static models are excluded for this event only because the source traverses two disjoint caustics, a rare (so far, unique) occurrence. Bennett et al. (1999) had earlier proposed that the odd light curve of MACHO 97-BLG-41 could be explained by a triple-lens system consisting of a binary plus a jovian-mass planet.
In the treatments of binary rotation given to date (Dominik 1998; Albrow et al. 2000), the light curve is actually calculated by considering a series of static binaries, each with the configuration of the binary being modeled at the instant when the light ray from the source passes the plane of the center of mass of the lens. That is, the deflection of light by the binary, $`\alpha \alpha `$, is taken to be the vector sum of the deflections produced by the two components of the binary, $`\alpha \alpha =\alpha \alpha _1+\alpha \alpha _2`$, according to the Einstein (1936) formula,
$$\alpha \alpha _i=\frac{4GM_i}{b_i^2c^2}𝐛_i.$$
$`(1)`$
Here $`M_i`$ is the mass, and $`𝐛_i`$ is the impact parameter of the $`i`$th component of the binary.
This approach is strictly valid only in the limit
$$ϵ1,ϵ\frac{\omega b}{c},$$
$`(2)`$
where $`2\pi /\omega `$ is the period of the binary. For MACHO 97-BLG-41, the only microlensing event for which rotation has been measured (Albrow et al. 2000), $`ϵ10^4`$, so this approach is certainly valid. However, in principle $`ϵ`$ can be close to unity or can even greatly exceed unity. In this case, it is necessary to take account of the binary motion during the time that the source light is passing close $`(b)`$ to the lens plane. The Einstein formula (1), which was calculated for a static lens, is then no longer valid.
Here we study the rotation effects on the caustic structure of close, rapidly-rotating binary lenses. We present our main idea and method in § 2 and discuss the rotation effects in § 3. In § 4, we summarize our results and discuss some possible applications.
## 2 Main Idea and Method
As mentioned in § 1, the binary phase varies during the time that the photons are traveling from the source to the observer. This modifies the calculation of the instantaneous magnification map, especially for $`ϵ1`$. The retarded gravitational potential then begins to differ significantly from the naive Newtonian potential, which would normally be adequate in the weak-field limit and which is used to derive equation (1).
### 2.1 Retarded Gravitational Potential
As an approximate result of Einstein’s field equations, the deflection of a light ray passing through a static gravitational field can be expressed as an integral of the gradient of the Newtonian gravitational potential performed along the trajectory of the light (Bourassa et al. 1973):
$$\alpha \alpha =\frac{2}{c}_{\mathrm{}}^+\mathrm{}\varphi \mathrm{d}t,$$
$`(3)`$
which yields equation(1). However, for the non-static case, the configuration of the gravitational field will propagate at light speed, and we must instead use the retarded potential.
In analogy to the results of classical electrodynamics (e.g. Jackson 1975), the gravitational potential at field point r and time $`t`$ is contributed by every mass point r at an earlier time $`t^{}=t|𝐫𝐫^{}|/c`$,
$$\varphi (𝐫,t)=\frac{G\rho (𝐫^{},t^{})}{|𝐫𝐫^{}|}\mathrm{d}^3\mathrm{r}^{},$$
$`(4)`$
where $`\rho `$ is the (time-dependent) mass distribution.
For a point mass $`M`$, the retarded potential (4) can be written, similarly to the Liénard-Wiechert potential,
$$\varphi (𝐫,t)=\frac{GM}{(1𝐧^{}\beta \beta ^{})|𝐫𝐫^{}|},$$
$`(5)`$
where $`𝐧^{}=(𝐫𝐫^{})/|𝐫𝐫^{}|`$, $`\beta \beta ^{}`$ is the velocity of the point mass divided by $`c`$, and the prime denotes the value at time $`t^{}`$,
$$t=t^{}+\frac{|𝐫𝐫^{}(t^{})|}{c}.$$
$`(6)`$
The Newtonian gravitational field, $`𝐠=\varphi `$, is then
$$𝐠(𝐫,t)=\left[\frac{GM}{(1𝐧\beta \beta )|𝐫𝐫^{}|}\right]_t^{}=\left[\frac{GM}{(1𝐧\beta \beta )^2|𝐫𝐫^{}|^2}(𝐧+\beta \beta )\right]_t^{}.$$
$`(7)`$
For those rays that pass across the lens at a distance much greater than the Schwarzschild radius $`2GM/c^2`$, the total deflection angle caused by several lens objects is a superposition of the individual deflections (Bourassa et al. 1973). From equations (3) and (7), we have,
$$\alpha \alpha =\underset{i}{}\frac{2GM_i}{c}_{\mathrm{}}^+\mathrm{}\left[\frac{𝐧_i+\beta \beta _i}{(1𝐧_i\beta \beta _i)^2|𝐫𝐫_i|^2}\right]_t^{}dt.$$
$`(8)`$
It is convenient to take the time when the photon crosses the lens plane as $`t=0`$. Then at a distance $`|ct|`$ to the lens plane, the photon will feel the potential caused by the point mass $`M_i`$ at time $`t^{}`$. We have
$$|𝐫𝐫_i(t^{})|=\sqrt{b_i^2(t^{})+(ct)^2},$$
$`(9)`$
where $`b_i(t^{})`$ is the distance from the point mass $`M_i`$ to the impact point at time $`t^{}`$. Substitution of equation (9) into equation (6) yields
$$t=\frac{1}{2}t^{}\frac{b_i^2(t^{})}{2c^2t^{}}.$$
$`(10)`$
This relation between $`t`$ and $`t^{}`$ reflects the retarded effect. For a finite $`b_i(t^{})`$, and under the condition $`t>t^{}`$, we find that when $`t\mathrm{}`$, $`t^{}\mathrm{}`$ and when $`t+\mathrm{}`$, $`t^{}0^{}`$. That is $`t^{}2t`$ when $`t0`$ and $`t^{}0`$ when $`t0`$. For a circular orbit, this means that as a photon moves towards the rotating system, it will “find” that the angular velocity of the system is nearly doubled, and as it moves away from the system, it will feel an almost static field. Equation (10) makes it possible to replace the integration variable in equation (8) with $`t^{}`$ (from $`\mathrm{}`$ to $`0`$).
### 2.2 Lens Equation
The lens equation, i.e. the light-ray deflection equation, tells one how the light-ray deflection maps points in the source plane into points in the image plane (lens plane). If a photon comes from point $`\eta \eta `$ in the source plane and hits point $`\zeta \zeta `$ in the lens plane, we have (Schneider & Weiss 1986),
$$\eta \eta =\frac{D_s}{D_l}\zeta \zeta D_{ls}\alpha \alpha (\zeta \zeta ),$$
$`(11)`$
where $`D_s`$ and $`D_l`$ are the distances from the observer to the source plane and to the lens plane, respectively, and $`D_{ls}=D_sD_l`$ is the distance between the source and the lens.
In this paper, we consider a simple case: the lens is composed of two stars with equal masses $`M_1=M_2`$, rotating about each other in a circular face-on orbit. We define the distance between two stars to be $`2a`$ and the angular velocity around the center of mass to be $`\omega `$ (see Fig. 1). We define the radius of the Einstein ring generated by the total mass $`M_1+M_2`$ to be (e.g. Schneider & Weiss 1986)
$$r_\mathrm{E}=\sqrt{\frac{4G(M_1+M_2)}{c^2}\frac{D_lD_{ls}}{D_s}},$$
$`(12)`$
and then normalize the coordinates of points at the lens plane and those at the source plane using this radius and the radius projected onto the source plane, respectively:
$$X=\frac{a}{r_\mathrm{E}},$$
$`(13)`$
$$𝐫=\frac{\zeta \zeta }{r_\mathrm{E}},$$
$`(14)`$
$$𝐱=\eta \eta \left[r_\mathrm{E}\frac{D_s}{D_l}\right]^1.$$
$`(15)`$
The consequent lens equation is given in Appendix A.
The focus of the present study is the case $`X1`$, for which there are two small triangular caustics lying at projected distances $`r_\mathrm{E}/(2X)`$ from the binary center of mass (Schneider & Weiss 1986). As the binary rotates at angular speed $`\omega `$, the caustics rotate with it. Hence the transverse speed of the caustics is $`v/(2X^2)`$, where $`v=\omega a`$ is the speed of the binary components. For very small $`X`$, the caustics can move much faster than the speed of light (“superluminal motion”) even when the binary itself is well within the non-relativistic regime.
Since we have fixed the mass ratio, $`M_2/M_1=1`$, and adopted a face-on, circular orbit, there are only three free parameters of the lens system, $`X`$, $`\omega `$, and $`a`$. However, from the standpoint of studying the caustic structure that appears in diagrams from which all physical dimensions have been scaled out, it is only necessary to consider dimensionless parameters. There are two independent such parameters. One is $`X`$. There are two obvious possible choices for the other,
$$\beta =\frac{\omega a}{c},ϵ=\frac{\omega r_\mathrm{E}}{2cX}=\frac{\beta }{2X^2}.$$
$`(16)`$
The first is the speed of the lenses as a fraction of the speed of light, and the second is approximately the speed of the caustics as a fraction of the speed of light. As we show in the next section, $`ϵ`$ is a more useful parameter than $`\beta `$ because the caustics are more directly affected by their own speed rather than that of the lenses.
To make a two-dimensional magnification map of the lens system (around the caustics), we use the inverse ray-shooting technique (e.g., Schneider & Weiss 1986; Wambsganss 1997). Uniformly distributed light rays in the lens plane are evolved back to the source plane according to the lens equation. The magnification of each point in the source plane is then proportional to the density of rays at this point. We study the cases of $`X=0.1`$, $`X=0.05`$ with various values of $`ϵ`$.
## 3 Effects of the Rotation
Figure 2 displays some examples of the outer caustics generated by adopting different parameters. Compared with the static case, these caustics show some new features.
### 3.1 “Orbit Position” of the caustic
In the static case, the outer caustics of an equal-mass binary are located along the perpendicular bisector of the binary. If the binary lens rotates around the center while the photon is traveling, it is not hard to imagine that the consequent caustics will drift with respect to the bisector of the binary lens at phase $`t=0`$. It turns out that the direction of this drift is opposite to the rotation. Physically, the reason for this “opposite” drifting is that at all times $`t`$, the phase of the binary corresponds to an earlier time $`t^{}<0`$. See equation (10) and the analysis following it. In Figure 3, we show the angular position of the caustic relative to the static case as a function of $`\beta `$. For $`\beta 0.1`$, we find a fitted formula for the “orbit” angle $`\theta _{orbit}`$ (in radians):
$$\theta _{orbit}=\frac{4}{7}\beta =\frac{8}{7}X^2ϵ.$$
$`(17)`$
We find that the fitted coefficient in equation (17) is very close to the ratio of two small integers (4/7), but we do not know whether this result is exact.
### 3.2 “Spin” - Pointing of the caustic
Apart from the outer caustics’ orbit motion as a whole, these “concave triangles” have their own rotation. We use the direction of a vertex (the one pointing to the center of mass in the static case) as a tracer for the caustic’s “spin”. Unlike the Moon which always shows the same hemisphere to the Earth, these triangles spin much faster than their orbital rotation. For the non-static case, they will no longer point to the center of mass. See Figure 4. The fitted formula for the spin angle $`\theta _{spin}`$ (in degrees) for $`ϵ<2`$ is
$$\mathrm{log}\theta _{spin}=31.32\mathrm{log}ϵ+0.02154.$$
$`(18)`$
### 3.3 “Expansion” - Enlargement of the caustic
For static case, it is known that the closer the binary is, the farther the triangular caustics are away from the center of mass and the smaller they become. The outer caustics shrink almost to a point in the case of very small $`X`$. However, after taking the rotation effect into account, we find that the tiny caustics are strikingly magnified. Meanwhile, unlike the static case, the shape of the caustics gradually loses its symmetry. We choose the area inside the caustic as a measure of the expansion effect. Since in the static case, the linear size of the outer caustic scales approximately as $`X^3`$, we normalize the area in our case by $`X^6`$. The expansion effect is illustrated in Figure 5.
### 3.4 Magnification Properties
The rapidly rotating binary lens makes the outer caustics move with a speed comparable to light speed which in turn brings some new phenomena to the magnification properties of the caustics.
First, the expansion of the outer caustic dilutes the magnification. In the static case, the magnification factor near a caustic curve can be described as $`A(u)=A_0+(\mathrm{\Delta }u_{}/u_r)^{1/2}`$, where $`\mathrm{\Delta }u_{}`$ is the perpendicular distance to the caustic and $`A_0`$ and $`u_r`$ are constants (Schneider & Weiss 1986; Schneider & Weiss 1987; Albrow et al. 1999b). Hence $`u_r`$ describes the strength of the caustic. We investigate the variation of $`u_r`$ as a function of $`ϵ`$ and find that rapid rotation weakens the strength of the caustics. The relative strength of the three caustic lines (of the triangular caustic) also changes (see Fig. 6 for the case $`X=0.1`$): one caustic line (left border, see Fig. 2) becomes the strongest one at large $`ϵ`$ (large area). With a dimension of linear size, $`u_r`$ is expected to scale as $`X^3`$. However, our calculation shows that there is a slight deviation from this scale law: the change of $`u_r`$ for $`X=0.05`$ is a little steeper than that for $`X=0.1`$. Taking account of our resolution, we are not sure whether this marginal effect is real or not.
Secondly, the velocity of the source with respect to the outer caustic is overwhelmingly determined by the caustic’s high speed. The relative source-caustic trajectory is then a small piece of arc centered at the binary center of mass. The rapid rotation implies that the timescale for crossing the caustic will be very short. We choose the square root of the area $`S`$ inside the caustic as the size of the caustic and the crossing time is then
$$t_c=\frac{\sqrt{S}}{ϵ}\frac{r_\mathrm{E}}{c},$$
$`(19)`$
where $`r_\mathrm{E}`$ is the Einstein ring defined in equation (12). Since the area scales very nearly as $`X^6`$ (see Fig 5), this time scale can be normalized by $`X^3`$. As shown in Figure 7, $`t_c`$ has a minimum near $`ϵ1`$. At smaller $`ϵ`$, $`t_c`$ becomes larger mainly due to the “low” speed of the caustic and at larger $`ϵ`$, mainly due to the expansion of the area.
## 4 Discussions
In this paper, we point out that a modification is necessary to get the instantaneous magnification map for close, rapidly-rotating binary lenses, in which case the outer caustics have a very high speed which can even be superluminal. Taking the retarded gravitational potential into consideration, we investigate the outer caustic behavior for such a lens with equal masses and face-on circular orbit. Compared with the static case, the caustic is displaced in orbit position and is rotated about its own axis. The most remarkable result is the enlargement of the caustic by the rapid motion of the lens. This increase in size induces a corresponding drop in the strength of the caustic.
Instead of starting from Einstein’s field equation, we use a retarded potential at the first step. Although strictly speaking this method has its limitations, it is a reasonable approach for our purpose since our analysis focuses on the high speed of the outer caustics while the binary itself is not in the extreme relativistic regime.
What would be required to observe superluminal caustics? Combining the definitions of $`ϵ`$ and $`r_\mathrm{E}`$ with Kepler’s Third Law yields,
$$ϵ=\frac{(2GM)^{3/2}D}{c^3a^{5/2}}=0.37\left(\frac{M}{M_{}}\right)^{3/2}\left(\frac{a}{0.1\mathrm{AU}}\right)^{5/2}\frac{D}{2\mathrm{kpc}},$$
$`(20)`$
where $`D=D_lD_{ls}/D_s`$, $`M`$ is the mass of one binary component, and where we have normalized to the case of a pair of solar-mass stars seen half way to the Galactic center. Equation (20) can be rewritten in terms of $`X`$,
$$X=2^{9/10}ϵ^{2/5}\left(\frac{GM}{Dc^2}\right)^{1/10}.$$
$`(21)`$
Hence, to obtain $`ϵ1`$ would require $`X0.012`$, a factor 4 smaller than even the lesser of the two values that we examined in this paper. Note that this result is extremely insensitive to either $`M`$ or $`D`$.
From Figure 5, the combined cross section (linear size) of the two caustics at $`ϵ1`$ is $`S^{1/2}8X^3r_\mathrm{E}`$, that is, a factor $`4X^3`$ smaller than for the lens itself. For $`X=0.012`$, this factor is $`10^5`$, so that at first sight it appears completely hopeless that superluminal caustics would ever be observed. However, the event rate is the product of the cross section with the transverse speed, and the caustic moves $`10^3`$ times faster than transverse speed of the binary center of mass. In fact, since the caustic is likely to be smaller than the source, the event rate is given by the source size times the transverse speed of the caustic. Thus, the ratio of superluminal-caustic events to normal events generated by the same binary is
$$\frac{\mathrm{\Gamma }_{\mathrm{super}}}{\mathrm{\Gamma }_{\mathrm{normal}}}\frac{2R_{}D_l/D_s}{r_\mathrm{E}}\frac{ϵc}{v_{}}1.2\frac{R_{}/R_{}}{r_\mathrm{E}/6\mathrm{AU}}\frac{ϵ}{v_{}/200\mathrm{km}\mathrm{s}^1},$$
$`(22)`$
where $`R_{}`$ is the radius of the source, and $`v_{}`$ is the transverse speed of the binary center of mass. That is, they are about equally likely. Note that only a small minority (roughly a fraction $`X`$) of the superluminal events occur in association with a normal event (where the source passes within the Einstein ring). The rest are isolated “spike events”. The real problem with observing superluminial events is not that they are uncommon, but that they are weak. For $`ϵ1`$ the caustic covers (and hence magnifies) only a small fraction of the source star. For example, for $`X=0.012`$, $`M=M_{}`$, $`D_l=D_{ls}=4`$kpc, and $`R=R_{}`$, the caustic covers only 0.01% of the source. Of course, for higher $`ϵ`$, the caustic area grows, but as we discuss in § 3.4, the strength of the caustic declines. Hence, at least for the present, superluminal caustics appear to be of mainly theoretical interest.
We thank Scott Gaudi for valuable discussions. This work was supported in part by grant AST 97-27520 from the NSF.
APPENDIX A
The lens equation can be written in terms of the normalized coordinates of points at the lens plane $`𝐫=(r_1,r_2)`$ and those at the source plane $`𝐱=(x_1,x_2)`$. It is convenient for calculation if we replace the integration variable $`t`$ in equation (8) with $`t^{}`$ using equation (10). Note that $`\mathrm{}<t^{}<0`$. According to the geometric relations in Figure 1, with the definition $`\tau =ct^{}/a`$, we then have the lens equation in component form:
$$x_i=r_i+\frac{1}{4}_{\mathrm{}}^0[F_{1i}(\tau )+F_{2i}(\tau )]Xd\tau ,$$
$`(A1)`$
where $`i=1,2`$ and
$$F_{i1}=\frac{4X\tau [r_1+(1)^iX\mathrm{cos}\beta \tau ]}{(1p_i)[(X\tau )^2+d_i^2(\tau )]^2}\frac{2(1)^i\beta \mathrm{sin}\beta \tau }{(1p_i)[(X\tau )^2+d_i^2(\tau )],}$$
$`(A2)`$
$$F_{i2}=\frac{4X\tau [r_2(1)^iX\mathrm{sin}\beta \tau ]}{(1p_i)[(X\tau )^2+d_i^2(\tau )]^2}\frac{2(1)^i\beta \mathrm{cos}\beta \tau }{(1p_i)[(X\tau )^2+d_i^2(\tau )],}$$
$`(A3)`$
$$d_i^2(\tau )=[r_1+(1)^iX\mathrm{cos}\beta \tau ]^2+[r_2(1)^iX\mathrm{sin}\beta \tau ]^2,$$
$`(A4)`$
$$p_i=(1)^{i1}\frac{2X\beta \tau }{(X\tau )^2+d_i^2(\tau )}(r_1\mathrm{sin}\beta \tau +r_2\mathrm{cos}\beta \tau ).$$
$`(A5)`$ |
warning/0001/hep-ph0001307.html | ar5iv | text | # References
INFNCA-TH0001
hep-ph/0001307
$`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ polarization in polarized DIS
M. Anselmino<sup>1</sup>, M. Boglione<sup>2</sup> and F. Murgia<sup>3</sup>
<sup>1</sup> Dipartimento di Fisica Teorica, Università di Torino and
INFN, Sezione di Torino, Via P. Giuria 1, I-10125 Torino, Italy
<sup>2</sup> Dept. of Physics and Astronomy, Vrije Universiteit Amsterdam,
De Boelelaan 1081, 1081 HV Amsterdam, The Netherlands
<sup>3</sup> Istituto Nazionale di Fisica Nucleare, Sezione di Cagliari
and Dipartimento di Fisica, Università di Cagliari
C.P. 170, I-09042 Monserrato (CA), Italy
Abstract:
We consider the polarization of $`\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}`$ baryons produced in polarized Deep Inelastic Scattering at leading order, with various spin configurations: longitudinally polarized leptons and unpolarized nucleon; unpolarized leptons and longitudinally or transversely polarized nucleons; longitudinally polarized leptons and nucleons. We show how the different results in the different cases are related to different aspects of the elementary dynamics and to the spin properties of the distribution and fragmentation functions and show how a combined analysis might give useful information. We give numerical results according to several sets of polarized fragmentation functions recently proposed.
1. Introduction $`\mathrm{\Lambda }`$ baryons produced in high energy interactions and resulting from quark fragmentation allow a unique test of spin transfer from partons to hadrons: the $`\mathrm{\Lambda }`$ polarization is easily measurable by looking at the angular distribution of the $`\mathrm{\Lambda }p\pi `$ decay (in the $`\mathrm{\Lambda }`$ helicity rest frame) and the fragmenting parton polarization is determined by the elementary Standard Model interactions, provided one knows the initial parton spin state. In this respect $`\mathrm{\Lambda }`$’s produced in lepton induced processes are particularly interesting and indeed several papers on this subject have recently been published or submitted to e-Print archives -.
We perform here a detailed analysis of the polarization of $`\mathrm{\Lambda }`$’s produced in polarized DIS; a general discussion of the helicity density matrix of hadrons produced in polarized lepton-nucleon interactions, at leading order, can be found in Refs. and . It can be written as:
$`\rho _{\lambda _h^{},\lambda _h^{}}^{(s,S)}(h){\displaystyle \frac{d\sigma ^{\mathrm{},s+N,S\mathrm{}+h+X}}{dxdydz}}`$ $`=`$ $`{\displaystyle \underset{q;\lambda _{\mathrm{}}^{},\lambda _q^{},\lambda _q^{}}{}}{\displaystyle \frac{1}{16\pi xs}}\times `$
$`\rho _{\lambda _{\mathrm{}}^{},\lambda _{\mathrm{}}^{}}^{\mathrm{},s}\rho _{\lambda _q^{},\lambda _q^{}}^{q/N,S}f_{q/N}(x)\widehat{M}_{\lambda _{\mathrm{}}^{},\lambda _q^{};\lambda _{\mathrm{}}^{},\lambda _q^{}}^q\widehat{M}_{\lambda _{\mathrm{}}^{},\lambda _q^{};\lambda _{\mathrm{}}^{},\lambda _q^{}}^qD_{\lambda _h^{},\lambda _h^{}}^{\lambda _q^{},\lambda _q^{}}(z),`$
where $`x`$ and $`y`$ are the usual DIS variables, $`x=Q^2/2pq`$, $`y=Q^2/xs`$ and, neglecting hadron masses, $`z=p_hp/pq`$, where $`p`$, $`q`$ and $`p_h`$ are, respectively, the nucleon, virtual photon and final hadron four-momenta. $`\rho ^{\mathrm{},s}`$ is the helicity density matrix of the initial lepton with spin $`s`$, $`f_{q/N}(x)`$ is the number density of unpolarized quarks $`q`$ with momentum fraction $`x`$ inside an unpolarized nucleon and $`\rho ^{q/N,S}`$ is the helicity density matrix of quark $`q`$ inside the polarized nucleon $`N`$ with spin $`S`$. The $`\widehat{M}_{\lambda _{\mathrm{}}^{},\lambda _q^{};\lambda _{\mathrm{}}^{},\lambda _q^{}}^q`$’s are the helicity amplitudes for the elementary process $`\mathrm{}q\mathrm{}q`$. The final lepton spin is not observed and helicity conservation of perturbative QCD and QED has already been taken into account in the above equation: as a consequence only the diagonal elements of $`\rho ^{\mathrm{},s}`$ contribute to $`\rho (h)`$ and non diagonal elements, present in case of transversely polarized leptons, do not contribute. $`D_{\lambda _h^{},\lambda _h^{}}^{\lambda _q^{},\lambda _q^{}}(z)`$ is a generalized fragmentation function related to the usual unpolarized fragmentation function $`D_{h/q}(z)`$, i.e. the density number of hadrons $`h`$ resulting from the fragmentation of an unpolarized quark $`q`$ and carrying a fraction $`z`$ of its momentum, by
$$D_{h/q}(z)=\frac{1}{2}\underset{\lambda _q^{},\lambda _h^{}}{}D_{\lambda _h^{},\lambda _h^{}}^{\lambda _q^{},\lambda _q^{}}(z)=\frac{1}{2}\underset{\lambda _q^{},\lambda _h^{}}{}D_{h_{\lambda _h^{}}/q_{\lambda _q^{}}}(z),$$
(2)
where $`D_{\lambda _h^{},\lambda _h^{}}^{\lambda _q^{},\lambda _q^{}}(z)D_{h_{\lambda _h^{}}/q_{\lambda _q^{}}}(z)`$ is a polarized fragmentation function, i.e. the density number of hadrons $`h`$ with helicity $`\lambda _h^{}`$ resulting from the fragmentation of a quark $`q`$ with helicity $`\lambda _q^{}`$. Angular momentum conservation and collinear configurations imply for the generalized fragmentation functions :
$$D_{\lambda _h^{},\lambda _h^{}}^{\lambda _q^{},\lambda _q^{}}=0\text{when}\lambda _q^{}\lambda _q^{}\lambda _h^{}\lambda _h^{}.$$
(3)
The elementary amplitudes $`\widehat{M}^q`$ are normalized so that the elementary unpolarized cross-section, $`\mathrm{}q\mathrm{}q`$, is:
$$\frac{d\widehat{\sigma }^q}{dy}e_q^2\frac{d\widehat{\sigma }}{dy}=\frac{1}{16\pi xs}\frac{1}{4}\underset{\lambda _{\mathrm{}}^{},\lambda _q^{}}{}|\widehat{M}_{\lambda _{\mathrm{}}^{},\lambda _q^{};\lambda _{\mathrm{}}^{},\lambda _q^{}}^q|^2=e_q^2\frac{2\pi \alpha ^2}{Q^2}\frac{1+(1y)^2}{y}$$
(4)
Notice that there are only two independent helicity amplitudes:
$`\widehat{M}_{++;++}^q=\widehat{M}_;^q`$ $``$ $`e_q\widehat{M}_{++;++}=e_q\mathrm{\hspace{0.25em}8}\pi \alpha {\displaystyle \frac{1}{y}}`$
$`\widehat{M}_{+;+}^q=\widehat{M}_{+;+}^q`$ $``$ $`e_q\widehat{M}_{+;+}=e_q\mathrm{\hspace{0.25em}8}\pi \alpha \left[{\displaystyle \frac{1}{y}}1\right]`$ (5)
so that
$$\frac{d\widehat{\sigma }^q}{dy}=\frac{1}{2}\left[\frac{d\widehat{\sigma }_q^{++}}{dy}+\frac{d\widehat{\sigma }_q^+}{dy}\right]$$
(6)
with
$$\frac{d\widehat{\sigma }_q^{++}}{dy}e_q^2\frac{d\widehat{\sigma }^{++}}{dy}=\frac{e_q^2}{16\pi xs}|\widehat{M}_{++;++}|^2=e_q^2\frac{4\pi \alpha ^2}{Q^2}\frac{1}{y}$$
(7)
and
$$\frac{d\widehat{\sigma }_q^+}{dy}e_q^2\frac{d\widehat{\sigma }^+}{dy}=\frac{e_q^2}{16\pi xs}|\widehat{M}_{+;+}|^2=e_q^2\frac{4\pi \alpha ^2}{Q^2}\frac{(1y)^2}{y}$$
(8)
Finally, the cross-section appearing in the l.h.s. of Eq. (S0.Ex1), which gives the correct normalization to $`\rho (h)`$, Tr$`\rho =1`$, can be written, using Eq. (3), as
$$\frac{d\sigma ^{\mathrm{},s+N,S\mathrm{}+h+X}}{dxdydz}=\underset{q;\lambda _{\mathrm{}}^{},\lambda _q^{}}{}\frac{1}{16\pi xs}\rho _{\lambda _{\mathrm{}}^{},\lambda _{\mathrm{}}^{}}^{\mathrm{},s}\rho _{\lambda _q^{},\lambda _q^{}}^{q/N,S}f_{q/N}(x)|\widehat{M}_{\lambda _{\mathrm{}}^{},\lambda _q^{};\lambda _{\mathrm{}}^{},\lambda _q^{}}^q(y)|^2D_{h/q}(z).$$
(9)
Eqs. (S0.Ex1)-(9) hold within QCD factorization theorem at leading twist and leading order in the coupling constants; the intrinsic $`𝒌_{}`$ of the partons have been integrated over and collinear configurations dominate both the distribution and the fragmentation functions. For simplicity of notations we have not indicated the $`Q^2`$ scale dependences in $`f`$ and $`D`$.
We shall use such equations for spin 1/2 $`\mathrm{\Lambda }`$ baryons produced starting from several particular initial spin configurations; we will discuss how the measurable components of the $`\mathrm{\Lambda }`$ polarization vector depend on different combinations of distribution functions, elementary dynamics and fragmentation functions: each of these terms predominantly depends on a single variable, respectively $`x`$, $`y`$ and $`z`$, and a careful analysis of different situations can yield precious information. Although $`\mathrm{\Lambda }`$ production in polarized DIS has been recently discussed in several papers, most of them only consider some initial spin configurations and specific models for fragmentation functions. Our analysis is more comprehensive and somewhat more general, emphasizing the physical meaning of possible measurements, and allowing also to obtain some general relationships between different polarization values.
2. Polarization vector of spin 1/2 baryons We fix our spin notations in the $`\mathrm{}p`$ center of mass frame: the lepton moves with four-momentum $`l`$ along the $`z`$-axis and the proton moves with four-momentum $`p`$ in the opposite direction; we choose $`xz`$ as the lepton-hadron production plane, with the $`y`$-axis parallel to $`𝒍\times 𝒑_h`$. We denote by $`S_L`$ the (longitudinal) nucleon spin oriented along the $`z`$-axis, by $`S_S`$ the (sideway) spin oriented along the $`x`$-axis and by $`S_N`$ the (normal) spin oriented along the $`y`$-axis. Notice that $`+S_L`$ corresponds to a $``$ helicity proton and $`S_L`$ to a $`+`$ helicity one; we will only consider longitudinally polarized leptons with spins $`\pm s_L`$ which correspond respectively to $`\pm `$ helicities.
From Eqs. (S0.Ex1)-(9) one obtains the explicit expression for the components of the helicity density matrix of a spin 1/2 baryon. These are related to the three components of the baryon polarization vector, as measured in its helicity rest frame by $`P_i(B)=`$ Tr$`[(\sigma ^i\rho (B)]`$ ($`i=x,y,z`$); details can be found in Ref. . We denote by $`(s,S)`$ \[or $`(h,H)`$\] the (lepton, nucleon) spins \[or helicities\], 0 stands for unpolarized particle; one finds, for a spin 1/2 hadron $`B`$:
$`P_z^{(0,S_L)}(B;x,y,z)`$ $`=`$ $`P_z^{(0,+)}(B;x,y,z)={\displaystyle \frac{_q\mathrm{\Delta }qd\widehat{\sigma }^q\mathrm{\Delta }D_{B/q}}{_qqd\widehat{\sigma }^qD_{B/q}}}`$ (10)
$`=`$ $`{\displaystyle \frac{_qe_q^2\mathrm{\Delta }q(x)\mathrm{\Delta }D_{B/q}(z)}{_qe_q^2q(x)D_{B/q}(z)}}`$
$`P_z^{(s_L,0)}(B;x,y,z)`$ $`=`$ $`P_z^{(+,0)}(B;x,y,z)={\displaystyle \frac{_q(1/2)q[d\widehat{\sigma }_q^{++}d\widehat{\sigma }_q^+]\mathrm{\Delta }D_{B/q}}{_qqd\widehat{\sigma }^qD_{B/q}}}`$ (11)
$`=`$ $`{\displaystyle \frac{_qe_q^2q(x)\mathrm{\Delta }D_{B/q}(z)}{_qe_q^2q(x)D_{B/q}(z)}}\widehat{A}_{LL}(y)`$
$`P_y^{(0,S_N)}(B;x,y,z)`$ $`=`$ $`P_x^{(0,S_S)}(B;x,y,z)={\displaystyle \frac{_q\mathrm{\Delta }_Tq[d\widehat{\sigma }_q^{}d\widehat{\sigma }_q^{}]\mathrm{\Delta }_TD_{B/q}}{_qqd\widehat{\sigma }^qD_{B/q}}}`$ (12)
$`=`$ $`{\displaystyle \frac{_qe_q^2\mathrm{\Delta }_Tq(x)\mathrm{\Delta }_TD_{B/q}(z)}{_qe_q^2q(x)D_{B/q}(z)}}\widehat{D}_{NN}(y)`$
$`P_z^{(s_L,S_L)}(B;x,y,z)`$ $`=`$ $`P_z^{(+,+)}(B;x,y,z)={\displaystyle \frac{_q[q_+d\widehat{\sigma }_q^{++}q_{}d\widehat{\sigma }_q^+]\mathrm{\Delta }D_{B/q}}{_q[q_+d\widehat{\sigma }_q^{++}+q_{}d\widehat{\sigma }_q^+]D_{B/q}}}`$ (13)
$`=`$ $`{\displaystyle \frac{_qe_q^2[q(x)\widehat{A}_{LL}(y)+\mathrm{\Delta }q(x)]\mathrm{\Delta }D_{B/q}(z)}{_qe_q^2[q(x)+\widehat{A}_{LL}(y)\mathrm{\Delta }q(x)]D_{B/q}(z)}}`$
$`P_z^{(s_L,+S_L)}(B;x,y,z)`$ $`=`$ $`P_z^{(+,)}(B;x,y,z)={\displaystyle \frac{_q[q_{}d\widehat{\sigma }_q^{++}q_+d\widehat{\sigma }_q^+]\mathrm{\Delta }D_{B/q}}{_q[q_{}d\widehat{\sigma }_q^{++}+q_+d\widehat{\sigma }_q^+]D_{B/q}}}`$ (14)
$`=`$ $`{\displaystyle \frac{_qe_q^2[q(x)\widehat{A}_{LL}(y)\mathrm{\Delta }q(x)]\mathrm{\Delta }D_{B/q}(z)}{_qe_q^2[q(x)\widehat{A}_{LL}(y)\mathrm{\Delta }q(x)]D_{B/q}(z)}}`$
where $`d\widehat{\sigma }_q`$ stands for $`d\widehat{\sigma }_q/dy`$.
Some comments are in order.
* The above results are known ; we have rederived and grouped them here for convenience and further discussion. Notice that all other spin configurations – at leading order – either give the same results or no polarization. We remind that $`q_\lambda =f_{q_\lambda /p_+}`$ is the number density of quarks with helicity $`\lambda `$ inside a + helicity proton, $`q(x)=q_+(x)+q_{}(x)`$ and $`\mathrm{\Delta }q(x)=q_+(x)q_{}(x)`$. Similarly $`\mathrm{\Delta }D_{B/q}=D_{B_+/q_+}D_{B_{}/q_+}`$; $`\mathrm{\Delta }_Tq`$ and $`\mathrm{\Delta }_TD_{B/q}`$ are respectively the analogue of $`\mathrm{\Delta }q`$ and $`\mathrm{\Delta }D_{B/q}`$ for transverse spins.
* The longitudinal $`B`$ polarization induced by a longitudinal nucleon polarization, Eq. (10), does not depend on the elementary dynamics, but only on the quark spin distribution and fragmentation properties. Neglecting the QCD $`Q^2`$-evolution, $`P_z^{(0,S_L)}`$ does not depend on the DIS variable $`y`$, but only on $`x`$ and $`z`$.
* The longitudinal $`B`$ polarization resulting from the scattering of longitudinally polarized leptons off unpolarized nucleons depends on the unpolarized distribution functions, the polarized fragmentation functions and the elementary dynamics, through the double spin asymmetry for the $`\mathrm{}q\mathrm{}q`$ process \[see Eqs. (6)-(8)\]:
$$\widehat{A}_{LL}(y)=\frac{d\widehat{\sigma }_q^{++}d\widehat{\sigma }_q^+}{d\widehat{\sigma }_q^{++}+d\widehat{\sigma }_q^+}=\frac{d\widehat{\sigma }_q^{++}d\widehat{\sigma }_q^+}{2d\widehat{\sigma }^q}=\frac{y(2y)}{1+(1y)^2}$$
(15)
Notice that $`\widehat{A}_{LL}`$ grows with $`y`$ from 0 (at $`y=0`$) to 1 (at $`y=1`$), so that $`P_z^{(s_L,0)}`$ is an increasing function of $`y`$, starting from 0 at $`y=0`$.
* The transverse $`B`$ polarization induced by a transverse nucleon polarization, Eq. (12), depends on the quark transverse spin distribution and fragmentation properties and on the elementary dynamics, through the double transverse spin asymmetry for the $`\mathrm{}q\mathrm{}q`$ process:
$`\widehat{D}_{NN}(y)`$ $`=`$ $`{\displaystyle \frac{d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}}{d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}+d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}}}`$ (16)
$`=`$ $`{\displaystyle \frac{d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}d\widehat{\sigma }^{\mathrm{}q^{}\mathrm{}q^{}}}{d\widehat{\sigma }^q}}={\displaystyle \frac{2(1y)}{1+(1y)^2}}`$
where $`=S_N`$ and $`=S_N`$.
Contrary to $`\widehat{A}_{LL}`$, $`\widehat{D}_{NN}`$ decreases with $`y`$, with $`\widehat{D}_{NN}=1`$ at $`y=0`$ and $`\widehat{D}_{NN}=0`$ at $`y=1`$; thus, $`P_y^{(0,S_N)}`$ is a decreasing function of $`y`$, reaching 0 at $`y=1`$. An experimental confirmation of the opposite $`y`$-dependences of $`P_z^{(s_L,0)}`$ and $`P_y^{(0,S_N)}`$ would supply a new, subtle and important test of the factorization scheme of Eq. (S0.Ex1). $`\widehat{D}_{NN}`$ is the transverse polarization of the final quark generated by an initial transversely polarized ($``$) quark in the $`\mathrm{}q\mathrm{}q`$ process, and it is usually referred to as the depolarization factor.
* When both the lepton and the nucleon are longitudinally polarized the $`B`$ resulting polarization depends on yet a different combination of polarized quark distribution functions, fragmentation functions and elementary dynamics. It is then clear why a combined study of $`P(B)`$ in different cases could yield unique information.
* Finally, we notice that, in general, $`P_i^{(s,S)}=P_i^{(s,S)}`$.
3. Polarization vector of $`\mathrm{\Lambda }`$ baryons
We now consider the particular case of $`\mathrm{\Lambda }`$ baryons and discuss possible ways of extracting information from combined measurements of $`P_i(\mathrm{\Lambda })`$; as we said the $`\mathrm{\Lambda }`$ polarization vector can be measured by looking at the proton angular distribution as resulting from $`\mathrm{\Lambda }\pi p`$ decay in the $`\mathrm{\Lambda }`$ helicity rest frame, that is the frame obtained by rotating the $`\mathrm{}p`$ c.m. frame around the $`y`$-axis so that the new $`z_h`$-axis is parallel to the $`\mathrm{\Lambda }`$ direction, and then boosting along $`z_h`$ with the same speed as the $`\mathrm{\Lambda }`$:
$`W(\theta _p,\varphi _p)`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}\left[1+\alpha (P_z\mathrm{cos}\theta _p+P_x\mathrm{sin}\theta _p\mathrm{cos}\varphi _p+P_y\mathrm{sin}\theta _p\mathrm{sin}\varphi _p)\right]`$ (17)
$`=`$ $`{\displaystyle \frac{1}{4\pi }}\left[1+\alpha 𝑷\widehat{𝒑}\right]`$
where $`\alpha =0.642\pm 0.013`$.
We follow Ref. and assume for the unpolarized fragmentation functions:
$$D_{\mathrm{\Lambda }/u}=D_{\mathrm{\Lambda }/d}=D_{\mathrm{\Lambda }/s}=D_{\mathrm{\Lambda }/\overline{u}}=D_{\mathrm{\Lambda }/\overline{d}}=D_{\mathrm{\Lambda }/\overline{s}}D_{\mathrm{\Lambda }/q}$$
(18)
where $`\mathrm{\Lambda }`$ means $`\mathrm{\Lambda }^0+\overline{\mathrm{\Lambda }}^0`$.
The heavy quark and gluon unpolarized fragmentation functions play a negligible role for $`z\text{ }>0.3`$ and we neglect them here: we have actually checked that our results, when comparable, are almost indistinguishable from those of Ref. where also heavy quark and gluon contributions to the unpolarized cross-sections are taken into account.
Similarly, we follow Ref. for the polarized fragmentation functions:
$$\mathrm{\Delta }D_{\mathrm{\Lambda }/u}(z,Q_0^2)=\mathrm{\Delta }D_{\mathrm{\Lambda }/d}(z,Q_0^2)=N_u\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z,Q_0^2).$$
(19)
Eqs. (19) holds also for light antiquarks and it remains valid through QCD $`Q^2`$-evolution; heavy quark contributions are neglected.
Using Eq. (18) and (19) into Eqs. (10)-(11) and (13)-(14) gives:
$`P_z^{(0,+)}(\mathrm{\Lambda };x,y,z)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }Q^{}(x)}{Q(x)}}{\displaystyle \frac{\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}}`$ (20)
$`P_z^{(+,0)}(\mathrm{\Lambda };x,y,z)`$ $`=`$ $`{\displaystyle \frac{Q^{}(x)}{Q(x)}}{\displaystyle \frac{\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}}\widehat{A}_{LL}(y)`$ (21)
$`P_z^{(+,+)}(\mathrm{\Lambda };x,y,z)`$ $`=`$ $`{\displaystyle \frac{Q^{}(x)\widehat{A}_{LL}(y)+\mathrm{\Delta }Q^{}(x)}{Q(x)+\mathrm{\Delta }Q(x)\widehat{A}_{LL}(y)}}{\displaystyle \frac{\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}}`$ (22)
$`P_z^{(+,)}(\mathrm{\Lambda };x,y,z)`$ $`=`$ $`{\displaystyle \frac{Q^{}(x)\widehat{A}_{LL}(y)\mathrm{\Delta }Q^{}(x)}{Q(x)\mathrm{\Delta }Q(x)\widehat{A}_{LL}(y)}}{\displaystyle \frac{\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}}`$ (23)
where
$`Q`$ $``$ $`4(u+\overline{u})+(d+\overline{d})+(s+\overline{s})`$ (24)
$`\mathrm{\Delta }Q`$ $``$ $`4(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})+(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})+(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s})`$ (25)
$`Q^{}`$ $``$ $`[4(u+\overline{u})+(d+\overline{d})]N_u+(s+\overline{s})`$ (26)
$`\mathrm{\Delta }Q^{}`$ $``$ $`[4(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})+(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})]N_u+(\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}).`$ (27)
If one assumes the same relation (19) to hold also for transversely polarized quark fragmentation functions, then Eq. (12) yields:
$$P_y^{(0,S_N)}(\mathrm{\Lambda };x,y,z)=\frac{\mathrm{\Delta }_TQ^{}(x)}{Q(x)}\frac{\mathrm{\Delta }_TD_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}\widehat{D}_{NN}(y)$$
(28)
with
$$\mathrm{\Delta }_TQ^{}[4(\mathrm{\Delta }_Tu+\mathrm{\Delta }_T\overline{u})+(\mathrm{\Delta }_Td+\mathrm{\Delta }_T\overline{d})]N_u+(\mathrm{\Delta }_Ts+\mathrm{\Delta }_T\overline{s}).$$
(29)
Eqs. (20)-(23) hold under assumptions (18) and (19) alone, independently of the actual value of $`N_u`$; as the polarized and unpolarized distribution functions are well known and several sets are available in the literature, we can exploit Eqs. (20) and (21) to obtain:
$$N_u=\frac{S}{U}\frac{\widehat{A}_{LL}P_z^{(0,+)}(\mathrm{\Delta }S/S)P_z^{(+,0)}}{\widehat{A}_{LL}P_z^{(0,+)}(\mathrm{\Delta }U/U)P_z^{(+,0)}}$$
(30)
and
$$\frac{\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z)}{D_{\mathrm{\Lambda }/q}(z)}=\frac{Q}{S}\frac{1}{\widehat{A}_{LL}}\frac{(\mathrm{\Delta }U/U)P_z^{(+,0)}\widehat{A}_{LL}P_z^{(0,+)}}{(\mathrm{\Delta }U/U)(\mathrm{\Delta }S/S)}$$
(31)
where we have defined
$$U4(u+\overline{u})+(d+\overline{d})Ss+\overline{s}\mathrm{\Delta }U4(\mathrm{\Delta }u+\mathrm{\Delta }\overline{u})+(\mathrm{\Delta }d+\mathrm{\Delta }\overline{d})\mathrm{\Delta }S\mathrm{\Delta }s+\mathrm{\Delta }\overline{s}.$$
(32)
Eqs. (22) and (23) can then be used to predict – within the general assumptions (18) and (19) – the following interesting relations between polarization observables:
$$P_z^{(+,\pm )}=\frac{P_z^{(+,0)}\pm P_z^{(0,+)}}{1\pm (\mathrm{\Delta }Q/Q)\widehat{A}_{LL}}$$
(33)
and
$$P_z^{(+,+)}P_z^{(+,)}=2\frac{P_z^{(0,+)}(\mathrm{\Delta }Q/Q)\widehat{A}_{LL}P_z^{(+,0)}}{1[(\mathrm{\Delta }Q/Q)\widehat{A}_{LL}]^2}$$
(34)
Notice that in the small $`y`$ region, due to the elementary dynamics, see Eq. (15), one has:
$$P_z^{(+,0)}0P_z^{(+,\pm )}\pm P_z^{(0,+)}(y1).$$
(35)
In the large $`y`$ region instead, again from Eq. (15) and from the fact that large $`y`$ implies small $`x`$, where $`Q(x)\mathrm{\Delta }Q(x)`$ and $`Q^{}(x)\mathrm{\Delta }Q^{}(x)`$, we expect:
$$P_z^{(+,+)}P_z^{(+,)}P_z^{(+,0)}(y1).$$
(36)
We conclude this Section by reminding that we have derived our results for $`\mathrm{\Lambda }=\mathrm{\Lambda }^0+\overline{\mathrm{\Lambda }}^0`$, which allows assumption (18) concerning $`q`$ and $`\overline{q}`$ fragmentation functions. However, Eqs. (20)-(23) and (28) hold identical also for single $`\mathrm{\Lambda }^0`$ production, provided one neglects the fragmentation function of a $`\overline{q}`$ into $`\mathrm{\Lambda }^0`$, i.e. one neglects all $`\overline{q}`$ terms in Eqs. (24)-(27) and (29). Anyway, the production of $`\overline{\mathrm{\Lambda }}^0`$, in $`\mathrm{}p`$ processes is strongly suppressed by the limited amount of initial $`\overline{q}`$, unless one considers very small $`x`$ values.
4. Numerical estimates
We give now some numerical estimates of Eqs. (20)-(23) and (28); we use the sets of unpolarized and polarized fragmentation functions introduced and discussed by the authors of Ref. : together with Eqs. (18), we use the expression for the unpolarized fragmentation functions they obtained by fitting $`e^+e^{}\mathrm{\Lambda }X`$ data. At initial $`Q_0^2=0.23`$ (GeV/$`c)^2`$ scale one has, from a leading order (LO) analysis :
$$D_{\mathrm{\Lambda }/q}(z,Q_0^2)=0.63z^{0.23}(1z)^{1.83}.$$
(37)
The polarized fragmentation functions are assumed to be of the initial form (19), with:
$$\mathrm{\Delta }D_{\mathrm{\Lambda }/s}(z,Q_0^2)=z^\alpha D_{\mathrm{\Lambda }/q}(z,Q_0^2).$$
(38)
Leading order QCD evolution is consistently taken into account in our numerical computations.<sup>1</sup><sup>1</sup>1We are very grateful to D. de Florian, M. Stratmann and W. Vogelsang for providing us with their FORTRAN package for unpolarized and polarized fragmentation functions. Next to leading order contributions to the $`\mathrm{\Lambda }`$ polarization have been shown to be tiny and we neglect them.
The parameter $`N_u`$ defined in Eq. (19), has been chosen according to three different scenarios typical of a wide range of plausible models, and the corresponding remaining parameter $`\alpha `$ of Eq. (38) has been fixed by fitting the few LEP data on $`\mathrm{\Lambda }`$ polarization, with the results :
1. $`𝑵_𝒖\mathbf{=}\mathrm{𝟎}\mathbf{,}𝜶\mathbf{=}\mathbf{0.62}`$. This scenario corresponds to $`SU(6)`$ non relativistic quark model, according to which the whole $`\mathrm{\Lambda }`$ spin is carried by the $`s`$ quark.
2. $`𝑵_𝒖\mathbf{=}\mathbf{}\mathbf{0.2}\mathbf{,}𝜶\mathbf{=}\mathbf{0.27}`$. Such value of $`N_u`$ is suggested in Ref. , based on a $`SU(3)`$ flavour symmetry analysis and on data on the first moment of $`g_1^p`$.
3. $`𝑵_𝒖\mathbf{=}\mathrm{𝟏}\mathbf{,}𝜶\mathbf{=}\mathbf{1.66}`$. A scenario in which, contrary to the non relativistic models, all light quarks contribute equally to the $`\mathrm{\Lambda }`$ polarization.
The unpolarized and polarized distribution functions are taken respectively from Refs. and . We have explicitely checked that a different choice of unpolarized and polarized distribution functions, like those of Refs. and , can change significantly the numerical values of $`P_i(\mathrm{\Lambda })`$ (but not their qualitative behaviour) only for $`x\text{ }>0.3`$.
A computation of $`P_y^{(0,S_N)}`$, Eq. (28), requires the knowledge of the quark transversely polarized distributions, $`\mathrm{\Delta }_Tq`$ or $`h_{1q}`$, and of the transversely polarized fragmentation functions, $`\mathrm{\Delta }_TD`$, which are not known. In order to give an estimate we fix $`\mathrm{\Delta }_Tq`$ for $`u`$ and $`d`$ quarks by saturating the Soffer’s bound (assuming the same signs for $`\mathrm{\Delta }_Tq`$ and $`\mathrm{\Delta }q`$):
$$\mathrm{\Delta }_Tu=\frac{1}{2}(u+\mathrm{\Delta }u)\mathrm{\Delta }_Td=\frac{1}{2}(d+\mathrm{\Delta }d).$$
(39)
All other transverse distributions ($`\mathrm{\Delta }_T\overline{q}`$ and $`\mathrm{\Delta }_Ts`$) are neglected here and we also assume $`\mathrm{\Delta }_TD_{\mathrm{\Lambda }/s}=\mathrm{\Delta }D_{\mathrm{\Lambda }/s}`$.
A sample of typical results is presented in Figs. 1-4, for HERMES kinematics, $`s=52.4`$ (GeV)<sup>2</sup> and $`Q^2\text{ }>1`$ (GeV/$`c)^2`$.
* In Fig. 1a we plot $`P_z^{(0,+)}`$ – at fixed $`x=0.1`$ and $`z=0.5`$ values – as a function of $`y`$, for each of the three scenarios; $`P_z^{(0,+)}`$, Eq. (20), can depend on $`y`$ only via the $`Q^2`$-evolution and indeed the three curves show an almost flat behaviour. The three scenarios yield quite different results. The minimum value of $`y`$ is given by $`y_{min}=Q_{min}^2/(xs)0.19`$.
The same plot for $`P_z^{(+,0)}`$ is presented in Fig. 1b; the $`y`$-dependence is essentially due to the factor $`\widehat{A}_{LL}`$ in Eq. (21); scenario 3, which assumes $`Q^{}(x)=Q(x)`$, together with Eq. (38) in which we neglect the mild $`Q^2`$-evolution (taken into account in our numerical computations), gives a particularly simple result:
$$P_z^{(+,0)}z^{1.66}\widehat{A}_{LL}.$$
(40)
* $`P_z^{(+,+)}`$ and $`P_z^{(+,)}`$, again as functions of $`y`$ at fixed $`x=0.1`$ and $`z=0.5`$ values, are shown in Figs. 2a and 2b respectively. At large $`y1`$ values Eq. (36) is satisfied. The small $`y0`$ behaviour cannot be seen with $`x=0.1`$; in Fig. 2c and 2d we plot respectively $`P_z^{(+,+)}`$ and $`P_z^{(+,)}`$, changing the $`x`$ value to $`x=0.3`$ and keeping $`z=0.5`$. The allowed minimum value of $`y`$ is now 0.06 and we have checked that Eq. (35) is indeed obeyed (by comparing with $`P_z^{(0,+)}`$ as a function of $`y`$ at $`x=0.3`$ and $`z=0.5`$, not shown in Fig. 1); the change in sign of $`P_z^{(+,)}`$ is particularly interesting.
* In Figs. 3a-3d we show respectively $`P_z^{(0,+)}`$, $`P_z^{(+,0)}`$, $`P_z^{(+,+)}`$ and $`P_z^{(+,)}`$, at fixed values of $`Q^2=1.7`$ (GeV/$`c)^2`$ and $`z=0.5`$, as functions of $`x`$. At fixed $`Q^2`$, $`y`$ decreases with increasing $`x`$ – and viceversa – and this explains why relations (35) and (36) hold respectively when $`x1`$ and $`x0`$. Once more, the three different scenarios give very different results. Using different sets of polarized distribution functions, like or instead of , gives almost identical results for $`x\text{ }<0.3`$ and larger (in magnitude) results for $`x\text{ }>0.3`$ (of course, $`P_z^{(+,0)}`$ is not affected at all by a change in $`\mathrm{\Delta }q`$).
* In Figs. 4a we plot $`P_y^{(0,S_N)}`$, Eq. (28), at fixed $`x=0.1`$ and $`z=0.5`$ values, as a function of $`y`$ for all scenarios; the $`y`$ dependence is almost entirely given by $`\widehat{D}_{NN}`$, Eq. (16), and indeed $`P_y^{(0,S_N)}0`$ when $`y1`$. Notice the opposite $`y`$ behaviour of $`P_z^{(+,0)}`$ and $`P_y^{(0,S_N)}`$ due to the opposite behavior of $`\widehat{A}_{LL}`$ and $`\widehat{D}_{NN}`$.
In Fig. 4b $`P_y^{(0,S_N)}`$ is plotted as a function of $`x`$ at fixed $`Q^2=1.7`$ (GeV/$`c)^2`$ and $`z=0.5`$.
* In general one finds very small values of $`P_i(\mathrm{\Lambda })`$ in scenario 1, negative values in scenario 2 and positive ones in scenario 3. This can easily be understood by the different values of $`N_u`$ in the three scenarios, which assign respectively zero, negative, and positive contributions to $`u`$ and $`d`$ quarks, which dominate in the proton. Experimental measurements can easily discriminate between them.
* Recently HERMES Collaboration have published a single experimental measurement of $`P_z^{(+,0)}/\widehat{A}_{LL}`$, as a function of $`z`$ and this seems to favour the scenario 1 prediction of Ref. , although errors and uncertainties are still too large to draw any reliable conclusions. Similarly, two values of $`P_z^{(+,0)}(z)`$ published by the E665 Collaboration still have much too large errors.
5. Conclusions The study of the angular distribution of the $`\mathrm{\Lambda }p\pi `$ decay allows a simple and direct measurement of the components of the $`\mathrm{\Lambda }`$ polarization vector. For $`\mathrm{\Lambda }`$’s produced in the current fragmentation region in DIS processes, the component of the polarization vector are related to spin properties of the quark inside the nucleon, to spin properties of the quark hadronization, and to spin dynamics of the elementary interactions. All this information, concerning quark distribution functions, quark fragmentation functions and spin properties of elementary dynamics are essentially factorized and separated as depending on three different variables, respectively $`x`$, $`z`$ and $`y`$. The $`Q^2`$-evolution and dependence of distribution and fragmentation functions somewhat mix the three variables, but smoothly, keeping separated the main properties of each of the different aspects of the process. Moreover, such $`Q^2`$ dependence is perturbatively well known and under control.
We have discussed all different polarization states of baryons, obtainable in the fragmentation of a quark in DIS with polarized initial leptons and nucleons, Eqs. (10)-(14), showing how they can reveal different quark features, weighted and shaped by elementary dynamics.
Adopting a simplifying – although rather general and representative of many possible choices – assumption about the quark fragmentation functions into a $`\mathrm{\Lambda }`$ , we are able to extract from measurements further information on the quark fragmentation process, Eqs. (30) and (31), and to predict relations among polarization states induced by different initial spin configurations, Eqs. (33) and (34).
Numerical estimates are given in Figs. 1-4, according to three largely different scenarios for fragmentation functions; each scenario has physical motivations and yields qualitatively different results: compatible with zero, large and negative, large and positive. Such results are stable against different choices of the polarized and unpolarized distribution functions, so that experimental information should immediately allow to draw clear conclusions and to learn about quark fragmentation properties.
The elementary dynamics fixes the small or large $`x`$ or $`y`$ behaviour of some of the polarization components; although expected, such behaviours should indeed be checked, as an independent and non trivial test of the QCD factorization scheme of Eq. (S0.Ex1); such a scheme has been widely used and tested for the computation of semi-inclusive unpolarized cross-sections, but not for more subtle spin observables.
We think that our comparative and comprehensive discussion of all possible $`\mathrm{\Lambda }`$ polarization measurements in polarized DIS is useful and can lead to a new and clear strategy which allows to obtain novel information.
Acknowledgements We would like to thank W. Vogelsang for useful discussions. |
warning/0001/astro-ph0001096.html | ar5iv | text | # Mapping the evolution of high redshift dusty galaxies with submillimeter observations of a radio-selected sample
## 1. Introduction
Recent detections of distant dusty galaxies with the SCUBA camera (the Submillimeter Common User Bolometer Array; Holland et al. 1999) on the 15 m James Clerk Maxwell Telescope<sup>1</sup><sup>1</sup>1The JCMT is operated by the Joint Astronomy Center on behalf of the parent organizations, the Particle Physics and Astronomy Research Council in the United Kingdom, the National Research Council of Canada, and the Netherlands Organization for Scientific Research. constitute a substantial fraction of the cosmic FIR background detected by the FIRAS and DIRBE experiments on the COBE satellite (Puget et al. 1996; Guiderdoni et al. 1997; Schlegel, Finkbeiner, & Davis 1998; Fixsen et al. 1998; Hauser et al. 1998; Lagache et al. 1999). Since the observed FIR background is comparable to the total unobscured emission at ultraviolet/optical wavelengths, a full determination of the global star formation history of the Universe requires a comprehensive understanding of this dust-enshrouded galaxy population. The 850 $`\mu `$m SCUBA surveys to date have reported galaxy number counts that are in general agreement (Smail, Ivison & Blain 1997; Barger et al. 1998; Hughes et al. 1998; Blain et al. 1999a; Eales et al. 1999; Barger, Cowie, & Sanders 1999); the cumulative surface density above 2 mJy is about $`3\times 10^3\mathrm{deg}^2`$. The discrete sources have bolometric luminosities that are characteristically $`10^{12}h_{65}^2\mathrm{L}_{}`$ if they lie at $`z1`$. Moreover, the sources for which measurements exist at multiple wavelengths (e.g. Ivison et al. 1998) show thermal spectral energy distributions (SEDs). Thus, the SCUBA sources are inferred to be the distant analogs of the local ultraluminous infrared galaxy (ULIG; Sanders & Mirabel 1996) population.
An essential observational goal is to determine the redshift distribution of the submillimeter population in order to trace the extent and evolution of obscured emission in the distant Universe. However, identifying the optical/NIR counterparts to the submillimeter sources is difficult due to the uncertainty in the SCUBA positions. Barger et al. (1999b) presented a spectroscopic survey of possible optical counterparts to a flux-limited sample of galaxies selected from the 850 $`\mu `$m survey of massive lensing clusters by Smail et al. (1998). Candidate optical counterparts in the SCUBA error-boxes were identified using moderately deep ground-based and HST exposures ($`I23.5`$ and $`I26`$, respectively). One-quarter of the sources could be reliably identified, and those had redshifts in the range $`z13`$. A lower limit of 20 per cent of the full sample showed signs of AGN activity. However, for the majority of the submillimeter sources there were either no optical counterparts or the optical associations were not secure. Such sources could either be at very high redshift or be so highly obscured that they emit their energy almost entirely in the submillimeter.
High resolution radio continuum maps with subarcsecond positional accuracy and resolution offer new opportunities for locating submillimeter sources and determining their physical properties. The unique advantage of centimeter and FIR observations is that galaxies and the intergalactic medium are transparent at these wavelengths, so observed flux densities are proportional to intrinsic luminosities. In galaxies without a powerful AGN, the radio luminosity is dominated by diffuse synchrotron emission from relativistic electrons accelerated in supernovae remnants from stars more massive than $`8\mathrm{M}_{}`$. These massive stars live $`3\times 10^7`$ yr; the relativistic electrons probably live $`10^8`$ yr (Condon 1992). Thus, radio observations probe very recent star formation. FIR observations of starburst galaxies are also a direct measure of massive star formation. As summarized by Condon (1992), radio continuum emission and thermal dust emission are empirically observed to be tightly correlated due to both being linearly related to the massive star formation rate. If the FIR-radio correlation applies to high redshift, as is plausibly the case (though at the very highest redshifts Compton cooling of the relativistic electrons by the microwave background may suppress the radio emission), then very sensitive radio observations can be used to pinpoint distant submillimeter sources.
In this paper we investigate the feasibility of using radio data to identify and characterize the bright submillimeter source population. Richards (1999b) recently obtained an extremely deep Very Large Array (VLA) 1.4 GHz image centered on the Hubble Deep Field (HDF). Richards et al. (1999) matched ground-based optical data from Barger et al. (1999a) to the radio image and found that $`20`$ per cent of the galaxies in the sample could not be identified to optical magnitude limits of $`I25`$. In other respects, such as radio size and spectral index, the optically-faint objects were not any different from the remaining population. Richards et al. (1999) proposed four possible scenarios to explain this population, including $`1<z<3`$ obscured starbursts (beyond $`z3`$ the sensitivity to star forming galaxies cuts off due to the flux density limits of the radio data), extreme redshift ($`z>6`$) AGN, $`z<2`$ obscured AGN, or one-sided radio jets.
In the first phase of our program we observed with LRIS on the Keck II 10 m telescope a complete subsample of the radio sources. Our primary objective was to determine the redshifts of the optical/NIR-faint ($`HK^{}>20.5`$) radio sources. Although we were able to spectroscopically identify nearly all the objects in our subsample to $`HK^{}20`$ (all had $`z1.3`$), we were unable to obtain redshifts for the fainter objects. Either these sources are distant ($`z>1.5`$) with spectral features lying outside the optical wavelength range, or the visibility of remarkable features is strongly affected by dust in the galaxies.
If the optical/NIR-faint radio sources are highly dust obscured systems, then it is possible that they will be detectable in the submillimeter. In the second phase of our program we observed with SCUBA 15 of the 22 optical/NIR-faint radio sources in the central $`80`$ square arcminute region of the radio map; another 4 were observed by Hughes et al. (1998; hereafter H98) in the HDF-proper. The jiggle map mode enabled simultaneous observations of a large fraction (31/48) of the optical/NIR-bright radio sources. Even with relatively shallow SCUBA observations (a $`3\sigma `$ detection limit of 6 mJy at 850 $`\mu `$m), we detected 5 of the optical/NIR-faint radio sources; a sixth with submillimeter flux $`<6`$ mJy was detected in the deep HDF-proper submillimeter map of H98. In contrast, none of the optical/NIR-bright radio sources were detected. We additionally detected two $`>6`$ mJy sources that did not have radio counterparts. Thus, our targeted SCUBA survey of optical/NIR-faint radio sources turned up $`70`$ per cent of the bright submillimeter sources in our surveyed areas.
In the final phase of our program, we explored the feasibility of obtaining redshift estimates from the submillimeter-to-radio flux ratios, as recently suggested by Carilli & Yun (1999). We find that the redshifted Arp 220 SED reasonably describes both redshifted local ULIG data and known high redshift submillimeter sources and hence can be used as a rough redshift estimator. We estimate that all of our bright submillimeter sources fall in the redshift range $`z=13`$, consistent with the redshifts for the lensed submillimeter sources of Barger et al. (1999b).
Once we have redshifts for the distant submillimeter source population, we can determine the global evolution of star formation in dust-obscured galaxies. Previous studies of the star formation rate density (SFRD) have primarily used rest-frame ultraviolet data (e.g. Madau et al. 1996). However, the ultraviolet emission from a galaxy is heavily affected by the presence of even small amounts of dust, and the extinction corrections are highly uncertain; for example, corrections for dust obscuration at $`z3`$ range from factors of $`3`$ (Pettini et al. 1997) to factors of $`15`$ (Meurer et al. 1997), though the more recent estimate of Meurer, Heckman, & Calzetti (1999) lowers the latter to a factor of $`5`$. Only direct measurements of the reradiated light at submillimeter wavelengths can securely address the SFRD at high redshifts.
The true SFRD will receive contributions from both the ultraviolet/optical and the submillimeter. We find that the submillimeter contribution to the SFRD in the $`z=13`$ range from our $`>6`$ mJy sources is comparable to the observed ultraviolet/optical SFRD contribution. The ultraviolet Lyman-break galaxies are on the average undetected in the submillimeter at a $`1\sigma `$ level of $`0.5`$ mJy, and thus the Lyman-break galaxy population is largely distinct from the bright submillimeter population (Chapman et al. 1999). If we assume that fainter submillimeter sources have the same redshift distribution and properties as the $`>6`$ mJy sample, then the SFRD of the entire population contributing to the submillimeter background is about an order of magnitude higher than the observed ultraviolet/optical SFRD. The contribution from ULIGs to the SFRD increases by more than two orders of magnitude from $`z0`$ to $`z=13`$, which supports a scenario in which the distant submillimeter sources are the progenitors of massive spheroidal systems. This rapid evolution in the ULIG population sampled by SCUBA can be reproduced in models (e.g., Blain et al. 1999b,c).
In § 2 we present our radio sample and optical/NIR imaging, along with our new SCUBA and LRIS observations. In § 3 we determine our radio-selected submillimeter source counts. In § 4 we introduce the predicted high redshift submillimeter-radio flux correlation and obtain millimetric redshift estimates for our sources. In § 5 we use the complementary information from radio and submillimeter fluxes to gain insights into the characteristics of our radio-selected submillimeter source population. In § 6 we compare the rest-frame colors of the radio sources that have spectroscopic redshifts with the rest-frame colors of our submillimeter sources and find that the latter are likely to fall in the extremely red object category. In § 7 we calculate the luminosities, number densities, and SFRDs of our submillimeter sources. We compare the submillimeter contribution to the SFRD with contributions from radio and ultraviolet/optical wavebands over a range of redshifts. In § 8 we summarize our main conclusions. We take $`\mathrm{H}_\mathrm{o}=65h_{65}`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and consider both $`\mathrm{\Omega }_\mathrm{M}=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_\mathrm{M}=1/3`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$, which should cover the full range of possible cosmologies.
## 2. Samples and Observations
The present study is based on deep radio maps centered on the HDF that were observed with the VLA at 1.4 GHz (Richards 1999b) and 8.5 GHz (Richards et al. 1998). The primary 1.4 GHz image covers a $`40^{}`$ diameter region with an effective resolution of $`1.8^{\prime \prime }`$ and a $`5\sigma `$ completeness limit of 40 $`\mu `$Jy. The 8.5 GHz images have an effective resolution of $`3.5^{\prime \prime }`$ and a $`5\sigma `$ completeness limit of 8 $`\mu `$Jy over a radius of $`1^{}`$ from the HDF center, rising to 40 $`\mu `$Jy at $`6.6^{}`$.
The 1.4 GHz HDF map was trimmed to match the 79.4 arcmin<sup>2</sup> NIR and optical imaging of the field obtained by Barger et al. (1999a). The absolute radio positions are known to $`0.1^{\prime \prime }0.2^{\prime \prime }`$ rms in the HDF; the alignment of the optical data to the radio data left residual astrometric uncertainties of $`0.2^{\prime \prime }`$ (Richards et al. 1999). There are 70 sources in our final radio sample.
Table 1 gives the 1.4 GHz radio catalog for the central 79.4 arcmin<sup>2</sup> region discussed above. The first five columns are catalog number, RA(2000), Dec(2000), 1.4 GHz radio flux, and radio flux uncertainty (Richards 1999b). The remaining columns, to be discussed in the following subsections, are $`HK^{}`$, $`I`$, $`V`$, $`R`$, $`B`$, and $`U^{}`$ magnitudes, redshift, submillimeter flux, submillimeter flux uncertainty, and $`6.75`$ and 15 $`\mu `$m ISOCAM fluxes (Aussel et al. 1999). The last column of Table 1 gives the radio spectral index $`\alpha _r`$ ($`S\nu ^{\alpha _r}`$) or limit for the sources where 8.4 GHz data are available and can be used to determine the origin of the radio emission (Richards 1999b). Inverted spectrum sources ($`\alpha _r>0`$) invariably have self-absorbed synchrotron emission associated with an AGN. Flat spectrum sources ($`0.5<\alpha _r<0`$) can be produced by AGN activity or by optically thin Bremsstrahlung radiation from star formation at higher radio frequencies ($`\nu >5`$ GHz). Steep spectrum sources ($`\alpha _r<0.5`$) consist of diffuse synchrotron emission often associated with either radio jets or star formation in galaxies.
In this paper we use the radio spectral index as a crude discriminator between AGN and star formation activity. We arbitrarily classify any source with an available spectral index that has $`\alpha _r>0.3`$ (there are six such sources in our sample) to be primarily powered by AGN activity.
### 2.1. Optical, Near-infrared, and Mid-infrared Imaging
Wide-field and deep $`HK^{}`$ observations of the HDF and flanking fields were obtained using the University of Hawaii Quick Infrared Camera (QUIRC; Hodapp et al. 1996) on the 2.2 m University of Hawaii (UH) telescope and the 3.6 m Canada-France-Hawaii Telescope (CFHT). The $`HK^{}`$ ($`1.9\pm 0.4`$ $`\mu `$m) filter is described in Wainscoat & Cowie (in preparation); the empirical relation between $`HK^{}`$ and $`K`$ is $`HK^{}K=0.13+0.05(IK)`$, which simplifies to $`HK^{}K=0.3`$, assuming the median $`IK`$ galaxy color (Barger et al. 1999a). In Table 1 we use a $`2\sigma `$ limit of $`HK^{}=21.5`$ for our wide-field image and a $`2\sigma `$ limit of $`HK^{}=22.6`$ for our deep image in the area around the HDF proper.
Deep Johnson $`V`$ and Kron-Cousins $`I`$-band observations were made with the CFHT over a much larger area using the UH8K CCD Mosaic Camera built by Metzger, Luppino, and Miyazaki. Details of the above $`HK^{}`$ and optical observations can be found in Barger et al. (1999a). In Table 1 we use $`2\sigma `$ limits of $`I=25.3`$ and $`V=26.4`$.
In February 1998 we observed four of the sources in our radio sample with the near-infrared camera (NIRC; Matthews & Soifer 1994) on the Keck I 10 m telescope. These sources previously had only $`HK^{}`$ limits but were detected in the NIRC observations. Conditions were photometric with seeing $`0.7^{\prime \prime }`$ FWHM. NIRC has a $`256\times 256`$ InSb array with $`0.15^{\prime \prime }\times 0.15^{\prime \prime }`$ pixels, giving a $`38^{\prime \prime }\times 38^{\prime \prime }`$ field of view. We imaged at 2.1 $`\mu `$m ($`K^{}`$) with total exposure times for each object of $`32404320`$ s. The data were obtained in sets of 120 s exposures, and the center of the field was moved in a $`3\times 3`$ grid pattern with $`3^{\prime \prime }`$ on a side. The centers of successive grids were moved by $`2^{\prime \prime }`$ between each set. A fifth source was subsequently detected with NIRC with a longer exposure under non-photometric conditions. The data were processed using median sky flats generated from the dithered images and calibrated onto the $`HK^{}`$ images using other galaxies in the field.
We used the Low-Resolution Imaging Spectrometer (LRIS; Oke et al. 1995) on the Keck II 10 m telescope in March 1997 and February 1998 to obtain $`B`$-band and Kron-Cousins $`R`$-band images, respectively, of a strip region $`6^{}\times 2.5^{}`$ in size that crosses the HDF. The total exposure times were 1680 s and 1600 s for $`B`$ and $`R`$, respectively, and the seeing was $`0.8^{\prime \prime }`$ in $`R`$ and $`1.3^{\prime \prime }`$ in $`B`$. The $`2\sigma `$ limits are $`B=26.6`$ and $`R=26.6`$. Details of the observations can be found in Cowie & Hu (1998) and Cowie, Songaila, & Barger (1999).
Finally, a deep $`U^{}(3400\pm 150`$Å) image of an area $`80`$ arcmin<sup>2</sup> centered on the HDF-proper was obtained using the ORBIT CCD on the UH telescope. The $`2\sigma `$ limit for the $`U^{}`$ image is $`U^{}=25.8`$. Details of the $`U^{}`$ observations can be found in Wilson et al. (in preparation).
All photometric magnitudes were measured in $`3^{\prime \prime }`$ diameter apertures and then corrected to $`6^{\prime \prime }`$ diameter (near total) magnitudes following the procedures of Cowie et al. (1994).
Serjeant et al. (1997) obtained deep ISOCAM observations of the HDF and flanking fields at 6.75 and 15 $`\mu `$m. The Aussel et al. (1999) reductions of the data produced a main source list of 49 objects ($`7\sigma `$) and a supplementary list of an additional 51 objects ($`3\sigma `$ for 15 $`\mu `$m and $`5\sigma `$ for 6.75 $`\mu `$m). Despite the large point spread function ($`15^{\prime \prime }`$ at 15 $`\mu `$m), in most cases the optical or NIR identification of the ISOCAM source was straightforward. The redshifts for sources in the sample are in the range $`z=0.078`$ to $`z=1.242`$ (median $`z=0.585`$).
### 2.2. Keck Spectroscopy
We used LRIS during a two night run on Keck II in March 1999 to obtain spectroscopic observations of a sub-sample of the 1.4 GHz sources in a strip region centered on the HDF. Although the HDF and flanking field region has been intensely studied by a number of groups (see Cohen et al. 1999 for a summary and references), most of the sources in our sample had not been previously observed due to their faint optical fluxes. We used $`1.4^{\prime \prime }`$ wide slits and the 400 lines mm<sup>-1</sup> grating blazed at 8500 Å, which gives a wavelength resolution of $`12`$ Å and a wavelength coverage of $`4000`$ Å. The wavelength range for each object depends on the exact location of the slit in the mask but is generally between $`4000`$ and 10000 Å. Three of the slit masks were constructed at a position angle of $`90^{}`$, and the remaining three were nearly identical versions constructed at a position angle of $`90^{}`$. This procedure enabled us to obtain sufficient wavelength coverage for all the objects in our sample, including those that fell close to the edges of the masks. The observations were 1.5 hr per slit mask, broken into three sets of 30 minute exposures. Three HDF slit masks were observed per night. Some of the objects were in all of the slit masks and hence were observed for 9 hrs (see Table 2). Conditions were photometric with seeing $`0.6^{\prime \prime }0.7^{\prime \prime }`$ FWHM both nights. The objects were stepped along the slit by $`10^{\prime \prime }`$ in each direction, and the sky backgrounds were removed using the median of the images to avoid the difficult and time-consuming problems of flat-fielding LRIS data. Details of the spectroscopic reduction procedures can be found in Cowie et al. (1996).
In Table 2 we list all the objects in the radio sample that fall in the LRIS strip region. The columns are radio catalog number from Table 1, $`HK^{}`$ mag, $`I`$ mag, redshift, and exposure time for any object targeted in our March 1999 spectroscopic survey. Over an area $`58`$ arcmin<sup>2</sup>, 19 of the 37 radio sources now have secure redshift identifications; all are at $`z1.3`$. Figure 1a, b shows $`HK^{}`$ versus redshift and $`I`$ versus redshift. Spectroscopic redshifts are in general relatively straightforward to obtain for radio objects with $`HK^{}20`$.
We note that Waddington et al. (1999) claim a redshift identification of $`z=4.42`$ for object 30 based on a Ly$`\alpha `$ detection. However, the position of their detection is $`1^{\prime \prime }`$ away from the radio source position and counterpart optical galaxy, and so an association with the radio source is not secure. The radio source is detected in the mid-IR, which would suggest $`z<1.3`$, consistent with the redshift estimated using the millimetric redshift technique described in § 4. However, we did not see any strong \[O II\] 3727 feature in our spectrum from $`40009200`$ Å, which would suggest $`z>1.5`$. Thus, the redshift of this object remains uncertain.
In all of our slit-masks (equivalent to a 9 hr integration) we also included the position of the brightest SCUBA source, HDF850.1, from the deep 850 $`\mu `$m map of the HDF-proper by H98. We centered our slit across the position of the 1.3 mm detection of HDF850.1 reported by Downes et al. (1999), which also coincides with the 8.5 GHz supplemental radio source 3651+1226 (Richards et al. 1998). We oriented the slit such that it fell across both the arc-like feature 3-593.0 that is favored by Downes et al. (1999) as the optical counterpart to HDF850.1 and the nearby red galaxy 3-586.0. However, we were unable to determine a secure redshift for either source from our spectroscopic data.
### 2.3. Submillimeter Observations
Our SCUBA jiggle map observations were flexibly scheduled in mostly excellent observing conditions during two runs in April and June 1999 for a total of five observing shifts. The maps were dithered to prevent any regions of the sky from repeatedly falling on bad bolometers. The chop throw was fixed at a position angle of $`90^{}`$ so that the negative beams would appear 45 arcsec on either side east-west of the positive beam. The data were reduced using beam weighted extraction routines that included both the positive and negative portions of the chopped images, thereby increasing the effective exposure times. Regular “skydips” (Lightfoot et al. 1998) were obtained to measure the zenith atmospheric opacities at 450 and 850 $`\mu `$m, and the 225 GHz sky opacity was monitored at all times to check for sky stability. The median 850 $`\mu `$m optical depth for all nights together was 0.265. Pointing checks were performed every hour during the observations on the blazars 0954+685, 1418+546, 0923+392, or 1308+326. The data were calibrated using 30 arcsec diameter aperture measurements of the positive beam in beam maps of the primary calibration source, Mars, and one of three secondary calibration sources, CRL618, IRC+10216, or OH231.8.
The data were reduced in a standard and consistent way using the dedicated SCUBA User Reduction Facility (SURF; Jenness & Lightfoot 1998). Due to the variation in the density of bolometer samples across the maps, there is a rapid increase in the noise levels at the very edges. We have clipped the low exposure edges from our images. We present our SCUBA maps in Fig. 2.
We also re-reduced the archival HDF-proper SCUBA data of H98 in order to make a consistent analysis with the present data. We could only make use of the maps that were taken with a fixed RA chop (90 per cent of the data sample), which included 34 hours with a 30 arcsec chop throw and 28 hours with a 45 arcsec chop throw. We combined these data separately to form two independent maps. In Table 1 we quote submillimeter fluxes determined from the weighted average of measurements made in each map.
The SURF reduction routines arbitrarily normalize all the data maps in a reduction sequence to the central pixel of the first map; thus, the noise levels in a combined image are determined relative to the quality of the central pixel in the first map. In order to determine the absolute noise levels of our maps, we first eliminated the $`3\sigma `$ real sources in each field by subtracting an appropriately normalized version of the beam profile. We then iteratively adjusted the noise normalization until the dispersion of the signal-to-noise ratio measured at random positions became $`1`$. Our noise estimate includes both fainter sources and correlated noise.
We centered on the positions of the radio sources and measured the submillimeter fluxes in both the positive and negative beams. The extracted fluxes were calibrated to the 30 arcsec diameter aperture fluxes of the brightest sources. The brighter submillimeter sources were subtracted from the maps before the fainter sources were extracted. In cases where the positions of two radio sources are very close, we may be slightly overestimating the submillimeter flux allocated to the bright source with this procedure and underestimating that allocated to the faint source, but since the total submillimeter flux should be reasonable, we will not be making any gross errors in our later estimates of the total star formation rate.
From the total sample of 70 radio-selected galaxies in the $`80`$ square arcminute central region of the radio map, we take the 22 with $`HK^{}>20.5`$ to be our optical/NIR-faint radio sample, for which there are now submillimeter observations of 19. Even though our SCUBA observations were relatively shallow (the 850 $`\mu `$m $`3\sigma `$ limit was 6 mJy), we detected 5 of the 15 optical/NIR-faint radio sources that we observed; a sixth significant source ($`S_{850\mu \mathrm{m}}=2.4\pm 0.7`$ mJy) was detected in the HDF-proper SCUBA map of H98 (HDF850.2 in their notation). The submillimeter flux we measure for the brightest source (HDF850.1) in the HDF-proper map is $`5.4\pm 0.6`$ mJy, which is slightly lower than the value of $`7.0\pm 0.4`$ found by H98 but is consistent within the statistical and systematic errors. The jiggle map observing mode enabled simultaneous observations of a large fraction (31/48) of the optical/NIR-bright radio sources, none of which were detected. There are two $`>6`$ mJy submillimeter sources in our jiggle maps that were not in the radio sample.
In Fig. 3a, b we show the 1.4 GHz fluxes of the radio sources in our sample versus their $`HK^{}`$ and $`I`$ magnitudes, indicating those with measured redshifts and those with submillimeter detections.
## 3. Submillimeter Source Counts
We find that the radio selection technique is effective in locating the majority of bright submillimeter sources. We document this in Fig. 4 where we compare the combined 850 $`\mu `$m source counts from blank field submillimeter surveys (H98; Eales et al. 1999; Barger, Cowie, & Sanders 1999) with our new radio-selected 850 $`\mu `$m source counts. A source of flux strength $`S_i`$ contributes $`N(S_i)=1/A_i`$ to the counts per unit area, where $`A_i`$ is the area over which there is $`3\sigma `$ sensitivity to $`S_i`$. The cumulative counts, $`N(>S)`$, are given by the sum of the inverse areas of all sources brighter than $`S`$. However, for our radio-targeted search, the effective area (65 square arcminutes) is greater than if we had done an untargeted search with the same number of pointings and is given by the fraction (number $`HK^{}>20.5`$ sources observed)/(total number $`HK^{}>20.5`$ sources in sample) times the radio survey area (79.4 square arcminutes). Because of the correspondence between the optical/NIR-faint radio population and the bright submillimeter population, we do not need to survey the entire radio field in order to observe the bright submillimeter population.
It is interesting to speculate whether the correspondence between optical/NIR-faint radio sources and bright submillimeter sources also holds at fainter submillimeter flux levels. A necessary requirement for this to be the case is that the surface density of faint radio sources be comparable to or exceed the surface density of submillimeter sources. Extrapolating a power-law parameterization of the 1.4 GHz counts over the range $`401000`$ $`\mu `$Jy (Richards 1999b) down to a 1 $`\mu `$Jy cut-off gives $`N(S_{1.4}>1\mu `$Jy) $`=127`$ arcmin<sup>-2</sup>. The empirically estimated surface density of submillimeter sources (Barger, Cowie, & Sanders 1999) is $`N(S_{353}>0)11`$ arcmin<sup>-2</sup>.
A possible counter-argument to a close correspondence between radio and submillimeter sources at faint submillimeter flux levels comes from the study by Smail et al. (1999b) who compared their lensed submillimeter survey data with radio data. They found that the Carilli & Yun (1999) radio-submillimeter spectral indices differed for their bright and faint submillimeter subsamples. They suggested that the difference could arise if the faint submillimeter sources lie at higher redshifts or have different fractions of radio-loud AGN or have different dust temperatures. However, with the low statistics, strong conclusions cannot be drawn at this time.
## 4. Millimetric Redshift Estimation
Optical spectroscopic surveys of submillimeter galaxies are difficult because of the very different behaviors of the $`K`$-corrections in the optical and submillimeter and the poor submillimeter resolution. Spectroscopic surveys to date (Barger et al. 1999b; Lilly et al. 1999) give potentially conflicting results and are very limited in size. Photometric redshift estimates of submillimeter sources have been made from candidate optical counterparts (H98; Lilly et al. 1999), but these are questionable both because the counterpart identification is not secure and because the SEDs of the submillimeter sources may be unlike those of the optical sources used in making the redshift estimates due to dust extinction and possible AGN contributions. Consequently, a new approach is required if the positions of the submillimeter sources are to be reliably determined and their nature and redshift distribution understood.
The remarkably tight local correlation between the global FIR and nonthermal radio luminosities provides a promising alternative method for identifying and studying individual submillimeter sources, provided that the correlation holds to high redshift. In § 3 we found that targeting with SCUBA optical/NIR-faint 1.4 GHz sources is an efficient technique for identifying the majority of the bright submillimeter source population. Our results indicate that a large fraction of bright sources in submillimeter surveys have extremely faint optical/NIR counterparts and hence are inaccessible to optical spectroscopy (see also Smail et al. 1999b). This conclusion is consistent with results from the Barger et al. (1999b) spectroscopic survey of lensed submillimeter sources discussed in the introduction.
Although we are unable to obtain spectroscopic redshifts for the optical/NIR-faint radio-selected submillimeter sources, we can use the submillimeter-to-radio flux ratios to obtain millimetric redshift estimates. Figure 5 illustrates how the slope of a dusty galaxy’s SED changes abruptly at frequencies higher than 100 GHz. Below 30 GHz synchrotron radio emission is dominant, free-free emission is largest in the range $`30200`$ GHz, and thermal dust emission dominates above 200 GHz. Because of the opposing spectral slopes of the blackbody spectrum in the submillimeter and the synchrotron spectrum in the radio, the submillimeter-to-radio flux ratio rises extremely rapidly with increasing redshift (see Fig. 6). Carilli & Yun (1999; hereafter CY99) have therefore suggested using the submillimeter-to-radio flux ratio as a redshift estimator.
The dust emission from local luminous infrared galaxies is well described by optically thin single-temperature modified blackbodies with extinction coefficient $`ϵ_\nu \nu ^\beta `$, where $`\beta 1`$ to 2. The dust temperatures derived lie between 30 and 60 K. Arp 220, with its high bolometric luminosity produced almost entirely from starburst activity (Downes & Solomon 1998), is an appropriate prototype for high redshift submillimeter sources, as we will justify subsequently. Over the interval $`100\mathrm{GHz}\nu 10^4`$ GHz ($`3000\mu \mathrm{m}\lambda 30\mu `$m) Arp 220’s SED is well represented by a modified blackbody with $`\beta =1`$ and a dust temperature $`T_d=47`$ K; the luminosity of this dominant cooler component is $`1.36\times 10^{12}h_{65}\mathrm{L}_{}`$ (Klaas et al. 1997). In terms of the blackbody distribution
$$B(\nu ,T)=\frac{S(\nu ,T)}{\pi }=\frac{2h\nu ^3}{c^2}\frac{1}{e^{h\nu /kT}1}$$
(1)
the submillimeter flux for an Arp 220-like galaxy at redshift $`z`$ is given by
$$S_s=f(0.57\mathrm{Jy})\pi B(\nu (1+z),T_d)\nu ^\beta (1+z)^{1+\beta }\left[\frac{d_L(z_{Arp})}{d_L(z)}\right]^2$$
(2)
where $`\nu `$ is the observed submillimeter frequency, $`T_d`$ is the dust temperature, and $`d_L`$ is the luminosity distance. We allow for an overall strength factor, $`f`$, relative to Arp 220, for which $`f=1`$. The above normalization constant is fixed by a $`\chi ^2`$ fit to the Arp 220 flux measurements. In the Rayleigh-Jeans long-wavelength limit, the submillimeter flux is approximately given by
$$S_s=f(0.825\mathrm{Jy})(1+z)^{1+\alpha _s}\left(\frac{T_d}{47\mathrm{K}}\right)\left(\frac{\nu }{375\mathrm{GHz}}\right)^{\alpha _s}\left[\frac{d_L(z_{Arp})}{d_L(z)}\right]^2$$
(3)
which is normalized to the measured 375 GHz (800$`\mu `$m) flux value for Arp 220 from Rigopoulou, Lawrence, & Rowan-Robinson (1996); here $`\alpha _s=\beta +2=3`$. The approximation in Eq. 3 is reasonably accurate for $`z3`$; at higher $`z`$ it extrapolates above the result of Eq. 2.
The radio flux is given by
$$S_r=f(0.148\mathrm{Jy})(1+z)^{1+\alpha _r}\left(\frac{\nu }{8.4\mathrm{GHz}}\right)^{\alpha _r}\left[\frac{d_L(z_{Arp})}{d_L(z)}\right]^2$$
(4)
which is normalized to the measured 8.4 GHz flux for Arp 220 from Condon et al. (1991). Since both $`S_s`$ and $`S_r`$ depend linearly on the star formation rate, the scale factor $`f`$ is the same. In local star forming galaxies there is an observed spectral flattening at 1.4 GHz due to free-free absorption (Condon 1992). Since 1.4 GHz measurements at higher redshifts sample higher frequencies that are not sensitive to these absorption effects, we adopt the standard value $`\alpha _r=0.8`$ from Condon (1992) to extrapolate local 8.4 GHz fluxes to 1.4 GHz fluxes, in order to be consistent with high redshift observations.
The ratio of Eqs. 3 and 4 with $`T_d=47`$K gives
$$\frac{S_{353\mathrm{GHz}}}{S_{1.4\mathrm{GHz}}}=1.1\times (1+z)^{3.8}$$
(5)
A relation of this type was obtained by CY99 based on the SED of M82. In their relation, the factor of 1.1 in Eq. 5 is replaced by 0.25 and the power-law index $`\alpha _s\alpha _r=3.8`$ is replaced by 4.3. However, M82, with more than an order of magnitude lower luminosity than Arp 220 ($`L_{FIR}=1.78\times 10^{10}h_{65}\mathrm{L}_{}`$; Hughes, Dunlop, & Rawlings 1997), is not a ULIG, and so we would argue that our result is more appropriate for the distant submillimeter sources. In addition to their analytic result, CY99 give empirical results for the flux ratio versus redshift for both Arp 220 and M82. Other authors have used the span of the four CY99 curves in their Fig. 1 to estimate redshift ranges; however, the use of M82 again introduces an inappropriate uncertainty in obtaining redshifts.
Inverting Eq. 5, the value of $`z`$ can be related to the flux ratio by
$$z+1=0.98\times \left(\frac{S_{353\mathrm{GHz}}}{S_{1.4\mathrm{GHz}}}\right)^{0.26}$$
(6)
which holds for $`z3`$. For higher redshifts the ratio of Eqs. 2 and 4 must be used. The flux ratio method for estimating redshifts has the advantage that the luminosity distance drops out, and thus there is no dependence on $`\mathrm{H}_\mathrm{o}`$ or on the cosmology.
We make the assumption that unevolved local ULIGs are representative of the distant ULIG population. We show in the following that the Arp 220 SED provides a good representation of the submillimeter-to-radio flux ratios of the ensemble of local ULIGs placed at appropriate redshifts. We use a compilation of submillimeter measurements at 3000, 857, 667, and 375 GHz from Rigopoulou et al. (1996) to infer 353 GHz fluxes at redshifts of 7.5, 1.44, 0.9, and 0.06. Fluxes at 1.4 GHz were extrapolated from the 8.4 GHz measured values of Condon et al. (1991) using a standard synchrotron emission spectrum with $`\alpha =0.8`$, as justified earlier. We plot the mean 353 GHz to 1.4 GHz flux ratios versus redshift as solid diamonds in Fig. 6a with ranges from the minimum to maximum flux ratios. The number of objects in the bins are, in order of decreasing redshift, 25, 2, 4, and 7. The data are to be compared with the redshifted Arp 220 flux ratio calculated from Eqs. 2 and 4 and plotted as a solid line in the figure. The spread in measured ratios is about a multiplicative factor of two relative to Arp 220, as indicated by the dashed curves. Thus, the redshifted Arp 220 flux ratio is likely to approximate the flux ratios of ULIGs at high redshifts to this accuracy.
The above conclusion is strengthened by Fig. 6b where we have plotted as solid (star forming galaxies) and open (AGN) squares available submillimeter observations with corresponding radio and redshift information. The star forming galaxies fall neatly on the curves. The average submillimeter-to-radio flux ratio of the galaxies in our survey which have spectroscopic redshifts and have been observed in the submillimeter are shown in two redshift bins as the large filled diamonds with $`1\sigma `$ uncertainties. The lowest point is consistent with a $`3\sigma `$ null detection, but in the higher redshift bin there is a strong positive detection consistent with the Arp 220 ratio, showing that these objects have cool dust emission and obey the FIR-radio relation at these redshifts. The short horizontal lines show the redshift ranges for our six significant radio-selected submillimeter detections; all are in the $`z=13`$ redshift range. The top portion of Table 3 gives the fluxes and the millimetric redshifts and ranges for these six sources, as well as their bolometric luminosities for both $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (see § 5).
The ULIG sources plotted in Fig. 6 have a range of emissivity indices, $`\beta `$, and dust temperatures in the vicinity of 50 K. It is therefore quite remarkable that all the sources follow a constrained envelope around the Arp 220 submillimeter-to-radio flux ratio. The implication of this empirical correlation is that for practical purposes a temperature-redshift degeneracy (Blain 1999) does not present a serious problem to millimetric redshift estimation.
A recent empirical analysis by Carilli & Yun (2000), which appeared subsequent to the submission of our paper, finds conclusions that are consistent with our analysis above.
## 5. Radio and Submillimeter Fluxes versus Redshift
In our analysis we assume a flat Universe, $`\mathrm{\Omega }_\mathrm{M}+\mathrm{\Omega }_\mathrm{\Lambda }=1`$, where $`\mathrm{\Omega }_\mathrm{M}`$ is the matter density and $`\mathrm{\Omega }_\mathrm{\Lambda }`$ the vacuum density. The absolute flux values in Eqs. 2 and 4 depend on the cosmology through the luminosity distance (Carroll, Press, & Turner 1992)
$$d_L(z)=c\mathrm{H}_{\mathrm{o}}^{}{}_{}{}^{1}(1+z)_0^z𝑑z^{}/[(1+z^{})^2(1+\mathrm{\Omega }_\mathrm{M}z^{})z^{}(2+z^{})\mathrm{\Omega }_\mathrm{\Lambda }]^{1/2}$$
(7)
The basic physics underlying the joint use of radio and submillimeter observations is contained in Figs. 7 and 8. Figure 7 shows predicted Arp 220 radio and submillimeter fluxes versus redshift for the $`\mathrm{\Omega }_\mathrm{M}=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology ($`\mathrm{H}_\mathrm{o}=65`$ km s<sup>-1</sup> Mpc<sup>-1</sup>) with selected overall strengths $`f=1/3`$, 1, 3, and 6 relative to Arp 220. $`S_{353\mathrm{GHz}}`$ is flat with redshift for $`z1`$ due to the negative $`K`$-correction, whereas $`S_{1.4\mathrm{GHz}}`$ falls sharply with increasing redshift due to the $`1/d_L^2`$ dependence. These behaviors are characteristic of other cosmologies as well. Thus, radio surveys detect a high proportion of low redshift sources whereas submillimeter surveys offer comparable sensitivity to both moderate ($`z1`$) and high redshifts.
In Fig. 8 we illustrate the dependence on cosmology of the absolute fluxes in the $`S_{1.4\mathrm{GHz}}`$ versus $`S_{353\mathrm{GHz}}`$ plane. Figure 8a shows predicted Arp 220 fluxes for the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology, and Fig. 8b shows the same for the $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ cosmology favored by recent distant supernovae type Ia observations (Perlmutter et al. 1999; Riess et al. 1998). Redshift labels are given on the $`f=1`$, 3, and 6 curves. In Fig. 8a the bright submillimeter sources fall near the $`f=3`$ contour, indicating that these sources have luminosities several times the luminosity of the Arp 220 prototype. For the $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ cosmology we would infer somewhat higher luminosities from the curves, although the redshift estimates remain the same.
From Fig. 8, we see that any submillimeter sources above our detection threshold of 6 mJy would not be detectable in the radio if $`z4`$. Of the significant submillimeter sources detected in our survey, two had no radio counterparts and one, HDF850.1, had only a supplemental 8.5 GHz radio detection. Based on the millimetric redshift estimator, these sources are potentially at very high redshift. The bottom three lines of Table 3 give the fluxes and the millimetric redshifts and ranges, along with the bolometric luminosities for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$, for the three significant SCUBA detections without corresponding 1.4 GHz detections; here the numerical entries are based on the 0.040 mJy ($`5\sigma `$) radio flux limit, which corresponds to the lowest possible redshift and $`L_{bol}`$.
Since we have covered only one-third of our $`80`$ square arcminute area with our targeted SCUBA observations, there could in principle be more high redshift submillimeter sources with no radio detections. However, our present radio counts already saturate the bright submillimeter counts distribution from the combined blank field submillimeter surveys (see Fig. 4) and hence argue against the probability of finding many such additional sources.
In contrast, sources of Arp 220 or sub-Arp 220 strength could be seen in the radio but not in the submillimeter with our present 6 mJy threshold if $`z2`$. Figure 7b shows that we have many such candidates. Our targeted SCUBA observations of optical/NIR-faint radio sources are therefore selecting the high redshift end of the faint radio source population.
## 6. Rest-frame Color versus Redshift
In Fig. 9 we plot rest-frame AB(2800)–AB(8140) color versus redshift for the radio sample with spectroscopic identifications. The colors of field galaxies in the HDF and SSA22 fields are indicated with tiny solid circles to illustrate how the radio galaxies (solid triangles for sources with $`\alpha _r>0.3`$ but otherwise solid circles) generally have very red AB colors. These colors range from $`1.5`$ to $`4.5`$. There is no evidence for a color trend with redshift. The rest-frame AB colors of the submillimeter galaxies with $`HK^{}`$ detections (solid pentagons), although mostly lower limits, are not inconsistent with the colors of the lower-redshift radio sources.
At $`z=2`$, AB(2800)–AB(8140) roughly corresponds to AB($`I`$)–AB($`HK^{}`$), and thus submillimeter sources with similar colors to the radio sources would have observed Vega-based $`IHK^{}`$ colors in the range $`3.16.1`$, which would mostly place them in the extremely red object (ERO) category ($`IHK^{}>3.7`$ or $`IK>4`$; e.g., McCracken et al. 1999 and references therein). This result is consistent with the recent detections of bright ($`K<20`$) EROs as submillimeter sources (Cimatti et al. 1998; Dey et al. 1999; Smail et al. 1999a). In general, however, the submillimeter sources are so optically faint that identifying them as EROs is difficult.
## 7. Properties of the Radio and Submillimeter Populations
The luminosities of our submillimeter sources were obtained by scaling the Arp 220 luminosity at redshift $`z`$ by the relative source strengths, $`f`$; these are given in Table 3. The FIR luminosity provides a direct measure of the current star formation rate (SFR); with present detection capabilities, the radio provides greater sensitivity to SFRs at low redshifts ($`z1.5`$) while the submillimeter is superior at high redshifts. In this sense radio and submillimeter data complement each other in getting the SFR over the full range of redshifts. This relative capability is illustrated in Fig. 10 where the strength factors, $`f`$, as determined from the radio luminosities relative to redshifted Arp 220 luminosities, are plotted versus redshift in the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology. Here the dotted curve represents the radio threshold of 40 $`\mu `$Jy, and the solid line at a value of two represents our 6 mJy submillimeter threshold of roughly two times the luminosity of Arp 220 for the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology (see Fig. 8). These two threshold curves cross at $`z2.3`$. Since the bulk of the radio detections are at $`z1.3`$ and have strength factors $`f<1`$, their non-detection in the submillimeter is as expected; our $`z2`$ sources are detected in both the radio and submillimeter. At redshifts $`z34`$ the submillimeter may provide the only access to distant ULIGs.
Figure 10 also provides a means to establish the completeness of the ULIG number distribution versus redshift. Our distant ULIG sample is complete over the redshift range $`z=13`$ (except possibly for the 4 out of 22 optically-faint radio sources that were not observed) because these $`z=13`$ sources are above both radio and submillimeter detection thresholds. At $`z3`$ there is the risk that we may lose sources that fall below the rising 1.4 GHz threshold relative to the Arp 220 luminosity at redshift $`z`$.
### 7.1. Space Density
The space density of our five significant $`S_{850\mu \mathrm{m}}>6`$ mJy submillimeter sources in the redshift range $`z=13`$ is
$`n`$ $`=`$ $`3.5_{1.5}^{+2.4}\times 10^5,\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (8a)
$`n`$ $`=`$ $`1.1_{0.5}^{+0.7}\times 10^5,\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (8b)
in units $`h_{65}^3\mathrm{Mpc}^3`$. In comparison, the space density of local ULIGs ($`z<0.15`$; Kim & Sanders 1998) with bolometric luminosities above $`10^{12}h_{65}^2\mathrm{L}_{}`$ is only $`1.7\times 10^7h_{65}^3\mathrm{Mpc}^3`$ for $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ and $`1.4\times 10^7h_{65}^3\mathrm{Mpc}^3`$ for $`\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$, about a factor of 200 and 80 lower than our high redshift results. The discrepancy in space density is still larger if we compare with local ULIGs that are comparably luminous to the $`>6`$ mJy population. Thus, enormous evolution in the ULIG population must take place between $`z2`$ and the present.
It is interesting to note a possible empirical similarity between the redshift distribution of quasars and that of submillimeter sources, which might be expected if quasars are a successor stage to merger processes forming submillimeter sources (Sanders & Mirabel 1996). The relative space density of a complete radio-selected quasar sample (Shaver et al. 1998), which should be unaffected by dust obscuration, is approximately described by a Gaussian as a function of redshift, given by
$$n(z)e^{(z2.28)^2/1.72}$$
(9)
The peaking of this distribution at $`z2.3`$ and the rapid fall-off at low redshift seems consistent with our measured submillimeter redshift distribution. Recent ROSAT X-ray studies have found a very high redshift tail ($`z4`$) to the quasar distribution that may be above the Gaussian expectation (Hasinger 1998). Since the ROSAT sources were not seen in the optical, it is plausible that they are dust obscured. Similarly, we have found a few submillimeter sources that plausibly lie at $`z3`$.
### 7.2. Conversion from Luminosity to Star Formation Rate
The observed ultraviolet light can most simply be translated to a metal production rate where, as first pointed out in Cowie (1988), there is an extremely tight relation since both metals and ultraviolet light are produced by the same massive stars. However, it is usual (particularly in the FIR) to translate the metal production rate to a total stellar mass production rate, which may be more intuitively interpreted. Unfortunately, this extrapolation requires knowledge of the shape of the initial mass function (IMF), which is quite uncertain.
In the following we shall assume a Salpeter IMF ($`\psi (M)M^{2.35}`$, $`0.1125\mathrm{M}_{}`$). The SFR is related to the rest-frame 2800 Å luminosity per unit frequency by
$$\dot{M}=\varphi \times L_{2800}$$
(10)
with $`L_{2800}`$ in ergs s<sup>-1</sup> Hz<sup>-1</sup> and $`\dot{M}`$ in M$`{}_{}{}^{}\mathrm{yr}_{}^{1}`$. The Cowie (1988) estimate of $`\varphi `$ from nucleosynthesis arguments and galaxy SED modelling is $`2.2\times 10^{28}`$ for the above IMF limits. Later Songaila, Cowie, & Lilly (1989) suggested that $`\varphi `$ should be about a factor 1.7 lower based on stellar synthesis. Madau et al. (1996) obtained a value of $`1.5\times 10^{28}`$ based on the evolutionary models of Bruzual & Charlot (1993). The relatively invariant estimates of $`\varphi `$ are a direct consequence of the nucleosynthesis arguments underlying the modelling; the primary uncertainty in $`\dot{M}`$ lies in the stellar IMF. We adopt the Madau et al. value for consistency with other recent papers that follow this normalization. In terms of the solar luminosity, the SFR relation is then
$$\dot{M}=\mathrm{\Phi }\times \nu L_{2800}/\mathrm{L}_{}$$
(11)
where $`\mathrm{\Phi }=5.3\times 10^{10}`$.
The conversion of the bolometric luminosity to a total stellar mass production rate is an even more invariant prediction of the nucleosynthesis arguments since knowledge of the exact shape of the ultraviolet SED of the galaxies is not needed. The bolometric luminosity in units of the solar luminosity is
$$L_{bol}=BC\times \nu L_{2800}/\mathrm{L}_{}$$
(12)
where $`BC`$ is the bolometric correction. For a flat $`S_\nu `$ over the wavelength range 912 Å to 22000 Å, $`BC=2.95`$. Rowan-Robinson et al. (1997) computed a bolometric correction of 3.5 from the Bruzual & Charlot (1993) models, which we adopt. Then the SFR is related to $`L_{bol}`$ by
$$\dot{M}=(\mathrm{\Phi }/BC)\times (L_{bol}/\mathrm{L}_{})=1.5\times 10^{10}\times L_{bol}/L_{}$$
(13)
$`\dot{M}`$ would be increased by a factor of 3.3 if we had instead assumed a Miller-Scalo IMF over the same mass range.
Equations 11 and 13 therefore provide a self-consistent description of the mass production rates seen in the optical and submillimeter. Light which directly escapes the galaxy can be mapped with the ultraviolet light, while reprocessed light, which escapes the galaxy in the submillimeter, can be calibrated with Eq. 13.
We now use the Arp 220 bolometric luminosity $`L_{bol}=1.36\times 10^{12}h_{65}^2\mathrm{L}_{}`$ (Klaas et al. 1997) and the 8.4 GHz radio flux of 0.148 Jy (Condon et al. 1991) to calibrate the SFRs for our radio sources. We find $`\dot{M}=200h_{65}^2\mathrm{M}_{}\mathrm{yr}^1`$; thus
$$\dot{M}=L_{8.4\mathrm{GHz}}/6.4\times 10^{27}$$
(14)
where the normalization factor is in units of ergs s<sup>-1</sup> Hz<sup>-1</sup>. From the equations of Condon (1992) translated to the same assumptions about the IMF, we would obtain a normalization of approximately $`1.8\times 10^{27}`$ or about a factor of 3.5 smaller than our result. Our normalization translated with the $`\nu ^{0.8}`$ spectral shape to 1.4 GHz is
$$\dot{M}=L_{1.4\mathrm{GHz}}/2.7\times 10^{28}$$
(15)
### 7.3. Contribution of the Radio and Submillimeter Sources to the Star Formation Rate Density
Using Eq. 15, we can convert our radio luminosities into SFRDs. We know that 95 per cent of the micro-Jansky radio sources in the HDF region have been resolved at $`0.2^{\prime \prime }`$ resolution with the Multi-Element Radio Linked Interferometer (MERLIN) at 1.4 GHz. The median angular size for the radio emission is $`1^{\prime \prime }2^{\prime \prime }`$ (Richards 1999a), which suggests radio emission on galactic or sub-galactic size scales. Thus, the radio emission in most of these systems may originate primarily from star formation, although we cannot exclude contributions to the radio flux densities from embedded AGN.
We divide our radio sources with spectroscopic redshifts into two bins, $`0.1<z0.7`$ and $`0.7<z1.3`$, excluding sources with radio spectral indices $`\alpha _r>0.3`$ that might be AGN. Of the 38 presumed star forming objects with $`HK^{}<20`$ in our sample, 23 have spectroscopic redshifts. Most of the remaining 15 sources have not been observed and thus may be expected to follow the same redshift distribution as the identified sources. We therefore calculate the differential comoving volume over a survey area $`(23/38)\times 79.4`$ arcmin<sup>2</sup>. The volume is integrated from $`z_{min}`$, the lower redshift limit of the bin, to $`z_{max}`$, either the radio luminosity limit of the survey or the upper redshift limit of the bin, whichever is smaller. The comoving volume element in Mpc<sup>3</sup> is
$$\mathrm{\Delta }V=\frac{1}{3}\times A\times 8.46\times 10^8\times \left[\left(\frac{d_L(z_{max})}{1+z_{max}}\right)^3\left(\frac{d_L(z_{min})}{1+z_{min}}\right)^3\right]$$
(16)
where $`A`$ is the area in arcmin<sup>2</sup> and $`d_L`$ is in Mpc.
We find that for the $`0.1<z0.7`$ bin, the SFRD in units $`h_{65}\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3`$ is
$`\mathrm{SFRD}`$ $`=`$ $`0.033_{0.009}^{+0.012},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (17a)
$`\mathrm{SFRD}`$ $`=`$ $`0.025_{0.007}^{+0.009},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (17b)
and for the $`0.7<z1.3`$ bin,
$`\mathrm{SFRD}`$ $`=`$ $`0.048_{0.015}^{+0.020},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (18a)
$`\mathrm{SFRD}`$ $`=`$ $`0.032_{0.010}^{+0.014},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (18b)
where the uncertainties are Poissonian based on the number of sources. These represent lower limits to the SFRD since we have not attempted to correct for the contributions of faint sources below our flux limits.
The SFRD from our submillimeter sources can be estimated under the assumption that star formation dominates AGN contributions. We can calculate the contribution of the $`S_{850\mu \mathrm{m}}>6`$ mJy submillimeter sources at $`z=13`$ to the SFRD using Eq. 13; we find
$`\mathrm{SFRD}`$ $`=`$ $`0.023_{0.010}^{+0.016},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (19a)
$`\mathrm{SFRD}`$ $`=`$ $`0.014_{0.006}^{+0.009},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (19b)
Unlike the number distribution versus redshift, the SFRD needs large corrections for completeness due to the fact that we are detecting only relatively bright sources here (at about six times the luminosity at which the submillimeter background is primarily resolved), and the distribution $`dN/dS`$ increases rapidly as $`S`$ decreases. To estimate the completeness correction, we assume that the dependences of $`N`$ on $`S`$ and $`z`$ factorize, $`d^2N/dSdz=g(S)h(z)`$. This is a plausible assumption in the submillimeter where the fluxes are nearly independent of redshift; nonetheless, in view of the Smail et al. (1999b) study, this assumption remains to be confirmed.
We can determine the completeness correction for the SFRD using the empirical number distribution at 850 $`\mu `$m versus $`S`$
$$dN/dS=3.0\times 10^4\mathrm{deg}^2\mathrm{mJy}^1/S^{3.2}$$
(20)
that describes the measured submillimeter counts above 2 mJy (Barger, Cowie, & Sanders 1999). We are making the assumption that the flux to $`L_{FIR}`$ conversion based on Arp 220 applies in the low submillimeter flux region; the justification is that even the dominant $`1`$ mJy population are near-ULIG sources. The completeness correction over all submillimeter fluxes is therefore the measured $`850\mu `$m extragalactic background light (EBL) divided by the $`850\mu `$m light above 6 mJy. The $`850\mu `$m EBL measurement of $`3.1\times 10^4`$ mJy deg<sup>-2</sup> from Puget et al. (1996) and $`4.4\times 10^4`$ mJy deg<sup>-2</sup> from Fixsen et al. (1998) imply correction factors in the submillimeter of 11 and 15, respectively. Thus, the estimated total submillimeter contribution to the SFRD in units of $`h_{65}\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3`$ is
$`\mathrm{SFRD}`$ $`=`$ $`0.25_{0.11}^{+0.17},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (21a)
$`\mathrm{SFRD}`$ $`=`$ $`0.15_{0.06}^{+0.10},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (21b)
where we have used the factor of 11 completeness correction.
More speculatively, we can determine the SFRD from higher redshift sources using our two $`>6`$ mJy submillimeter sources without radio counterparts. In this case we use the volume from $`z=36`$ and the actual area surveyed in the submillimeter. We find
$`\mathrm{SFRD}`$ $`=`$ $`0.029_{0.019}^{+0.038},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (22a)
$`\mathrm{SFRD}`$ $`=`$ $`0.016_{0.011}^{+0.022},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (22b)
After including the factor of 11 completeness correction, this becomes
$`\mathrm{SFRD}`$ $`=`$ $`0.32_{0.20}^{+0.42},\mathrm{\Omega }_\mathrm{\Lambda }=0`$ (23a)
$`\mathrm{SFRD}`$ $`=`$ $`0.18_{0.12}^{+0.24},\mathrm{\Omega }_\mathrm{\Lambda }=2/3`$ (23b)
H98 used photometric redshift estimates to infer that four of their five $`S_{850\mu \mathrm{m}}>2`$ mJy sources were in the redshift range $`z=24`$. While these source identifications were problematic, it appears likely from the present work that the redshifts do lie in this rough redshift range. Using our parameters and a scaled Arp 220 SED, we find that the SFRD for their sources with $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ is $`0.10_{0.05}^{+0.08}h_{65}\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3`$. If we make a completeness correction to include the contribution below 2 mJy, we obtain SFRD=$`0.28_{0.14}^{+0.22}h_{65}\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3`$, in good agreement with our result in Eq. 21a.
The presence of a substantial fraction of AGN-dominated ULIG sources would reduce the above SFRDs. In a recent near-infrared spectroscopic study of 64 local ULIGs, Veilleux, Sanders, & Kim (1999) found AGN characteristics in $`2025`$ per cent of the sample, which increased to $`3550`$ per cent for the sample with $`L_{IR}>10^{12.3}\mathrm{L}_{}`$. Thus, our $`>6`$ mJy contributions to the SFRD may need to be reduced by a factor $`1.52`$. However, the lower AGN fraction in fainter ULIGs seen locally suggests that AGN contamination may be less of an issue for the extrapolated SFRD of the whole submillimeter population.
### 7.4. Comparison with the Optical Star Formation Rate Density Diagram
The determination of the SFRD from optical observations has been a subject of intense investigation. Observations first indicated a rather rapid rise in the SFRD from $`z=01`$ followed by a sharp decline at higher redshifts with the peak SFRD being $`z1.5`$ (Madau et al. 1996). A recent modification in the inferred optical SFRD at low redshifts was made by Cowie, Songaila, & Barger (1999), whose data indicated a more gradual rise in the SFRD than had previously been found by Lilly et al. (1996).
It was realized that dust obscuration effects could result in factors of 3 to 5 (Pettini et al. 1997; Meurer, Heckman, & Calzetti 1999) increases in the SFRD at high redshift. With these rather uncertain dust corrections taken into account, it has been argued that the SFRD flattens at a constant $`\mathrm{SFRD}0.2\mathrm{h}_{65}\mathrm{M}_{}\mathrm{yr}^1\mathrm{Mpc}^3`$ in the $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$ cosmology for $`z=15`$ (Steidel et al. 1999).
In Fig. 11 we compare the star formation history in the optical (without extinction corrections) with that which we obtain in the submillimeter both before (Eqs. 19a and 22a; solid triangles) and after (Eqs. 21a and 23a; solid circles) correcting for incompleteness. We also include our newly determined radio SFRD limits (Eqs. 17a and 18a) on the figure as solid squares.
The submillimeter contribution to the SFRD inferred from our $`>6`$ mJy observations is comparable to the ultraviolet/optical contribution to the SFRD. The two wavelength regimes are likely sampling different stages in galaxy formation. The submillimeter detects the formation of massive spheroids while the ultraviolet/optical detects the formation of smaller disk or bulge systems. The approximate equality of the optical and submillimeter backgrounds supports this hypothesis; the metal density in present-day disks is roughly comparable to that in the spheroidal components of galaxies, so comparable amounts of light are expected to be produced in their formation (Cowie 1988).
The completeness-corrected submillimeter SFRD, shown by the solid circles in Fig. 11, is based on the assumption that fainter submillimeter sources have the same redshift distribution and properties as the $`>6`$ mJy sample. Hence the SFRD from the entire population contributing to the submillimeter background is about an order of magnitude higher than the extinction-uncorrected ultraviolet/optical SFRD. Since the submillimeter measures reradiated optical light, the observed ultraviolet/optical SFRD contribution should be added to the submillimeter contribution.
We also plot in Fig. 11 the SFRD determined from the sum of the SFRD of the local ULIG data ($`z<0.15`$; $`L_{bol}>10^{12}h_{65}^2L_{}`$) from Kim & Sanders (1998) and the local near-ULIG data ($`z<0.02`$; $`2\times 10^{11}h_{65}^2L_{}<L_{bol}<10^{12}h_{65}^2L_{}`$) from Sanders et al. (in preparation). The latter $`L_{bol}`$ selection is imposed in order to be able to compare with our completeness-corrected submillimeter data points, which should include all sources with luminosities $`L_{bol}>2\times 10^{11}h_{65}^2L_{}`$. We note that the radio data in Fig. 11 cannot be straightforwardly compared to the ULIG data because the radio sources may not satisfy the same luminosity criteria.
The completeness-corrected SFRD in Fig. 11 shows a very rapid evolution, $`(1+z)^6`$, in the SFRD of ULIGs from $`z0`$ to $`z13`$. Fast evolution is not surprising if the distant submillimeter sources are associated with major merger events giving rise to the formation of massive spheroidal systems (Smail et al. 1998; Eales et al. 1999; Lilly et al. 1999; Trentham et al. 1999). Barger, Cowie, & Sanders (1999) showed that the volume density of submillimeter sources at high redshift ($`n=5\times 10^3h_{65}^3`$ Mpc<sup>-3</sup> for $`q_o=0.5`$) is comparable to the volume density of present-day elliptical galaxies ($`n=10^3h_{65}^3`$ Mpc<sup>-3</sup>).
## 8. Conclusions
We have carried out an observational program designed to establish the overlap between the optical/NIR-faint radio population and the submillimeter population and to exploit the complementary information that the radio and submillimeter wavelengths provide to gain insights into the evolution of the galaxy populations. Our major conclusions are as follows
$``$ We have found that submillimeter sources at bright flux levels ($`>6`$ mJy) are associated with optical/NIR-faint radio sources ($`40300\mu `$Jy). This association provides a powerful means to conduct submillimeter surveys by preselecting potential bright submillimeter sources through high-resolution deep radio maps.
$``$ We have shown that the redshifted submillimeter-to-radio flux ratio of an unevolved Arp 220 SED reproduces the ensemble of local ULIG data placed at appropriate redshifts, and thus Arp 220 provides a prototype for ULIG sources at high redshift. From the Arp 220 model we have derived a millimetric redshift estimator to determine the redshifts of the submillimeter sources using the ratio of the submillimeter to radio fluxes. Our estimator is consistent with another recent study by Carilli & Yun (2000). The flux strengths relative to the redshifted Arp 220 SED indicate that the $`z=13`$ bright ($`>6`$ mJy) submillimeter sources are $`3\times `$ the luminosity of Arp 220 for $`q_o=0.5`$ and $`6\times `$ the luminosity of Arp 220 for $`\mathrm{\Lambda }`$-dominated models.
$``$ Through radio versus submillimeter flux plots of the Arp 220 predictions with redshift, we have shown that our present survey is sensitive only to sources with strengths comparable to or greater than Arp 220 and that our $`40\mu `$Jy radio threshold precludes the detection of very high redshift sources ($`z>4`$) if they are radio sources; however, we serendipitously observed two bright submillimeter sources without radio detections that may lie at extreme redshifts. An alternate possibility is that they could be colder galaxies than Arp 220 located at $`z<4`$. It remains to be established whether the very faint radio population down to $`1\mu `$Jy have submillimeter counterparts.
$``$ At $`z1`$ the ultraviolet/optical contribution to the SFRD dominates. At $`z>1`$ the submillimeter contribution from $`>6`$ mJy sources is comparable to the observed ultraviolet/optical contribution. The ultraviolet/optical and bright submillimeter are likely sampling two different stages in galaxy formation. The bright submillimeter detects the formation of massive spheroids while the ultraviolet/optical detects the formation of smaller disk or bulge systems.
Joint radio/submillimeter surveys are a powerful way to explore dust-obscured galaxies in the distant Universe since the radio and submillimeter approaches are complementary and allow us to infer redshifts and SFRDs for galaxy sources that we now conclusively see are inaccessible in shorter wavelength observations.
We thank Dave Sanders for valuable discussions, Nicolas Biver for expert observing assistance, and an anonymous referee for helpful comments about the manuscript. Support for this work was provided by NASA through Hubble Fellowship grants HF-01117.01-99A and HF-01123.01-99A awarded by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5-26555. |
warning/0001/hep-lat0001013.html | ar5iv | text | # The Overlap-Dirac Operator: Topology and Chiral Symmetry Breaking JLAB-THY-00-01FSU-SCRI-99-74
## I Overlap and domain wall Dirac operators
In these proceedings, we review some basic properties of the overlap Dirac operator and how its index can be computed by spectral flow techniques. One of the side results is that for fermions in the adjoint representation of $`SU(N)`$ we find evidence for fractional topological charge. The presentation is pedagogical with the intent of illustrating the origin of numerical difficulties in simulating overlap and domain wall fermions. Recent results from our work using overlap fermions can be found in references .
The massive overlap Dirac operator derived from the overlap formalism is
$$D_{\mathrm{ov}}(\mu )=\frac{1}{2}\left[1+\mu +(1\mu )\gamma _5ϵ(\mathrm{H}_L(m))\right]$$
(1)
where $`\mathrm{H}_L(m)`$ is a lattice hermitian Dirac-like operator describing a single fermion species with a large negative mass. The mass $`m`$ is a regulator parameter for the theory. In this work, we use the hermitian Wilson-Dirac operator $`\mathrm{H}_w(m)=\gamma _5D_{\mathrm{Wilson}}(m)`$, although we have tested other fermion actions. The mass parameter $`1<\mu <1`$ is related to the fermion mass by
$$m_f=Z_m^1\mu (1+𝒪(a^2)).$$
(2)
The propagator for external fermions is given by
$$\stackrel{~}{D}^1(\mu )=(1\mu )^1\left[D_{\mathrm{ov}}^1(\mu )1\right],$$
(3)
i.e. it has a contact term subtracted, which makes the massless propagator chiral: $`\{\stackrel{~}{D}^1(0),\gamma _5\}=0`$.
A massless vector gauge theory can also be obtained from domain wall fermions , where an extra, fifth dimension, of infinite extent is introduced. In the version of ref. , one can show that the physical (light) fermions contribute $`\mathrm{log}detD_{\mathrm{DW}}`$ to the effective action with the 4-d action
$$D_{\mathrm{DW}}=\frac{1}{2}\left[1+\mu +(1\mu )\gamma _5\mathrm{tanh}\left(\frac{L_s}{2}\mathrm{log}T\right)\right]$$
(4)
where $`T`$ is the transfer matrix in the extra dimension and $`L_s`$ its size. As long as $`\mathrm{log}T0`$ we obtain in the limit as $`L_s\mathrm{}`$
$$D_{\mathrm{DW}}\frac{1}{2}\left[1+\mu +(1\mu )\gamma _5ϵ(\mathrm{log}T)\right].$$
(5)
This is just the massive overlap Dirac operator up to the replacement $`\mathrm{H}_w\mathrm{log}T`$. It is easy to see that in the limit $`a_s0`$, where $`a_s`$ is the lattice spacing in the extra dimension (set to 1 above), one obtains $`\mathrm{log}T=\mathrm{H}_w\left(1+𝒪(a_s)\right)`$.
## II Some properties of the overlap Dirac operator
In many cases it is more convenient to use the hermitian version of the overlap Dirac operator (1):
$$H_o(\mu )=\gamma _5D_{\mathrm{ov}}(\mu )=\frac{1}{2}\left[(1+\mu )\gamma _5+(1\mu )ϵ(H_w)\right].$$
(6)
The massless version satisfies,
$$\{H_o(0),\gamma _5\}=2H_o^2(0).$$
(7)
It follows that $`[H_o^2(0),\gamma _5]=0`$, i.e. the eigenvectors of $`H_o^2(0)`$ can be chosen as chiral. Since
$$H_o^2(\mu )=(1\mu ^2)H_o^2(0)+\mu ^2$$
(8)
this holds also for the massive case.
The only eigenvalues of $`H_o(0)`$ with chiral eigenvectors are 0 and $`\pm 1`$. Each eigenvalue $`0<\lambda ^2<1`$ of $`H_o^2(0)`$ is then doubly degenerate with opposite chirality eigenvectors. In this basis $`H_o(\mu )`$ and $`D_{\mathrm{ov}}(\mu )`$ are block diagonal with $`2\times 2`$ blocks, e.g
$$D_{\mathrm{ov}}(\mu ):\left(\begin{array}{cc}(1\mu )\lambda ^2+\mu & (1\mu )\lambda \sqrt{1\lambda ^2}\\ (1\mu )\lambda \sqrt{1\lambda ^2}& (1\mu )\lambda ^2+\mu \end{array}\right),\gamma _5=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
(9)
For a gauge field with topological charge $`Q0`$, there are, in addition, $`|Q|`$ exact zero modes with chirality $`\mathrm{sign}(Q)`$, paired with eigenvectors of opposite chirality and eigenvalue 1. These are also eigenvectors of $`H_o(\mu )`$ and $`D_{\mathrm{ov}}(\mu )`$:
$$D_{\mathrm{ov}}(\mu )_{\mathrm{zero}\mathrm{sector}}:\left(\begin{array}{cc}\mu & 0\\ 0& 1\end{array}\right)\mathrm{or}\left(\begin{array}{cc}1& 0\\ 0& \mu \end{array}\right)$$
(10)
depending on the sign of $`Q`$.
We remark that from eigenvalues/vectors of $`H_o^2(0)`$ those of both $`H_o(\mu )`$ and $`D_{\mathrm{ov}}(\mu )`$ are easily obtained. There is no need for a non-hermitian eigenvalue/vector solver! For example, the Ritz algorithm will do just fine.
## III Implementations of the overlap Dirac operator
In practice, we only need the application of $`D(\mu )`$ on a vector, $`D(\mu )\psi `$, and therefore only the sign function applied to a vector, $`ϵ(H_w)\psi `$. Since we need the sign function of an operator (a large sparse matrix) this is still a formidable task.
Methods proposed for this computation are:
* A Chebyshev approximation of $`ϵ(x)=\frac{x}{\sqrt{x^2}}`$ over some interval $`[\delta ,1]`$ . For small $`\delta `$ a large number of terms are needed.
* A fractional inverse method using Gegenbauer polynomials for $`\frac{1}{\sqrt{x^2}}`$ . This has a poor convergence since these polynomials are not optimal in the Krylov space.
* Use a Lanczos based method to compute $`\frac{1}{\sqrt{x^2}}`$ based on the sequence generated for the computation of $`\frac{1}{x}`$ .
* Use a rational polynomial approximation for $`ϵ(x)`$ which can then be rewritten as a sum over poles:
$$ϵ(x)x\frac{P(x^2)}{Q(x^2)}=x\left(c_0+\underset{k}{}\frac{c_k}{x^2+b_k}\right)$$
(11)
The application of $`\chi ϵ(H_w)\psi `$ can be done by the simultaneous solution of the shifted linear systems
$$(H_w^2+b_k)\varphi _k=\psi ,\chi =H_w(c_0\psi +\underset{k}{}c_k\varphi _k).$$
(12)
One such approximation, based on the polar decomposition , was introduced in this context by Neuberger . We use optimal rational polynomials . The accuracy of this approximation is shown in Fig. II.
We note that in all methods listed above, one can enforce the accuracy of the approximation of $`ϵ(x)`$ for small $`x`$ by projecting out the lowest few eigenvectors of $`H_w`$ and adding their correct contribution exactly.
$$ϵ(H_w)=\underset{i=1}{\overset{n}{}}|\psi _iϵ(\lambda _i)\psi _i|+𝒫_{}^{(n)}\mathrm{App}[ϵ(H_w)]𝒫_{}^{(n)},𝒫_{}^{(n)}=\mathrm{𝟏}\underset{i=1}{\overset{n}{}}|\psi _i\psi _i|.$$
(13)
To invert $`D^{}D`$ for overlap fermions, we have, generically, an outer CG method (a 4-d Krylov space search) and an independent inner search method for $`ϵ(H_w)\psi `$ – maybe CG again. For domain wall fermions, on the other hand, a 5-d Krylov space search method is used. It may pay off to try to combine inner and outer CGs for overlap fermions by reformulating them into a 5-d problem .
## IV Index defined via the Overlap formalism
The massless overlap Dirac operator is
$$D_{\mathrm{ov}}=\frac{1}{2}\left[1+\gamma _5ϵ(\mathrm{H}_L)\right];\mathrm{H}_L=\gamma _5D_w(m).$$
(14)
The index is given by $`Q=\mathrm{tr}ϵ(\mathrm{H}_L)/2`$. We see $`Q`$ simply counts the deficit of the number of positive energy states of $`\mathrm{H}_L`$. A simple method of computing $`Q`$ at some fixed $`m`$ is via the spectral flow method . Consider the eigenvalue problem
$$\mathrm{H}_L(m)\varphi _k(m)=\lambda _k(m)\varphi _k(m),\frac{d\lambda _k(m)}{dm}=\varphi _k^{}(m)\gamma _5\varphi _k(m).$$
(15)
An efficient way to compute $`Q`$ is to compute the lowest eigenvalues of $`\mathrm{H}_L(m)`$ for $`m>0`$. We can prove the number of positive and negative eigenvalues of $`\mathrm{H}_L(m)`$ for $`m<0`$ must be the same, so we slowly vary the mass $`m`$ from $`m=0`$ while keeping track of the levels crossing zero and direction of crossings up to some $`m`$. In this way, we get the topological charge as a function of $`m`$.
We note the mass $`m`$ must be greater than the usual critical mass of $`\mathrm{H}_L(m)`$, otherwise no topology change occurs and the overlap operator does not describe a massless chiral fermion. This critical mass value shifts from its free field value $`0`$ to some positive value for non-zero gauge coupling. We should also choose a mass below $`2`$ so that in the continuum limit there are no doubler contributions.
In Fig. II, we show spectral flow results for a smooth background field of a single instanton . There is a reflection symmetry about $`m=4`$, namely the spectrum for $`8m`$ is opposite that of $`m`$. We see a mode crosses down in eigenvalue, then crosses up again near $`2`$. There is a degeneracy of $`3`$ for the modes just beyond $`2`$. Hence, as we increase $`m`$ we find a sequence of the index $`Q`$ of $`1`$, $`3`$, $`3`$, $`1`$ and $`0`$ for $`m=1`$, $`3`$, $`5`$ and $`8`$. The one-dimensional profile for the modes associated with the crossings are
$$z(t)=\underset{\stackrel{}{n}}{}\varphi _k^{}(\stackrel{}{n},t)\varphi _k(\stackrel{}{n},t)$$
(16)
which we compare with the continuum solution in the center panel of Fig. II
$$z(t;c,\rho )=\frac{1}{\left[2\rho \left(1+(\frac{(tc_4)}{\rho })^2\right)^{3/2}\right]}.$$
(17)
Also shown is the zero crossing point $`m`$ as a function of $`\rho `$. We see that as the size of the instanton decreases the crossing point (the mass) increases.
As we turn on gauge fields, the picture of the flows complicates and we can find crossings throughout the mass region beyond the critical mass . Since we are interested in the zero modes at the crossings, we compute the density of zero eigenvalues $`\rho (0;m)`$ of $`\mathrm{H}_w(m)`$ by fitting linearly to the integrated density
$$_0^\lambda \rho (\lambda ^{})𝑑\lambda ^{}=\rho (0)\lambda +\frac{1}{2}\rho _1\lambda ^2+\mathrm{}$$
(18)
In Fig. II we show $`\rho (0;m)`$ for quenched SU(3) lattices. We see that for $`m`$ beyond the critical mass region, the density $`\rho (0;m)`$ rises sharply to a peak, then drops but is never zero, hence the spectral gap is always closed. A similar result is also found for two flavor dynamical fermion backgrounds (simulated with positive physical quark mass, e.g. not simulated in the super-critical mass region). From a size distribution, we observe the zero modes are on the order of one to two lattices spacings for $`m`$ beyond the main band of crossings, namely for $`m`$ in the “flat” region of $`\rho (0;m)`$. We find that a physical quantity, like the topological susceptibility, appears constant within errors in this “flat” region, indicating that these small modes make no physical contribution.
For topology to change in a gauge field evolution, we must create dislocations. These produce the small modes observed above which force the spectral gap of $`\mathrm{H}_w(m)`$ to be closed. In the right panel of Fig. II an empirical fit of the density to an exponential of the inverse lattice spacing is shown. This result implies the density is only zero in the continuum limit. The density of zero eigenvalues of $`H_w(m)`$, $`\rho (0;m)`$, is non-zero in the quenched case, but decreases rapidly with decreasing coupling . Very roughly, we find $`\rho (0;m)/\sigma ^{3/2}e^{e^\beta }`$. We note the gauge action can be modified to reduce $`\rho (0;m)`$, or even eliminate it altogether at some fixed $`m`$ .
In Fig. II, we show the size distribution for quenched $`16^4`$ SU(2) ensembles. The profile from Eq. (16) is used to define a size motivated by the t’Hooft zero mode in Eq. (17)
$$\rho _z=\frac{1}{2}\frac{\underset{t}{}z(t)}{z_{\mathrm{max}}}$$
(19)
We see there is a large number of modes about 1 to 2 lattice units in size. There is a corresponding secondary peak around 5 lattice units and all the distributions are bounded in size. This result indicates that the size distribution of zero modes does not show evidence for a peak at a physical scale (as suggested by some models) even after we remove the small modes which are most likely lattice artifacts . Instead, the observed scaling in lattice units is suggestive of a finite volume effect.
## V Evidence for fractional topological charge
In a continuum background field with topological charge $`Q`$, the index of the massless Dirac operator in the adjoint representation is equal to $`2NQ`$ . Classical instantons carry integer topological charge and can thus only cause condensation of an operator with $`2N`$ Majorana fermions. Witten argued that a bilinear gluino condensate exists in SUSY YM. Self-dual twisted gauge field configurations, with fractional topological charge $`1/N`$ exist (t’Hooft), and could explain a bilinear gluino condensate. What about non-classical gauge field configurations?
The adjoint representation is real $``$ spectrum of $`\mathrm{H}_L`$ is doubly degenerate: adjoint index can only be even. Are all even values realized, or only multiples of $`2N`$ ? To this end, we studied the flow in the adjoint representation on two SU(2) ensembles at the same fixed physical volume . We do find configurations with $`I_a=4I_f`$ (number of adjoint and fundamental crossings), but we also find configurations with $`I_a4I_f`$. An example is shown in Fig. II.
To check if this evidence for fractional topological charge is a lattice artifact, we plot $`\mathrm{\Delta }=I_a4I_f`$ in Fig. II. Note that $`\mathrm{\Delta }`$ takes on only even values. The probability of finding a certain value of $`\mathrm{\Delta }`$, $`p(\mathrm{\Delta })`$, is plotted for two ensembles in Fig. II. We find that $`p(\mathrm{\Delta })`$ for $`|\mathrm{\Delta }|>2`$ decreases as one goes toward the continuum limit at a fixed physical volume. However, $`p(\pm 2)`$ does not decrease indicating that it might remain finite in the continuum limit.
## VI Main problem for Overlap and Domain Wall fermions
The existence of small eigenvalues illustrated in Fig. II hampers the approximation accuracy and convergence properties of implementations of $`ϵ(H_w)`$. Eigenvector projection both increases the accuracy of the approximation and decreases the condition number, e.g. of the inner CG.
The existence of small eigenvalues has implications also for domain wall fermions. One can show that the spectrum of $`\mathrm{log}T(m)`$ of Eq. (4) around zero is the same as the spectrum of $`H_w(m)`$ . While the small eigenvalues of $`\mathrm{log}T(m)`$ don’t appear to cause algorithmic problems for domain wall fermions, they can induce rather strong $`L_s`$ dependence of physical quantities, and hence causing the need for large $`L_s`$.
### VI-1 Domain Wall and Overlap-Dirac operator spectral flow for smooth SU(2)
As an illustration of the effects of low lying modes of $`H_w(m)`$, we show in Fig. II the spectral flow of the hermitian domain wall operator $`H_{\mathrm{DW}}(m)=\mathrm{\Gamma }_5D_{\mathrm{DW}}(m)`$ on a smooth $`4^2\times 8^2`$ single instanton background with Dirichlet boundary conditions (BC). The $`\mathrm{\Gamma }_5`$ includes a parity operator. Also shown is the corresponding hermitian Wilson $`H_w(m)`$ spectral flow. The “zero” DWF eigenvalue sets in slowly in $`L_s`$. We need $`L_s1/\lambda _{\mathrm{min}}`$ where $`\lambda _{\mathrm{min}}`$ is the lowest eigenvalue of $`H_w(m)`$ (which is similar to the lowest eigenvalue of $`\mathrm{log}(T(m))`$) for $`ϵ(\mathrm{log}T(m))\mathrm{tanh}\left(\frac{L_s}{2}\mathrm{log}T(m)\right)`$. In fact, the (almost) chiral zero mode eigenvalue $`\lambda _{\mathrm{DW}}(m)\mathrm{const}\times (1\mathrm{tanh}(\lambda _{\mathrm{min}}(m)L_s/2))`$ for $`m`$ beyond the crossing indicating a sensitivity to the hermitian Wilson operator eigenvalue.
In Fig. II, we show a similar plot of the hermitian overlap Dirac operator $`H_o(m)`$ on a smooth $`8^4`$ single instanton background with Dirichlet BC. There are strict zero modes after the crossing, $`m=0.6`$, $`0.7`$, and $`0.8`$. Also shown are the zero mode profiles for these masses which are quite similar and nicely follow the t’Hooft zero mode solution.
## VII Conclusions
The overlap Dirac operator has the same chiral symmetries as continuum fermions. It has exact zero modes in topologically non-trivial gauge fields. It is therefore well suited for a study of the interplay of topology, with its associated exact zero modes, and chiral symmetry breaking, determined by the density of small eigenvalues.
The creation of dislocations necessary for change of topology causes numerical difficulties with overlap and domain wall simulations. These dislocations are purely a property of the gauge actions used and are not a fundamental limitation of the chiral fermion formalisms. In practice, the projection technique used in the overlap simulations is vital to precisely control the adverse influence of the low lying zero modes of the hermitian Wilson operator. The same technique can be used for domain wall simulations, but it is more cumbersome. Further work is directed towards lowering (and possibly eliminating) the density of low lying zero modes. A dynamical HMC algorithm for the overlap Dirac operator has also been developed .
RGE would like to thank the organizers for a splendid workshop. The work of RGE and UMH has been supported in part by DOE contracts DE-FG05-85ER250000 and DE-FG05-96ER40979. RGE was also supported by DOE contract DE-AC05-84ER40150 under which SURA operates the Thomas Jefferson National Accelerator Facility. |
warning/0001/gr-qc0001045.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Ever since the discovery that thermodynamic properties of black holes in anti–de Sitter (AdS) space-time are dual to those of a field theory in one dimension fewer, there has been of much interest in Reissner–Nordström (RN)–AdS black hole , which now becomes a prototype example to study this AdS/CFT correspondence . On the other hand, after Unruh’s work , it has been known that a thermal Hawking effect on a curved manifold can be looked at as an Unruh effect in a higher dimensional flat space-time. Recently, non-trivial works of isometric embeddings of the RN black hole and M2-, D3-, M5-branes into flat spaces with two times has been studied to get some insight of the global aspect of the space-time geometries in the context of brane physics. Moreover, several authors have also shown that global embedding Minkowski space (GEMS) approach of which a hyperboloid in a higher dimensional space corresponds to original curved space could provide a unified derivation of temperature for a wide variety of curved spaces. These include the static, rotating, charged BTZ , the Schwarzschild together with its AdS extensions, and the RN black holes. Therefore, it is interesting to study the geometry of the RN–AdS and their thermodynamics in this GEMS approach.
In this paper we will analyze the Hawking and Unruh effects of the $`D=4`$ RN–AdS space, which has not been tackled up to now due to the complicated structure of this system, in terms of the GEMS approach covering the usual Kruskal extension . In Sec.2, we discuss the $`D=4`$ RN–AdS embedding into a seven dimensional flat space. In Sec.3, we show that our results in the GEMS of the RN–AdS space systematically include those of the various limiting GEMS geometries, which are the RN, Schwarzschild–AdS, Schwarzschild, purely charged and AdS space-times, through the successive truncation procedure of parameters in the original curved space. These correspond to the dimensional reduction in the GEMS approach. Finally, we present summary in Sec.4.
## 2 Geometric Structure of RN-AdS <br>in the GEMS Approach
Let us consider the line element of the four dimensional RN–AdS space<sup>4</sup><sup>4</sup>4We restrict our discussion to the non-extremal case.
$$ds_4^2=f(r,m,e,R)dt^2f^1(r,m,e,R)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(1)
where $`f(r,m,e,R)`$ is given by
$$f(r,m,e,R)=1\frac{2m}{r}+\frac{e^2}{r^2}+\frac{r^2}{R^2}.$$
(2)
This space-time is asymptotically described by AdS, and there is an outer horizon at $`r=r_H`$. The case of $`e=0`$ yields the Schwarzschild–AdS metric, the case of $`m=e=0`$ yields the metric on the universal covering space of AdS , the case of $`R\mathrm{}`$ yields the RN metric, and the case of $`m=0`$ and $`R\mathrm{}`$ yields the purely charged metric.
To embed this space-time into a higher dimensional flat one, we first note that by introducing three coordinates $`(z^3,z^4,z^5)`$ in Eq. (2) (see below) the last term in the metric (1) can be written to give $`(dz^3)^2(dz^4)^2(dz^5)^2=dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)`$. Then, making use of an ansatz of two coordinates ($`z^0`$, $`z^1`$) in Eq. (2), we have obtained
$`(dz^0)^2(dz^1)^2(dz^3)^2(dz^4)^2(dz^5)^2`$
$`=f(r,m,e,R)dt^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)\left(1+{\displaystyle \frac{(\frac{f^{}(r,m,e,R)}{2})^2}{k_H^2f(r,m,e,R)}}\right)dr^2,`$
where $`f^{}(r,m,e,R)`$ denotes the derivative with respect to $`r`$ and
$$k_H(r_H,e,R^2)=\frac{(r_H^2e^2+\frac{3r_H^4}{R^2})}{2r_H^3}>0$$
(4)
is the surface gravity at the root of $`f(r,m,e,R)_{r=r_H}=0`$. In order to make the form of $`ds_4^2`$ in Eq. (1), we subtract the $`f^1(r,m,e,R)dr^2`$ term from Eq. (2) on the right-hand side and add it again to Eq. (2). Then, the remaining extra radial part of
$$f^1(r,m,e,R)dr^2\left(1+\frac{(\frac{f^{}(r,m,e,R)}{2})^2}{k_H^2f(r,m,e,R)}\right)dr^2$$
(5)
can be seperated into positive and negative definite parts with $`r>r_H`$ as follows:
$`R({\displaystyle \frac{e^2}{[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}}`$
$`+{\displaystyle \frac{r_H^2(r^2+rr_H+r_H^2)[(r_H^2e^2)^2R^4+r_H^6(r_H^2+2R^2)]}{r^2[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}})dr^2`$
$`\left(e^2{\displaystyle \frac{R^4r_H^6[4(rr_He^2)R^2+10r^4+2rr_H(r^2+rr_H+2r_H^2)]}{r^4[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}}\right)dr^2`$
$`\left({\displaystyle \frac{rr_H(r^2+rr_H+r_H^2)(4r_H^6R^2+[3r_H^4+(r_H^2e^2)R^2]^2)}{[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}}\right)dr^2`$
$`=(dz^2)^2(dz_e)^2(dz_R)^2,`$ (6)
where we have used the relation between the Arnowitt–Deser–Misner mass of the RN–AdS black hole and its event horizon radius $`r=r_H`$, i.e., $`2m=r_H+r_H^3/R^2+e^2/r_H`$. At this stage, it should be note that due to the existence of the last two terms, $`e`$-sensitive $`(dz_e)^2`$ and $`R`$-dominant $`(dz_R)^2`$, one may think that superficially two additional time dimensions are needed for a global flat embedding. However, it is in fact enough to introduce only one time dimension $`(dz^6)^2`$ by combining these two terms as
$$(dz^6)^2=(dz_e)^2+(dz_R)^2,$$
(7)
for a desired minimal GEMS<sup>5</sup><sup>5</sup>5In the region of $`r>r_H`$, it can be easily verified that the $`(dz^2)^2`$ and $`(dz^6)^2`$ are positive definite functions, when combined with the condition in Eq. (4). with an additional spacelike dimension $`(dz^2)^2`$. Note also that the $`(dz_e)^2`$ (or, $`(dz_R)^2`$) term is shown to be vanished in the limit of $`e0`$ (or, $`R\mathrm{}`$), and the $`(dz^6)^2`$ becomes $`(dz_R)^2`$ (or, $`(dz_e)^2`$). As a result, we have obtained a flat global embedding in (5+2)-dimensions of the corresponding curved 4-metric as
$`ds_7^2`$ $`=`$ $`(dz^0)^2{\displaystyle \underset{i=1}{\overset{5}{}}}(dz^i)^2+(dz^6)^2`$
$`=`$ $`f(r,m,e,R)dt^2f^1(r,m,e,R)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)`$
$`=`$ $`ds_4^2.`$ (9)
This equivalence between the (5+2)-dimensional flat embedding space and original curved space is the very definition of isometric embedding, mathematically developed by several authors .
It seems appropriate to comment on the lowest embedding dimensions in terms of the number of parameters. It is known from the previous works that whenever one parameter is increased in the original space, the embedding dimensions are either unchanged or increased depending on this parameter. In particular, the embedding dimension is already $`D=7`$ for the case of the RN or Schwarzschild–AdS, which have one less parameters than those of the RN–AdS case. Therefore, for the case of the RN–AdS the possibly lowest embedding dimension is $`D=7`$.
In summary, through the GEMS approach which makes the curved space possibly embedded in a higher dimensional flat space , we have found a $`D=7`$ dimensional isometric embedding of the RN–AdS space as
$`z^0=k_H^1\sqrt{f(r,m,e,R)}\mathrm{sinh}(k_Ht),`$
$`z^1=k_H^1\sqrt{f(r,m,e,R)}\mathrm{cosh}(k_Ht),`$
$`z^2={\displaystyle }drR({\displaystyle \frac{e^2}{[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}}`$
$`+{\displaystyle \frac{r_H^2(r^2+rr_H+r_H^2)[(r_H^2e^2)^2R^4+r_H^6(r_H^2+2R^2)]}{r^2[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}})^{1/2},`$
$`z^3=r\mathrm{sin}\theta \mathrm{cos}\varphi ,`$
$`z^4=r\mathrm{sin}\theta \mathrm{sin}\varphi ,`$ (10)
$`z^5=r\mathrm{cos}\theta ,`$
$`z^6={\displaystyle }dr({\displaystyle \frac{e^2R^4r_H^6[4(rr_He^2)R^2+10r^4+2rr_H(r^2+rr_H+2r_H^2)]}{r^4[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}}`$
$`+{\displaystyle \frac{rr_H(r^2+rr_H+r_H^2)(4r_H^6R^2+[3r_H^4+(r_H^2e^2)R^2]^2)}{[3r_H^4+(r_H^2e^2)R^2]^2[rr_H(r^2+rr_H+r_H^2)+(rr_He^2)R^2]}})^{1/2},`$
with an additional spacelike $`z^2`$ and a timelike $`z^6`$ dimensions. Therefore, the (3+1)-dimensional curved space is seen as the hyperboloid embedded in a (5+2)-dimensional flat space. It would be easily verified inversely that the flat metric (2) in the (5+2)-dimensional space defined as the coordinates (2) gives the original RN–AdS metric (1) correctly.
Now, following the trajectory of $`z^2=\mathrm{}=z^6=0`$ in Eq. (2) which corresponds to a static trajectory ($`r,\theta ,\varphi =\mathrm{constant}`$) in the curved space, the relevant $`D=7`$ acceleration $`a_7`$ is described as the Rindler-like motion of the form of $`(z^1)^2(z^0)^2=a_7^2`$ in the embedded flat space, i.e.,
$$a_7=\{(z^1)^2(z^0)^2\}^{1/2}=\frac{r_H^2e^2+\frac{3r_H^4}{R^2}}{2r_H^3\sqrt{f(r,m,e,R)}}.$$
(11)
As a result, the detector of the above Rindler-like motion would measure the correct Hawking temperature through the relation of $`T=a_7/2\pi `$ as follows
$$T=\frac{r_H^2e^2+\frac{3r_H^4}{R^2}}{4\pi r_H^3\sqrt{f(r,m,e,R)}},$$
(12)
in the GEMS approach. Then, the desired BH temperature is given by
$$T_0=\sqrt{g_{00}}T=\frac{r_H^2e^2+\frac{3r_H^4}{R^2}}{4\pi r_H^3}.$$
(13)
It is by now well–known that entropy, which is the extensive companion of the temperature, is given by one quarter of the horizon area . On the other hand, R. Laflamme showed that entropy seen by an accelerated observer in Minkowski space can be obtained from the consideration of the transverse area of a null surface on the wedge. This transverse area would diverge for otherwise unrestricted Rindler motion due to the integration over the whole transverse dimensions. In an embedded higher dimensional flat space, however, since there are “embedding” constraints, the resulting integral may not be divergent and make entropy finite.
Our RN–AdS case, where there are three additional dimensions in the transverse area, $`𝑑z^2\mathrm{}𝑑z^6`$, is correspondingly subject to four constraints as follows
$`(z^1)^2(z^0)^2=0,`$ (14)
$`z^2=f_1(r),z^6=f_2(r),`$ (15)
$`(z^3)^2+(z^4)^2+(z^5)^2=r^2,`$ (16)
where $`f_i(r)`$ are explicitly given in Eq. (2). Note that Eq. (14) leads to $`r=r_H`$. Since the $`z^2`$ and $`z^6`$ integrals subject to the constraints (15), $`𝑑z^2𝑑z^6\delta (z^2f_1(r))\delta (z^6f_2(r))`$, is unity, the remaining integrals of $`z_i(i=3,4,5)`$ well reproduce the desired area $`4\pi r_H^2`$ of the $`r=r_H`$ sphere. This ends the global flat embedding of the RN–AdS space giving the correct thermodynamics.
## 3 Various Limiting Geometries
Now, we are ready to analyze the various limiting geometries through the successive truncation procedure of the parameters, $`e`$, or $`R`$ (or, both) in the original curved space.
### 3.1 RN limit
Let us first consider the RN limit , which is the case of $`R\mathrm{}`$ in the metric (1),
$$ds_4^2=f_e(r,m,e)dt^2f_e^1(r,m,e)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(17)
where
$$f_e(r,m,e)=f(r,m,e,R\mathrm{})1\frac{2m}{r}+\frac{e^2}{r^2}.$$
(18)
The global flat embedding coordinates can be obtained either in the GEMS approach starting from the RN metric, Eq. (17), or in the limit of $`R\mathrm{}`$ from Eq. (2) directly as
$`z^0=k_H^1\sqrt{f_e(r,m,e)}\mathrm{sinh}(k_Ht),`$
$`z^1=k_H^1\sqrt{f_e(r,m,e)}\mathrm{cosh}(k_Ht),`$
$`z^2={\displaystyle 𝑑r\left(\frac{r^2(r_++r_{})+r_+^2(r+r_+)}{r^2(rr_{})}\right)^{1/2}},`$
$`z^3=r\mathrm{sin}\theta \mathrm{cos}\varphi ,`$
$`z^4=r\mathrm{sin}\theta \mathrm{sin}\varphi ,`$
$`z^5=r\mathrm{cos}\theta ,`$
$`z^6={\displaystyle 𝑑r\left(\frac{4r_+^5r_{}}{r^4(r_+r_{})^2}\right)^{1/2}},`$ (19)
where the surface gravity is given by $`k_H=k_H(r_H,e,\mathrm{})=(r_+r_{})/2r_+^2`$ with the outer horizon $`r_+=r_H`$, and $`r_\pm =m\pm \sqrt{m^2e^2}`$. In this limit, the $`R`$-dominant part of $`z^6`$ in Eq. (2) vanishes and the resulting GEMS becomes exactly the known $`D=7`$ RN one . Note that in the limit of $`R\mathrm{}`$ the corresponding event horizon becomes the usual RN one by rewriting the charge $`e^2`$ to $`r_+r_{}`$.
Moreover, the relevant $`D=7`$ acceleration and the Hawking temperature can be obtained either directly from Eqs. (11) and (12) by taking the limit of $`R\mathrm{}`$ and replacing $`e^2`$ with $`r_+r_{}`$, or from the Rindler-like motion in the $`D=7`$ GEMS, Eq. (3.1), following a static trajectory ($`r,\theta ,\varphi =\mathrm{constrant}`$) in the curved space as before,
$`a_7`$ $`=`$ $`\{(z^1)^2(z^0)^2\}^{1/2}={\displaystyle \frac{r_+r_{}}{2r_+^2\sqrt{f_e(r,m,e)}}},`$ (20)
$`T`$ $`=`$ $`{\displaystyle \frac{r_+^2e^2}{4\pi r_+^3\sqrt{f_e(r,m,e)}}}={\displaystyle \frac{r_+r_{}}{4\pi r_+^2\sqrt{f_e(r,m,e)}}}.`$ (21)
The entropy calculation of the RN is essentially the same as the previous RN–AdS case. In this case there are three additional dimensions, and four constraints, i.e., $`(z^1)^2(z^0)^2=0`$ leads to $`r=r_+`$, $`z^2=f_1(r,R\mathrm{})`$, $`z^6=f_2(r,R\mathrm{})`$ in Eqs. (2) and $`(z^3)^2+(z^4)^2+(z^5)^2=r^2`$. Thus, since the $`z^2,z^6`$ integrals, $`𝑑z^2𝑑z^6\delta (z^2f_1(r))\delta (z^6f_2(r))`$, are unity, the remaining integrals give the desired area 4$`\pi r_H^2`$, that of the corresponding $`r=r_H`$ sphere.
### 3.2 Schwarzschild-AdS limit
Secondly, the RN–AdS solution (2) is also easily reduced to the Schwarzschild–AdS space, which is the limiting case of $`e0`$,
$$ds_4^2=f_R(r,m,R)dt^2f_R^1(r,m,R)dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(22)
where
$$f_R(r,m,R)=f(r,m,e=0,R)1\frac{2m}{r}+\frac{r^2}{R^2},$$
(23)
giving another $`D=7`$ GEMS with the vanishing $`e`$-sensitive part of $`z^6`$ in Eq. (2),
$`z^0=k_H^1\sqrt{f_R(r,m,R)}\mathrm{sinh}(k_Ht),`$
$`z^1=k_H^1\sqrt{f_R(r,m,R)}\mathrm{cosh}(k_Ht),`$
$`z^2={\displaystyle 𝑑r\frac{R^3+Rr_H^2}{R^2+3r_H^2}\sqrt{\frac{r_H(r^2+rr_H+r_H^2)}{r^3(r^2+rr_H+r_H^2+R^2)}}},`$
$`z^3=r\mathrm{sin}\theta \mathrm{cos}\varphi ,`$
$`z^4=r\mathrm{sin}\theta \mathrm{sin}\varphi ,`$
$`z^5=r\mathrm{cos}\theta ,`$
$`z^6={\displaystyle 𝑑r\sqrt{\frac{(R^4+10R^2r_H^2+9r_H^4)(r^2+rr_H+r_H^2)}{(R^2+3r_H^2)^2(r^2+rr_H+r_H^2+R^2)}}}.`$ (24)
The surface gravity, $`k_H=k_H(r_H,0,R)=(R^2+3r_H^2)/2r_HR^2`$, is now either obtained at the root $`r_H`$ of $`f_R(r,m,R)_{r=r_H}=0`$, or reduced directly from the Eq. (4) with $`e=0`$. This seemingly complicated embedding space is firstly obtained in Ref. , and we have also reached to the exactly same results by the systematic reduction process from Eq. (2).
On the other hand, similar to the RN limit case, we directly obtain the Hawking temperature from Eqs. (11) and (12) by taking the limit of $`e0`$ as follows
$$T=\frac{a_7}{2\pi }=\frac{1+\frac{3r_H^2}{R^2}}{4\pi r_H\sqrt{f_R(r,m,R)}},$$
(25)
which again equals to that calculated in .
### 3.3 Schwarzschild limit
Thirdly, we can obtain the Schwarzschild limit without the cosmological constant from the RN embedding of (3.1) with $`e0`$ limit or the Schwarzschild–AdS embedding of (3.2) with $`R\mathrm{}`$ one. As a result, it is successfully reduced to the $`D=6`$ flat GEMS as follows ,
$`z^0=k_H^1\sqrt{12m/r}\mathrm{sinh}(k_H^1t),`$
$`z^1=k_H^1\sqrt{12m/r}\mathrm{cosh}(k_H^1t),`$
$`z^2={\displaystyle 𝑑r\sqrt{r_H(r^2+rr_H+r_H^2)/r^3}},`$
$`z^3=r\mathrm{sin}\theta \mathrm{sin}\varphi ,`$
$`z^4=r\mathrm{sin}\theta \mathrm{cos}\varphi ,`$
$`z^5=r\mathrm{cos}\theta ,`$ (26)
where the event horizon is $`r_H=2m`$, and the surface gravity is $`k_H(r_H,0,\mathrm{})=1/2r_H`$. Note that the analyticity of $`z^2(r)`$ in $`r>0`$ covers the region of $`r<r_H`$. Thus, it should be cautioned that the use of incomplete embedding spaces, that cover only $`r>r_H`$ (as, for example, in ), will lead to observers there for whom there is no event horizon, no loss of information, and no temperature.
We then obtain the Hawking temperature from Eqs. (20) and (21) by taking the limit $`e0`$ as follows
$`T={\displaystyle \frac{a_6}{2\pi }}={\displaystyle \frac{1}{8\pi m\sqrt{12m/r}}},`$
$`T_0=\sqrt{g_{00}}T={\displaystyle \frac{1}{8\pi m}}.`$ (27)
It seems appropriate to comment on a global flat embedding of $`D=4`$ covering of the AdS,
$$ds_4^2=(1+\frac{r^2}{R^2})dt^2(1+\frac{r^2}{R^2})^1dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(28)
which corresponds to the case of $`m0`$ in Eq. (22). In this case we cannot directly obtain a global embedding from Eq. (3.2) since in the limit of $`m0`$ the surface gravity $`k_H=1/2r_H=1/4m`$ yields a divergence. As discussed in Ref. in details, this problem originally comes from the fact that there is no intrinsic horizon of this space-time. However, there is of course the other direct route to embed this space-time into the $`D=5`$ flat space-time starting from the metric (22) with $`m=0`$. Based on the accelerating coordinate system, the correct temperature of $`2\pi T=(a^2R^2)^{1/2}`$ has been already found (For further details, see Ref. ).
Similar to the pure AdS case, we can directly analyze the purely charged case with $`m=0`$ in the metric (17) as
$$ds_4^2=(1+\frac{e^2}{r^2})dt^2(1+\frac{e^2}{r^2})^1dr^2r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(29)
As like the above $`D=4`$ covering of the AdS, this has also no event horizon. However, we can also embed this space-time into $`D=5`$ flat one in view of the accelerating coordinate frame as follows
$`z^0`$ $`=`$ $`\sqrt{\rho ^2e^2}\mathrm{sinh}(\eta /e),`$
$`z^1`$ $`=`$ $`\sqrt{\rho ^2e^2}\mathrm{cosh}(\eta /e),`$
$`z^2`$ $`=`$ $`\rho \mathrm{sinh}\mathrm{\Phi }\mathrm{cos}\theta ,`$
$`z^3`$ $`=`$ $`\rho \mathrm{sinh}\mathrm{\Phi }\mathrm{sin}\theta ,`$
$`z^4`$ $`=`$ $`\rho \mathrm{cosh}\mathrm{\Phi },`$ (30)
where $`\mathrm{}<\eta ,\mathrm{\Phi }<\mathrm{},\pi <\theta <\pi `$. While this coordinate patch only covers the region $`\rho >e`$, it can be extended to the entire space similar to the four dimensional AdS case . Then, we can easily obtain the temperature as $`2\pi T=(a^2e^2)^{1/2}`$ where the four acceleration $`a`$ is given by $`a=\rho ^2/e^2(\rho ^2e^2)`$.
Furthermore, if we take the limit $`R\mathrm{}`$ in the metric (28), or $`e0`$ in the metric (29), we finally reach to the flat four dimensional Minkowski space.
We have summarized all these results in Fig. 1 as a compact diagram, which can be obtained through the systematic truncation procedure from the seven dimensional flat embedding space.
## 4 Summary
In summary, we have shown that the Hawking thermal properties map into their Unruh equivalents in the (5+2)-dimensional GEMS, which is the lowest possible global embedding dimensions of the curved RN–AdS space. The relevant curved space detectors become Rindler ones, whose temperatures and entropies reproduce the originals. Our results of the RN–AdS in the GEMS approach include the various limiting geometries, which are the Reissner–Nordström, Schwarzschild–AdS, and Schwarzschild space-times through the successive reduction procedure of the parameters in the original space. As a result, the (5+2)-dimensional GEMS in Eq. (2) serves a unifying description of the global flat embedding of the various geometries. It would be interesting to consider other interesting applications of the GEMS, for example, the rotating Kerr type geometries .
We are grateful for interesting discussions to Prof. G. W. Gibbons. Y.W.K. acknowledges financial support from KOSEF, Y.J.P. from the Ministry of Education, BK21 Project No. D-0055/99, and K.S. for S.N.U. CTP and Ministry of Education for BK-21 Project. |
warning/0001/astro-ph0001537.html | ar5iv | text | # THE EQUATION OF STATE OF NEUTRON-STAR MATTER IN STRONG MAGNETIC FIELDS
## 1 INTRODUCTION
Recent observational and theoretical studies motivate the investigation of the effects of ultra-strong magnetic fields ($`B>10^{14}`$ Gauss) on neutron stars. Several independent arguments link the class of soft $`\gamma `$ray repeaters and perhaps certain anomalous X-ray pulsars with neutron stars having ultra strong magnetic fields – the so-called magnetars (Paczyński 1992; Thompson & Duncan 1995, 1996; Melatos 1999). In addition, two of the four known soft $`\gamma `$ray repeaters directly imply, from their periods and spin-down rates, surface fields in the range $`28\times 10^{14}`$ Gauss. Kouveliotou et al. (1998, 1999) argue from the population statistics of soft $`\gamma `$ray repeaters that magnetars constitute about 10% of the neutron star population. While some observed white dwarfs have large enough fields to give ultra-strong neutron star magnetic fields through flux conservaton, it does not appear likely that such isolated examples could account for a significant fraction of ultra-strong field neutron stars. Therefore, an alternative mechanism seems necessary for the creation ultra-strong magnetic fields in neutron stars. Duncan & Thompson (1992, 1996) suggested that large fields (up to $`3\times 10^{17}\times (1\mathrm{ms}/\mathrm{P}_\mathrm{i})`$ Gauss, where $`P_i`$ is the initial rotation period) can be generated in nascent neutron stars through the smoothing of differential rotation and convection.
These developments raise the intriguing questions:
(1) What is the largest frozen-in magnetic field a stationary neutron star can sustain?, and,
(2) What is the effect of such ultra-strong magnetic fields on the maximum neutron star mass?
The answers to both of these questions hinge upon the effects strong magnetic fields have both on the equation of state (EOS) of neutron-star matter and on the structure of neutron stars. In this paper, we will focus on the effects of strong magnetic fields on the EOS. Subsequent work will be devoted to investigating the effects of strong fields on the structure of neutron stars, incorporating the EOSs developed in this work.
The magnitude of the magnetic field strength $`B`$ needed to dramatically affect neutron star structure directly can be estimated with a dimensional analysis (Lai & Shapiro 1991) equating the magnetic field energy $`E_b(4\pi R^3/3)(B^2/8\pi )`$ with the gravitational binding energy $`E_{B.E.}GM^2/R`$, yielding $`B2\times 10^{18}\left(M/1.4\mathrm{M}_{}\right)\left(R/10\mathrm{km}\right)^2`$ Gauss, where $`M`$ and $`R`$ are, respectively, the neutron star mass and radius.
The magnitude of $`B`$ required to directly influence the EOS can be estimated by considering its effects on charged particles. Charge neutral, beta-equilibrated, neutron-star matter contains both negatively charged leptons (electrons and muons) and positively charged protons. Magnetic fields quantize the orbital motion (Landau quantization) of these charged particles. Relativistic effects become important when the particle’s cyclotron energy $`e\mathrm{}B/(mc)`$ is comparable to it’s mass (times $`c^2`$). The magnitudes of the so-called critical fields are $`B_c^e=(\mathrm{}c/e)\text{ }\text{-}\lambda _e^2=4.414\times 10^{13}`$ Gauss and $`B_c^p=(m_p/m_e)^2B_c^e=1.487\times 10^{20}`$ Gauss for the electron and proton, respectively ( $`\text{ }\text{-}\lambda _e=\mathrm{}/m_ec386`$ fm is the Compton wavelength of the electron). It will be convenient to measure the field strength $`B`$ in units of $`B_e^c`$, viz., $`B^{}B/B_e^c`$. When the Fermi energy of the proton becomes significantly affected by the magnetic field, the composition of matter in beta equilibrum is significantly affected. In turn, the pressure of matter is significantly affected. We show that this occurs when $`B^{}10^5`$, and will lead to a general softening of the EOS.
In neutron stars, magnetic fields may well vary in strength from the core to the surface. The scale lengths of such variations are, however, usually much larger than the microscopic magnetic scale $`l_m`$, which depends on the magnitude of $`B`$. For low fields, for which the quasi-classical approximation holds, $`l_m(\text{ }\text{-}\lambda _e^2/B^{})(3\pi ^2n_e)^{1/3}10^5(n_e/n_s)^{1/3}/B^{}`$ fm, where $`n_e`$ is the number density of electrons and $`n_s`$ is the normal nuclear saturation density (about 0.16 fm<sup>-3</sup>). For high fields, when only a few Landau levels are occupied, $`l_m2\pi ^2n_e(\text{ }\text{-}\lambda _e^2/B^{})^27\times 10^9(n_e/n_s)/B^2`$ fm. In either case, the requirement that $`R>>l_m`$ is amply satisfied; hence, the magnetic field $`B`$ may be assumed to be locally constant and uniform as far as effects on the EOS are concerned.
In non-magnetic neutron stars, the pressure of matter ranges from $`25\mathrm{MeV}\mathrm{fm}^3`$ at nuclear density to $`200600\mathrm{MeV}\mathrm{fm}^3`$ at the central density of the maximum mass configuration, depending on the EOS (Prakash et al. 1997). These values may be contrasted with the energy density and pressure from the electromagnetic field: $`\epsilon _f=P_f=B^2/(8\pi )=4.814\times 10^8B^2\mathrm{MeV}\mathrm{fm}^3`$. The field contributions can dominate the matter pressure for $`B^{}>10^4`$ at nuclear densities and for $`B^{}>10^5`$ at the central densities of neutron stars, and must therefore be included whenever the field dramatically influences the star’s composition and matter pressure.
In strong magnetic fields, contributions from the anomalous magnetic moments of the nucleons must also be considered. Experimentally, $`\kappa _p=\mu _N\left(g_p/21\right)`$ for the proton, $`\kappa _n=\mu _Ng_n/2`$ for the neutron, where $`\mu _N`$ is the nuclear magneton and $`g_p=5.58`$ and $`g_n=3.82`$ are the Landé g-factors for the proton and neutron, respectively. The energy $`|\kappa _n+\kappa _p|B1.67\times 10^5B^{}`$ MeV measures the changes in the beta equilibrium condition and to the baryon Fermi energies. Since the Fermi energies range from a few MeV to tens of MeV for the densities of interest, it is clear that contributions from the anomalous magnetic moments also become significant for $`B^{}>10^5`$. We demonstrate that for such fields, complete spin polarization of the neutrons occurs, which results in an overall stiffening of the EOS that overwhelms the softening induced by Landau quantization.
In magnetized matter, the stress energy tensor contains terms proportional to $`HB`$, where $`H=B+4\pi `$ and $``$ is the magnetization (Landau, Lifshitz & Pitaevskiĭ 1984). Thus, extra terms, in addition to the usual ones proportional to $`B^2`$, are introduced into the structure equations (Cardall et al. 1999). The magnetization in a single component electron gas has been studied extensively (Blandford & Hernquist 1982) for neutron star crust matter. We generalize this formulation to the case of interacting multicomponent matter with and without the effects of the anomalous magnetic moments. We find that deviations of $`H`$ from $`B`$ occur for field strengths $`B^{}10^5`$.
Although the effects of magnetic fields on the EOS at low densities, relevant for neutron star crusts, has been extensively studied (see for example, Canuto & Ventura 1977; Fushiki, Gudmundsson & Pethick 1989; Fushiki et al. 1992; Abrahams & Shapiro 1991; Lai & Shapiro 1991; Rögnvaldsson et al. 1993, Thorlofsson et al. 1998), only a handful of previous works have considered the effects of very large magnetic fields on the EOS of dense neutron star matter (Chakrabarty 1996; Chakrabarty, Bandyopadhyay, & Pal 1997, Yuan & Zhang 1999). Lai and Shapiro (1991) considered non-interacting, charge neutral, beta-equilibrated matter at subsaturation densities, while Chakrabarty and co-authors studied dense matter including interactions using a field-theoretical description. These authors found large compositional changes in matter induced by ultra-strong magnetic fields due to the quantization of orbital motion. Acting in concert with the nuclear symmetry energy, Landau quantization substantially increases the concentration of protons compared to the field-free case, which in turn leads to a softening of the EOS. This lowers the maximum mass relative to the field-free value. In these works, however, the electromagnetic field energy density and pressure, which tend to stiffen the EOS, were not included. In addition, changes in the general relativistic structure induced by the magnetic fields (studied in detail by Bocquet et al. 1995 who, however, omitted the compositional changes in the EOS due to Landau quantization) were also ignored. Thus, the combined effects of the magnetic fields on the EOS and on the general relativistic structure remain to be determined.
Compared to these earlier works, we make several improvements in the calculation of the EOS. These improvements include (1) a study of a larger class of field-theoretical models in order to extract the generic trends induced by Landau quantization, (2) the development of a covariant description for the inclusion of the anomalous magnetic moments of the nucleons, and (3) a detailed study of magnetization of interacting multicomponent matter with and without the inclusion of the anomalous magnetic moments. We also provide simple analytical estimates of when each of these effects begin to significantly influence the EOS. Our future work will employ the EOSs developed in this work to complete a fully self-consistent calculation of neutron star structure including the combined effects of the direct effects of magnetic fields on the EOS and general relativistic structure.
In §2, we present the field-theoretical description of dense neutron star matter including the effects of Landau quantization and the nucleon anomalous magnetic moments. Section 3 contains a detailed study of the effects of Landau quantization on the EOS for two classes of Lagrangians. In addition to providing contrasts with earlier work, our results highlight the extent to which the underlying interactions affect the basic findings. This section also includes new theoretical developments concerning the magnetization of interacting, multicomponent matter. Section 4 is devoted to the effects of the anomalous magnetic moments on the EOS. Here results for a charge neutral neutron, proton, electron, and muon gas are compared with those for interacting matter to asses the generic trends. Our conclusions and outlook, including the possible effects of additional components such as hyperons, Bose condensates and quarks, are presented in §5. The covariant description for the inclusion of the anomalous magnetic moments of the nucleons is presented in the Appendix, where explicit formulae for the nucleon Dirac spinors and energy spectra are derived. Except where necessary, we use units wherin $`\mathrm{}`$ and $`c`$ are set to unity.
## 2 THEORETICAL FRAMEWORK
For the description of the EOS of neutron-star matter, we employ a field-theoretical approach in which the baryons (neutrons, $`n`$, and protons, $`p`$) interact via the exchange of $`\sigma \omega \rho `$ mesons. We study two classes of models, which differ in their behavior at high density. The Lagrangian densities associated with these two classes are (Boguta & Bodmer 1977, Zimanyi & Moszkowski 1990)
$`_I`$ $`=`$ $`_b\left(1{\displaystyle \frac{g_{\sigma _b}\sigma }{m_b}}\right)\overline{\mathrm{\Psi }}_bm_b\mathrm{\Psi }_b+_m+_l,`$
$`_{II}`$ $`=`$ $`\left(1+{\displaystyle \frac{g_{\sigma _b}\sigma }{m_b}}\right)_b\overline{\mathrm{\Psi }}_bm_b\mathrm{\Psi }_b+_m+_l.`$ (1)
The baryon ($`b=n,p`$), lepton ($`l=e,\mu `$), and meson ($`\sigma ,\omega ,\mathrm{and}\rho `$) Lagrangians are given by
$`_b`$ $`=`$ $`\overline{\mathrm{\Psi }}_b\left(i\gamma _\mu ^\mu +q_b\gamma _\mu A^\mu g_{\omega _b}\gamma _\mu \omega ^\mu g_{\rho _b}\tau _{3_b}\gamma _\mu \rho ^\mu \kappa _b\sigma _{\mu \nu }F^{\mu \nu }\right)\mathrm{\Psi }_b,`$
$`_l`$ $`=`$ $`\overline{\psi }_l\left(i\gamma _\mu ^\mu +q_l\gamma _\mu A^\mu \right)\psi _l,`$
$`_m`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \sigma ^\mu \sigma {\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2U(\sigma )+{\displaystyle \frac{1}{2}}m_\omega ^2\omega _\mu \omega ^\mu {\displaystyle \frac{1}{4}}\mathrm{\Omega }^{\mu \nu }\mathrm{\Omega }_{\mu \nu }`$ (2)
$`+{\displaystyle \frac{1}{2}}m_\rho ^2\rho _\mu \rho ^\mu {\displaystyle \frac{1}{4}}P^{\mu \nu }P_{\mu \nu }{\displaystyle \frac{1}{4}}F^{\mu \nu }F_{\mu \nu },`$
where $`\mathrm{\Psi }_b`$ and $`\psi _l`$ are the baryon and lepton Dirac fields, respectively. The nucleon mass and the isospin projection are denoted by $`m_b`$ and $`\tau _{3_b}`$, respectively. The mesonic and electromagnetic field strength tensors are given by their usual expressions: $`\mathrm{\Omega }_{\mu \nu }=_\mu \omega _\nu _\nu \omega _\mu `$, $`P_{\mu \nu }=_\mu \rho _\nu _\nu \rho _\mu `$, and $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$. The strong interaction couplings are denoted by $`g`$, the electromagnetic couplings by $`q`$, and the meson masses by $`m`$ all with appropriate subscripts. The anomalous magnetic moments are introduced via the coupling of the baryons to the electromagnetic field tensor with $`\sigma _{\mu \nu }=\frac{i}{2}[\gamma _\mu ,\gamma _\nu ]`$ and strength $`\kappa _b`$. We will contrast results for cases with $`\kappa _b=0`$ and $`\kappa _b`$ taken to be their measured values. The quantity $`U(\sigma )`$ denotes possible scalar self-interactions. It is straightforward to include self interactions between both the vector $`\omega `$ and the iso-vector $`\rho `$ mesons (Müller & Serot 1996). Although the electromagnetic field is included in $`_I`$ and $`_{II}`$, it assumed to be externally generated (and thus has no associated field equation) and only frozen-field configurations will be considered.
The thermodynamic quantities will be evaluated in the mean field approximation, in which the mesonic fields are assumed to be constant. The field equations are
$$m_\sigma ^2\sigma +\frac{U(\sigma )}{\sigma }=\{\begin{array}{cc}_bg_{\sigma _b}n_b^s& \text{for }_I\hfill \\ _bg_{\sigma _b}\left(\frac{m_b^{}}{m_b}\right)^2n_b^s& \text{for }_{II}\text{ ,}\hfill \end{array}$$
(3)
$$m_\omega ^2\omega ^0=\underset{b}{}g_{\omega _b}n_b,$$
(4)
$$m_\rho ^2\rho ^0=\underset{b}{}g_{\rho _b}\tau _{3_b}n_b,$$
(5)
$$\left[\stackrel{}{\alpha }\stackrel{}{}\left(\stackrel{}{p}q_l\stackrel{}{A}\right)+\beta m_l\right]\psi _l=E\psi _l,$$
(6)
$$\left[\stackrel{}{\alpha }\stackrel{}{}\left(\stackrel{}{p}q_b\stackrel{}{A}\right)+\beta m_b^{}\right]\mathrm{\Psi }_b=(Eg_{\omega _b}\omega ^0g_{\rho _b}\tau _{3_b}\rho ^0)\mathrm{\Psi }_b,$$
(7)
where the effective baryon masses are
$$\frac{m_b^{}}{m_b}=\{\begin{array}{cc}\left(1\frac{g_{\sigma _b}\sigma }{m_b}\right)& \text{for }_I\hfill \\ \left(1+\frac{g_{\sigma _b}\sigma }{m_b}\right)^1& \text{for }_{II}\text{ ,}\hfill \end{array}$$
(8)
and $`n_b^s`$ is the scalar number density. The scalar self-interaction is taken to be of the form (Boguta & Bodmer 1977; Glendenning 1982, 1985)
$$U(\sigma )=\frac{1}{3}bm_n(g_{\sigma _n}\sigma )^3+\frac{1}{4}c(g_{\sigma _n}\sigma )^4,$$
(9)
where the $`m_n`$ in the first term is included to make $`b`$ dimensionless. In charge neutral, beta equilibrated matter, the conditions
$$\mu _n\mu _p=\mu _e=\mu _\mu ,$$
(10)
$$n_p=n_e+n_\mu ,$$
(11)
also apply. Given the nucleon-meson coupling constants and the coefficients in the scalar self-interaction, equations (3) through (11) may be solved self consistently for the chemical potentials, $`\mu _i`$, and the field strengths, $`\sigma `$, $`\omega ^0`$, and $`\rho ^0`$.
## 3 EFFECTS OF LANDAU QUANTIZATION
From equation (6), the energy spectrum for the leptons is (see, for example, Canuto & Ventura 1977)
$$E_l=\sqrt{k_z^2+\stackrel{~}{m}_{n,\sigma _z}^{l2}},$$
(12)
where
$$\stackrel{~}{m}_{n,\sigma _z}^{l2}=m_l^2+2\left(n+\frac{1}{2}\frac{1}{2}\frac{q_l}{|q_l|}\sigma _z\right)|q_l|B.$$
(13)
Here, $`n`$ is the principal quantum number and $`\sigma _z`$ (not to be confused with the scalar field $`\sigma `$) is the spin along the magnetic field axis. $`k_z`$ is the component of the momentum along the magnetic field axis. The quantity $`\nu =n+1/2(1/2)(q_i/|q_i|)\sigma _z`$ characterizes the so-called Landau level. Equation (7) gives the energy spectrum for the protons as
$$E_p=\sqrt{k_z^2+\stackrel{~}{m}_{n,\sigma _z}^{p2}}+g_{\omega _p}\omega ^0g_{\rho _p}\frac{1}{2}\rho ^0,$$
(14)
where $`\stackrel{~}{m}_{n,\sigma _z}^p`$ is obtained by replacing $`m_l`$ on the right hand side of equation (13) by $`m_p^{}`$. The neutron energy spectrum is that of the free Dirac particle, but with shifts arising from the scalar, vector, and isovector interactions:
$$E_n=\sqrt{k^2+m_n^2}+g_{\omega _n}\omega ^0+g_{\rho _n}\frac{1}{2}\rho ^0.$$
(15)
At zero temperature and in the presence of a constant magnetic field $`B`$, the number and energy densities of charged particles are given by
$$n_{i=l,p}=\frac{|q_i|B}{2\pi ^2}\underset{\sigma _z}{}\underset{n=0}{\overset{n_{max}}{}}k_{f,n,\sigma _z}^i,$$
(16)
$$\epsilon _{i=l,p}=\frac{|q_i|B}{4\pi ^2}\underset{\sigma _z}{}\underset{n=0}{\overset{n_{max}}{}}\left[E_f^ik_{f,n,\sigma _z}^i+\stackrel{~}{m}_{n,\sigma _z}^{i2}\mathrm{ln}\left(\left|\frac{E_f^i+k_{f,n,\sigma _z}^i}{\stackrel{~}{m}_{n,\sigma _z}^i}\right|\right)\right].$$
(17)
Above, $`k_{f,n,\sigma _z}^i`$ is the Fermi momentum for the level with the principal quantum number $`n`$ and spin $`\sigma _z`$ and is given by
$$k_{f,n,\sigma _z}^{i2}=E_f^{i2}\stackrel{~}{m}_{n,\sigma _z}^{i2}.$$
(18)
The summation in equation (16) is terminated at $`n_{max}`$, which is the integer preceeding the value of $`n`$ for which $`k_{f,n,\sigma _z}^{i2}`$ is negative. The Fermi energies are fixed by the chemical potentials
$$E_f^l=\mu _l,$$
(19)
$$E_f^b=\mu _bg_{\omega _b}\omega ^0g_{\rho _b}\tau _{3_b}\rho ^0.$$
(20)
For the protons, the scalar number density may be determined to be (Chakrabarty 1996)
$$n_p^s=\frac{|q_p|Bm_p^{}}{2\pi ^2}\underset{\sigma _z}{}\underset{n=0}{\overset{n_{max}}{}}\mathrm{ln}\left(\left|\frac{E_f^p+k_{f,n,\sigma _z}^p}{\stackrel{~}{m}_{n,\sigma _z}^p}\right|\right).$$
(21)
The number, energy, and scalar number densities of the neutrons are unchanged in form from the field-free case
$$n_n=\frac{k_f^{n3}}{3\pi ^2},$$
(22)
$$n_n^s=\frac{m_n^{}}{2\pi ^2}\left[E_f^nk_f^nm_n^2\mathrm{ln}\left(\left|\frac{E_f^n+k_f^n}{m_n^{}}\right|\right)\right],$$
(23)
$$\epsilon _n=\frac{1}{8\pi ^2}\left[2E_f^{n3}k_f^nm_n^2E_f^nk_f^nm_n^4\mathrm{ln}\left(\left|\frac{E_f^n+k_f^n}{m_n^{}}\right|\right)\right].$$
(24)
The total energy density of the system is
$`\epsilon `$ $`=`$ $`{\displaystyle \frac{1}{2}}m_\omega ^2\omega _0^2+{\displaystyle \frac{1}{2}}m_\rho ^2\rho _0^2+{\displaystyle \frac{1}{2}}m_\sigma ^2\sigma ^2+U(\sigma )`$ (25)
$`+{\displaystyle \underset{b}{}}\epsilon _b+{\displaystyle \underset{l}{}}\epsilon _l+{\displaystyle \frac{B^2}{8\pi ^2}},`$
where the last term is the contribution from the electromagnetic field. Use of equations (10) and (11), which are satisfied in charge neutral beta-equilibrated matter, in the general expression for the pressure, $`P=_i\mu _in_i\epsilon `$ ($`i=n,p,e,\mathrm{and}\mu `$), allows the pressure to be written only in terms of the neutron chemical potential through the relation $`P=\mu _nn_b\epsilon `$. In fact, utilizing the appropriate relations satisfied by the various chemical potentials and the number densities involved in the charge neutrality condition, it is easily verified that this relation is satisfied even in the presence of additional components such as strangeness-bearing hyperons, Bose condensates (pion or kaon), and quarks, which may likely exist in dense neutron-star matter.
### 3.1 Magnetization
The magnetic field strength, $`H`$, is related to the energy density by (Landau, Lifshitz & Pitaevskiĭ 1984)
$$H=4\pi \left(\frac{\epsilon }{B}\right)_{n_b}=B+4\pi ,$$
(26)
where $``$ is the magnetization. This is equivalent to the set of equations
$`{\displaystyle \frac{dn_b}{dB}}`$ $`=`$ $`0,`$
$``$ $`=`$ $`{\displaystyle \frac{\epsilon _l}{B}}+{\displaystyle \frac{\epsilon _l}{E_f^l}}{\displaystyle \frac{dE_f^l}{dB}}+{\displaystyle \frac{\epsilon _p}{B}}+{\displaystyle \frac{\epsilon _p}{E_f^p}}{\displaystyle \frac{dE_f^p}{dB}}+{\displaystyle \frac{\epsilon _n}{B}}+{\displaystyle \frac{\epsilon _n}{E_f^n}}{\displaystyle \frac{dE_f^n}{dB}}`$ (27)
$`+m_\rho ^2\rho ^0{\displaystyle \frac{d\rho ^0}{dB}}+m_\omega ^2\omega ^0{\displaystyle \frac{d\omega ^0}{dB}}+{\displaystyle \frac{\epsilon }{\sigma }}{\displaystyle \frac{d\sigma }{dB}}.`$
The first of these gives
$$\frac{n_n}{E_f^n}\frac{dE_f^n}{dB}=\frac{n_n}{B}\frac{n_p}{B}\frac{n_p}{E_f^p}\frac{dE_f^p}{dB}.$$
(28)
Using the conditions of charge neutrality and chemical equilibrium, one has
$$\frac{n_l}{E_f^l}\frac{dE_f^l}{dB}=\frac{n_l}{B}+\frac{n_p}{B}+\frac{n_p}{E_f^p}\frac{dE_f^p}{dB}.$$
(29)
From the field equations and the definition of the scalar density,
$`{\displaystyle \frac{d\omega ^0}{dB}}`$ $`=`$ $`{\displaystyle \frac{g_\omega }{m_\omega ^2}}{\displaystyle \frac{dn_b}{dB}}=0,`$
$`{\displaystyle \frac{d\rho ^0}{dB}}`$ $`=`$ $`{\displaystyle \frac{g_\rho }{m_\rho ^2}}{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{dB}}\left(n_nn_p\right)={\displaystyle \frac{g_\rho }{m_\rho ^2}}{\displaystyle \frac{dn_p}{dB}},`$
$`{\displaystyle \frac{\epsilon }{\sigma }}`$ $`=`$ $`0.`$ (30)
Note also that
$$\frac{\epsilon _i}{E_f^i}=E_f^i\frac{n_i}{E_f^i}.$$
(31)
Utilizing these results, equation (26) becomes
$$=T_I+T_{II}\frac{n_p}{E_f^p}\frac{dE_f^p}{dB},$$
(32)
where
$`T_I`$ $`=`$ $`{\displaystyle \frac{\epsilon _l}{B}}E_f^l{\displaystyle \frac{n_l}{B}}+{\displaystyle \frac{\epsilon _n}{B}}E_f^n{\displaystyle \frac{n_n}{B}}+{\displaystyle \frac{\epsilon _p}{B}}`$
$`+\left(E_f^lE_f^ng_\rho \rho ^0\right){\displaystyle \frac{n_p}{B}},`$
$`T_{II}`$ $`=`$ $`E_f^l+E_f^pE_f^ng_\rho \rho ^0.`$ (33)
Note that chemical equilibrium ensures that $`T_{II}=0`$ whence the magnetization takes the general form
$$=\underset{i=e,\mu ,p,n}{}\left(\frac{\epsilon _i}{B}E_f^i\frac{n_i}{B}\right).$$
(34)
In the case under current consideration, inserting the explicit forms of the energy density and number density yields the result
$$=\underset{i=e,\mu ,p}{}\left[\frac{\epsilon _iE_f^in_i}{B}+\frac{B}{2\pi ^2}\underset{\sigma _z}{}\underset{n=0}{\overset{n_{max}}{}}\left(n+\frac{1}{2}\frac{1}{2}\sigma _z\right)\mathrm{ln}\left(\left|\frac{E_f^i+k_{f,n,\sigma _z}^i}{\stackrel{~}{m}_{n,\sigma _z}^i}\right|\right)\right].$$
(35)
This result generalizes the result of Blandford & Hernquist (1982) for an electron gas to the case of a multi-component system including interacting nucleons. That the functional form of $``$ for interacting nucleons is the same as that for non-interacting particles stems from the fact that, in the mean field approximation, the field equations for the nucleons reduces to the Dirac equation for a free particle, but with an effective mass $`m^{}`$.
### 3.2 Results
In Table 1, we list the various nucleon-meson and meson self-interaction couplings for the two classes of models chosen for this study. In each case, the couplings were chosen to reproduce commonly accepted values of the equilibrium nuclear matter properties: the binding energy per particle $`B/A`$, the saturation density $`n_s`$, the Dirac effective mass $`m_n^{}/m_n`$, the compression modulus $`K_0`$, and the symmetry energy $`a_{\text{sym}}`$. The high-density behavior of the EOS is sensitive to the strength of the meson couplings employed and the models chosen encompass a fairly wide range of variation. The HS81 model, which has a rather high compression modulus, allows us to contrast our results with those of Chakrabarty (1996) who also used HS81 in the case when Landau quantization is considered, and to assess the effects of the inclusion of magnetic moments. Models HS81 and GM1–GM3 employ linear scalar couplings ($`_I`$), while the ZM model employs a nonlinear scalar coupling ($`_{II}`$), which is reflected in the high density behaviors of $`m_n^{}/m_n`$. Thus, comparison of the HS81, GM1–GM3 and ZM models allows us to contrast the effects of the underlying EOS.
In Figure 1, we show results of some physical quantitites of interest for our baseline case, model GM3. At supernuclear densities and in the absence of a magnetic field, the matter pressure is dominated by the baryons principally due to the repulsive nature of the strong interactions. Even up to the central density in a neutron star, the proton fraction remains sufficiently small that the neutrons dominate the total pressure.
The magnitude of the magnetic field $`B`$ required to induce significant changes in the EOS may be estimated in a straightforward manner. In the presence of a magnetic field, the contributions from the protons become significant when only one Landau level is occupied, i.e., when the protons are completely spin polarized. This happens when $`q_pB/\mathrm{}c>(2\pi ^4n_p^2)^{1/3}`$. Therefore, we arrive at the estimate $`B^{}>\text{ }\text{-}\lambda _e^2(2\pi ^4Y_p^2n_b^2)^{1/3}`$ for quantum effects to dominate. The proton fraction $`Y_p=n_p/n_b`$, which depends upon both the density and the magnetic field, typically lies in the range 0.1–0.7. As a result the term in parentheses is of order $`1\text{fm}^2`$ for densities $`n_b>0.1n_s`$, and $`\text{ }\text{-}\lambda _e^21.5\times 10^5\text{fm}^2`$. Thus, the magnetic field necessary to introduce significant contributions from the protons is of order $`B^{}10^5`$, which is well below the proton critical field $`B_c^p=(m_p/m_e)^2B_c^e=1.49\times 10^{20}`$ Gauss (or $`B^{}=3\times 10^6`$) for which protons begin to become relativistic.
The results in Figure 1 were obtained by accounting for all of the allowed Landau levels. Indeed, we notice that the matter pressure $`P_m`$, the effective mass $`m_n^{}`$, and the concentrations $`Y_i=n_i/n_b`$ begin to differ significantly from their field-free values only for $`B^{}>10^5`$. The results in the right panels, shown as a function of $`B^{}`$ for four values of $`u=n_b/n_s`$, show that the density dependence of this threshold value is also qualitatively correct.
The upper left panel shows that there is a substantial decrease in the pressure associated with increasing magnetic fields for $`B^{}>10^5`$. This is also evident from the inset, which clearly shows extensive softening of the EOS. The onset of changes in the pressure as a function of the magnetic field may be more clearly seen in the upper right panel, in which results for representative densities are shown.
The neutron effective mass $`m_n^{}`$ is shown in the lower left panel, and demonstrates the extent to which the scalar field $`\sigma `$ is influenced by the presence of magnetic fields. Note that $`m_n^{}`$ also enters in the calculation of all thermodynamic quantities. Again, it is clear that effects due to magnetic fields do not become significant until $`B^{}>10^5`$.
The lower right panel shows that the composition of neutron-star matter changes significantly at high magnetic fields. The striking feature is the large increase in the proton fraction for $`B^{}>10^5`$. This has two significant effects upon the EOS. First, the protons, which are spin polarized, begin to dominate the contributions to thermodynamics arising from the baryons. This leads to a substantial softening of the EOS (see upper left panel). The second effect stems from the requirement of charge neutrality. Because the leptons provide the only source of negative charge, the lepton fraction rises commensurately with the proton fraction. As a result, the lepton contributions to the pressure and energy density are somewhat increased relative to the field-free case. However, the contributions from the baryons remain dominant.
It is important to note that in order to obtain the total energy density and presssure relevant for neutron star structure, contributions from the electromagnetic field $`\epsilon _f=P_f=B^2/(8\pi )=4.814\times 10^8B^2\mathrm{MeV}\mathrm{fm}^3`$ must be added to the matter energy density $`\epsilon _m`$ and pressure $`P_m`$. This has not always been done in the literature. For $`B^{}>10^5`$, the field contributions can dominate the matter pressure, for the densities of interest, as shown in the upper right panel of Figure 1.
Figure 2 shows the dependance of $`H/B`$ on the field strength $`B^{}`$ and the density $`u`$ for the baseline model GM3. The so-called de Haas-van Alphen oscillations are evident and highlight the multi-component nature of the system. The origin of the increasing complexity in the oscillations may be understood by first inspecting the oscillation period when only a single charged species is present. When the quantity $`(E_f^2m^2)/(2qB)`$ successively approaches integer values, successive Landau orbits begin to get populated resulting in an oscillatory structure in $`H/B`$. The width of these oscillations may be estimated by considering the change in magnetic field $`\mathrm{\Delta }B`$ required to increase $`n_{max}`$ by 1. It is found to be dependent upon both the strength of the magnetic field and the Fermi momenta, and is given by
$$\mathrm{\Delta }B=B_+B_{}=\frac{2qB_+^2}{k_f^2+2qB_+},$$
(36)
where $`B_{}`$ and $`B_+`$ denote the fields at the beginning and end of an oscillation. In the low field limit
$$\mathrm{\Delta }B2q\left(\frac{B_+}{k_f}\right)^2,$$
(37)
and the period goes to zero. At subnuclear densities, where muons are generally absent, charge neutrality forces the Fermi momenta of protons and electrons to be equal and only a single oscillation period exists. However, as the density increases above nuclear densities, the appearance of muons introduces further structure in the oscillations as a result of the superposition arising from each of the three charged species present. Furthermore, with increasing density the Fermi momenta of all particle species increase, which decreases the oscillation periods. The insets in each of the panels clearly show these features.
At large enough magnetic fields, only one Landau level is occupied, and the value of $`H/B`$ saturates. Beyond this point, the fraction of $`H`$ that the magnetization comprises becomes increasingly small. This is demonstrated in the lower right panel of Figure 2, in which the ratio $`H/B`$ approaches unity for both $`B^{}=0`$ and $`B^{}3.3\times 10^6B_c^p`$. However, for $`B^{}10^5`$ there is a noticeable deviation of $`H/B`$ from unity. Nevertheless, in all cases considered, $`|4\pi /B|`$ does not exceed 4%, which does not represent a significant deviation from the case in which the magnetization is neglected.
In Figure 3, we compare results of the matter pressure and effective mass for the models HS81 (employed earlier by Chakrabarty (1996)), GM1, GM2, and ZM with a view towards extracting generic trends induced by strong magnetic fields. Although quantitative differences persist between the models, the qualitative trends of the effects of the field on these EOSs are shared with those of model GM3. Remaining differences are principally due to variations in the underlying stiffnesses, effective masses, and symmetry energies of the individual models.
## 4 EFFECTS INCLUDING ANOMALOUS MAGNETIC MOMENTS
We turn now to the inclusion of the anomalous magnetic moments of the nucleons. Johnson & Lippman (1950) first considered the inclusion of anomalous magnetic moments in the Dirac equation, but their formulation was noncovariant. Here, we employ the covariant form, suggested by Bjorken & Drell (1964), using $`_I`$ and $`_{II}`$ to evaluate the effects of magnetic fields. With the inclusion of the anomalous magnetic moments, the baryon field equations become
$$\left[\stackrel{}{\alpha }\stackrel{}{}\left(\stackrel{}{p}q_b\stackrel{}{A}\right)+\beta m_b^{}\right]\mathrm{\Psi }_b=\left(Eg_{\omega _b}\omega ^0g_{\rho _b}\tau _{3_b}\rho ^0+\kappa _b\frac{i}{2}\gamma _0[\gamma _\mu ,\gamma _\nu ]F^{\mu \nu }\right)\mathrm{\Psi }_b.$$
(38)
The derivation of the Dirac spinors is presented in Appendix A.
The energy spectrum for the protons is given by
$`E_{p,n,s}`$ $`=`$ $`\sqrt{k_z^2+\left(\sqrt{m_p^2+2\left(n+{\displaystyle \frac{1}{2}}+s{\displaystyle \frac{1}{2}}\right)q_pB}+s\kappa _pB\right)^2}+g_{\omega _p}\omega ^0{\displaystyle \frac{1}{2}}g_{\rho _p}\rho ^0`$ (39)
which may be compared with the result
$`E_{p,n,s}^{JL}`$ $`=`$ $`\sqrt{2\left(n+{\displaystyle \frac{1}{2}}+s{\displaystyle \frac{1}{2}}\right)q_pB+\left(\sqrt{m_p^2+k_z^2}+s\kappa _pB\right)^2}+g_{\omega _p}\omega ^0{\displaystyle \frac{1}{2}}g_{\rho _p}\rho ^0`$ (40)
obtained by using the Johnson & Lippman form $`\kappa _b(i/2)[\gamma _\mu ,\gamma _\nu ]F^{\mu \nu }`$ in equation (38) for the inclusion of the magnetic moment. In both cases $`n`$ and $`s`$ are the principle quantum number and “spin” quantum number respectively. As will be shown in the appendix, unlike the $`\kappa _b=0`$ case, the “big” components of the Dirac spinor are no longer eigenstates of the spin operator along the magnetic field ($`\sigma _z`$). However, as $`\kappa _p`$ tends toward zero, it is clear that the proton energy spectum reduces to the expression given in equation (14) and $`s`$ corresponds to the $`\sigma _z`$ eigenvalue. In the nonrelativistic limit, $`k_z^2m^2`$ and $`2\nu q_pBm^2`$ and for $`(\kappa _pB)^2m^2`$, both of the above results reduce to $`E_{p,n,s}m_p^{}+k_z^2/2m_p^{}+\left(n+1/2+s/2\right)(q_pB)/m_p^{}+s\kappa _pB`$, the standard nonrelativistic expression.
The evaluation of the thermodynamic quantities proceeds along the lines already presented in §2, but with a new definition of the Fermi momentum to account for the presence of magnetic moments, which cause an asymmetry in the phase space in addition to that caused by the charged particle interactions with $`B`$. With the energy spectrum in equation (39),
$$k_{f,n,s}^p=\sqrt{E_f^{p2}\left(\sqrt{m_p^2+2\left(n+\frac{1}{2}+\frac{1}{2}s\right)q_pB}+s\kappa _pB\right)^2}.$$
(41)
Setting
$$\overline{m}=\sqrt{m_p^2+2\left(n+\frac{1}{2}+\frac{1}{2}s\right)q_pB}+s\kappa _pB,$$
(42)
and substituting $`k_{f,n,s}`$ and $`\overline{m}`$ into the formulas for the number and energy density for the leptons (see § 2), one finds the analogous quantities for the protons:
$`n_p`$ $`=`$ $`{\displaystyle \frac{|q_p|B}{2\pi ^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{s}{}}k_{f,n,s}^p,`$ (43)
$`\epsilon _p`$ $`=`$ $`{\displaystyle \frac{|q_p|B}{2\pi ^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{s}{}}E_f^pk_{f,n,s}^p+\overline{m}^2\mathrm{ln}\left(\left|{\displaystyle \frac{E_f^p+k_{f,n,s}^p}{\overline{m}}}\right|\right).`$ (44)
The scalar density may be defined in terms of the energy spectrum by
$$n^s=d^3k\frac{E}{m_p^{}}.$$
(45)
Utilizing equation (39) results in
$`n_p^s`$ $`=`$ $`{\displaystyle d^3k\frac{\overline{m}}{\overline{m}s\kappa _pB}\frac{m_p^{}}{E}}`$ (46)
$`=`$ $`{\displaystyle \frac{|q_p|B}{2\pi ^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{s}{}}m_p^{}{\displaystyle \frac{\overline{m}}{\overline{m}s\kappa _pB}}\mathrm{ln}\left(\left|{\displaystyle \frac{E_f^p+k_{f,n,s}^p}{\overline{m}}}\right|\right).`$
The appearance of the factor $`\overline{m}/(\overline{m}s\kappa _pB)`$ may be understood by inspecting the zeroth component of the current four-vector
$$\mathrm{\Psi }^p\mathrm{\Psi }^p=j_p^0=\gamma j_p^0|_{k=0}=\gamma \overline{\mathrm{\Psi }^p}\mathrm{\Psi }^p$$
(47)
with
$$\gamma =\frac{E}{E|_{k=0}}=\frac{E}{m_p^{}+s\kappa _pB^{}},$$
(48)
where $`B^{}`$ is taken in the rest frame of the particle. Hence, the scalar density becomes
$$n_p^s=\overline{\mathrm{\Psi }^p}\mathrm{\Psi }^p=\frac{m_p^{}+s\kappa _pB^{}}{E}\mathrm{\Psi }^p\mathrm{\Psi }^p=d^3k\frac{\overline{m}}{\overline{m}s\kappa _pB}\frac{m_p^{}}{E}.$$
(49)
In a similar way, the energy spectrum of the neutrons is given by
$$E_{n,s}=\sqrt{k_z^2+\left(\sqrt{m_n^2+k_x^2+k_y^2}+s\kappa _nB\right)^2}+g_{\omega _n}\omega ^0+\frac{1}{2}g_{\rho _n}\rho ^0.$$
(50)
Including the magnetic moment for the neutron, the integral over phase space for any thermodynamical quantity $`Q`$ may be easily evaluated by noting that at zero temperature, it is simply the integral over all momenta within the Fermi surface defined by
$$E_f^n=E_{n,s}(k_x,k_y,k_z).$$
(51)
The integral may be written in terms of parallel and perpendicular components,
$$Q=\underset{s}{}\frac{1}{2\pi ^2}_0^bk_{}𝑑k_{}_0^a𝑑k_{}Q,$$
(52)
where $`a`$ and $`b`$ are determined by the Fermi surface to be
$`a`$ $`=`$ $`\sqrt{E_f^{n2}\left(\sqrt{k_{}^2+m_n^2}+s\kappa _nB\right)^2},`$ (53)
$`b`$ $`=`$ $`\sqrt{(E_f^ns\kappa _nB)^2m_n^2}.`$ (54)
With the substitution
$`x=\sqrt{k_{}^2+m_n^2}+s\kappa _nB,`$ (55)
the integral is transformed into
$$Q=\underset{s}{}\left(_{\overline{m}}^{E_f^n}x𝑑x_0^{\sqrt{E_f^{n2}x^2}}𝑑k_{}Q\right)s\kappa _nB\left(_{\overline{m}}^{E_f^n}𝑑x_0^{\sqrt{E_f^{n2}x^2}}𝑑k_{}Q\right),$$
(56)
where
$$\overline{m}=m_n^{}+s\kappa _nB.$$
(57)
Note that the first term is precisely the same as for the $`\kappa _n=0`$ case, but with a shifted mass. This form for the integral over phase space is particularly useful for calculating the number and energy densities. Defining $`k_{f,s}`$ by
$$k_{f,s}=\sqrt{E_f^{n2}\overline{m}^2}$$
(58)
the number and energy densities take the form
$`n_n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{s}{}}{\displaystyle \frac{1}{3}}k_{f,s}^3+{\displaystyle \frac{1}{2}}s\kappa _nB\left[\overline{m}k_{f,s}+E_f^{n2}\left(\mathrm{arcsin}{\displaystyle \frac{\overline{m}}{E_f^n}}{\displaystyle \frac{\pi }{2}}\right)\right],`$ (59)
$`\epsilon _n`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{s}{}}{\displaystyle \frac{1}{2}}E_f^{n3}k_{f,s}+{\displaystyle \frac{2}{3}}s\kappa _nBE_f^{n3}\left(\mathrm{arcsin}{\displaystyle \frac{\overline{m}}{E_f^n}}{\displaystyle \frac{\pi }{2}}\right)`$ (60)
$`+`$ $`\left({\displaystyle \frac{1}{3}}s\kappa _nB{\displaystyle \frac{1}{4}}\overline{m}\right)\left[\overline{m}k_{f,s}E_f^n+\overline{m}^3\mathrm{ln}\left(\left|{\displaystyle \frac{E_f^n+k_{f,s}}{\overline{m}}}\right|\right)\right].`$
The scalar number density reads
$$n_n^s=d^3k\left(1+\frac{s\kappa _nB}{\sqrt{k_{}^2+m_n^2}}\right)\frac{m_n^{}}{E},$$
(61)
which may be recast as
$$n_n^s=\underset{s}{}_{\overline{m}}^{E_f^n}x𝑑x_0^{\sqrt{E_f^{n2}x^2}}𝑑k_{}\frac{m_n^{}}{\sqrt{k_{}^2+x^2}}.$$
(62)
Performing the integration gives
$$n_n^s=\frac{m_n^{}}{4\pi ^2}\underset{s}{}k_{f,s}E_f^n\overline{m}^2\mathrm{ln}\left(\left|\frac{E_f^n+k_{f,s}}{\overline{m}}\right|\right).$$
(63)
As in the case without magnetic moments, the pressure in beta equilibrium is given by $`P=\mu _nn_b\epsilon `$.
### 4.1 Magnetization
Utilizing the expressions for the energy and number densities derived above, the magnetization including the effects of the anomalous magnetic moments may be calculated using the general relation in equation (34). For protons and neutrons, the results are given by
$`_p`$ $`=`$ $`\left\{{\displaystyle \frac{\epsilon _pE_f^pn_p}{B}}+{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{n}{}}{\displaystyle \underset{s}{}}\overline{m}\mathrm{ln}\left(\left|{\displaystyle \frac{E_f^p+k_{f,n,s}^p}{\overline{m}}}\right|\right)\left[{\displaystyle \frac{\left(n+\frac{1}{2}+\frac{1}{2}s\right)}{\overline{m}s\kappa _pB}}+s\kappa _pB\right]\right\},`$ (64)
$`_n`$ $`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{s}{}}\kappa _ns\{({\displaystyle \frac{1}{6}}\overline{m}{\displaystyle \frac{1}{2}}s\kappa _nB)E_f^nk_{f,s}{\displaystyle \frac{1}{6}}E_f^{n3}(\mathrm{arcsin}{\displaystyle \frac{\overline{m}}{E_f^n}}{\displaystyle \frac{\pi }{2}})`$ (65)
$`+({\displaystyle \frac{1}{2}}s\kappa _nB{\displaystyle \frac{1}{3}}\overline{m})\overline{m}^3\mathrm{ln}\left(\right|{\displaystyle \frac{E_f^n+k_{f,s}}{\overline{m}}}\left|\right)\}.`$
The extent to which the anomalous magnetic moments alter the magnetization relative to the case in which they are absent may be gauged by the magnitudes of $`H/B`$ in single component systems. For example, at fields below $`2.2\times 10^{19}`$ Gauss, $`H/B`$ is reduced by approximately 1% in a proton gas and by about 0.3% in a neutron gas.
### 4.2 Results for the $`npe\mu `$ Gas
To assess the influence of the anomalous magnetic moments on the EOS, it is instructive to consider a charge neutral $`npe\mu `$ gas in beta equilibrium. In addition to providing contrasts with the case in which only the effects of Landau quantization are considered (Lai & Shapiro 1991), it sets the stage for the effects to be expected for the case in which baryonic interactions are included.
The magnitude of the magnetic field required to induce significant effects on the EOS due to the inclusion of the magnetic moments may be inferred by considering the field strength at which neutrons become completely polarized. From equation (60), it is clear that complete polarization occurs when $`|\kappa _n|B=k_{f,+1}^2/(4m_n)(6\pi ^2n_n)^{\frac{2}{3}}/(4m_n)`$. At nuclear density, this leads to $`B^{}1.6\times 10^5`$. Note that this is approximately where the effects due to Landau quantization become large. This implies that a complete description of neutron-star matter in the presence of intense magnetic fields must necessarily include the nucleon anomalous magnetic moments.
The equations governing the thermodynamics of the gas are simply the non-interacting limits of equations (3) through (11), and equation (34) for the magnetization. The results are presented in Figure 4 in which the darker (lighter) shade curves show results with (without) the inclusion of the anomalous magnetic moments.
The left panels, in which the matter pressure is shown as functions of $`u`$ and $`\epsilon _m`$, clearly show that the EOS is stiffened upon the inclusion of magnetic moments. For example, in the extreme case when the field strength approaches the proton critical field, the pressure is increased by an order of magnitude over the zero field case (and two orders of magnitude over the case in which only the effects of Landau quantization are considered). The upper right panel, in which the matter pressure is shown as a function of $`B^{}`$, shows that above $`B^{}=10^5`$ the effects of the magnetic moments are more significant than those due to Landau quantization, and cannot be ignored.
The lower right panel provides some insight into the origin of the stiffening. At field strengths of $`B^{}=10^5`$, the composition of matter is dominated by neutrons, the proton fraction being small, about 0.1. Neutrons, however, are spin (up) polarized due to the interaction of the magnetic moment with the magnetic field. With increasing $`B`$, the fraction of neutrons that are polarized increases leading to a corresponding increase in the degeneracy pressure. Upon complete polarization, this increase is halted due to the absence of neutrons needed to fill further spin up energy levels. This is evident from the turnover in the matter pressure, occuring precisely at the point when the neutrons become completely spin-polarized, shown in the upper right panel.
### 4.3 Results for Interacting Matter
In this section, we include the effects of baryonic interactions, Landau quantization, and anomalous magnetic moments. In the absence of magnetic fields, the dominant effect of interactions between the baryons is to substantially stiffen the EOS compared to the case in which interactions are omitted. This is chiefly due to the repulsive nature of the baryonic interactions in beta stable matter. Notwithstanding the fact that the absolute magnitudes of the energy density and pressure are larger than the case in which the baryonic interactions are omitted, magnetic fields have many of the the qualitative effects discussed in the previous section.
The results for the baseline model GM3 are shown in Figure 5, which should be compared with Figure 1 to assess the role of magnetic moments. The upper left panel shows that the stiffening of the EOS observed for the $`npe\mu `$ gas (for $`B^{}>10^5`$) is also present in the case when interactions are included. The effects of magnetic moments are such that the softening caused by Landau quantiziation alone is overwhelmed, leading to an overall stiffening of the EOS. In fact, for fields on the order of the critical proton field, the EOS approaches the causal limit, $`p_m=\epsilon _m`$. As in Figure 1, the matter pressure $`P_m`$, the effective mass $`m_n^{}`$, and the concentrations $`Y_i=n_i/n_b`$ begin to differ significantly from their field-free values only for $`B^{}>10^5`$.
The neutron effective mass $`m_n^{}`$ is shown in the lower left panel. The behavior of $`m_n^{}`$ with $`B^{}`$ is opposite to that shown in Figure 1. The effects of magnetic moments cause $`m_n^{}`$ to increase at a rate approximately equal to $`\kappa _pB/m_n`$ and to become independent of density for $`B^{}>10^6`$. Note that this feature is also a consequence of complete spin polarization.
The lower right panel shows the relative concentrations. Comparing with Figure 1, it is evident that the composition of matter is principally controlled by the effects of Landau quantization. In contrast, the stiffening of the EOS is caused primarily by terms that are explicitly dependent upon the magnetic moments in the pressure and energy density.
Figure 6, to be compared with Figure 2, shows $`H/B`$ as functions of both $`B^{}`$ and $`u`$ for the baseline model GM3. The origin of the oscillations is similar to that discussed in conjunction with Figure 2, but there is an overall reduction of approximately 1% in $`H/B`$ caused chiefly by the magnetization of the neutron.
In Figure 7 (to be compared with Figure 3), we compare results among the models HS81, GM1, GM2, and ZM with the intention of extracting generic trends induced by the inclusion of magnetic moments. The pressure and effective masses share the qualitative trends exhibited by model GM3 (shown in Figure 5), although quantitative differences persist between the models. The stiffness induced by the inclusion of magnetic moments emerges as a general trend, and remaining differences are principally due to variations in the underlying stiffness, effective mass, and symmetry energies of these models..
## 5 SUMMARY AND OUTLOOK
We have developed the methodology necessary to consistently incorporate the effects of magnetic fields on the EOS in multicomponent, interacting matter, including a covariant description for the inclusion of the anomalous magentic moments of nucleons. This methodology is necessary because in the presence of the field all thermodynamic quantities inherit the dimensionful scale set by the magnetic field, which necessarily affects the composition and hence the EOS of matter. By employing a field theoretical-apporach which allows the study of models with different high density behaviors, we found that the results of incorporating strong magnetic fields were not very dependent upon the precise form of the model for the nucleon-nucleon interaction. The generic effects included softening of the EOS due to Landau quantization, which is, however, overwhelmed by stiffening due to the incorporation of the anomalous magnetic moments of the nucleons. These effects become significant for fields in excess of $`B^{}10^5`$, for which neutrons become completely spin polarized. Note that this field strength is substantially less than the proton critical field. In addition, the inclusion of ultra-strong magnetic fields leads to a reduction in the electron chemical potential and an increase in proton fraction. These compositional changes have implications for neutrino emission via the direct Urca process and, thus, for the cooling of neutron stars. The magnetization of the matter never appears to become very large, as the value of $`|H/B|`$ never deviates from unity by more than a few percent. However, it remains to be seen what effects the magnetization of matter will have on the structure and transport properties of neutron stars.
It is worthwhile to note here that the qualitative effects of strong magnetic fields found in the relativistic field-theoretical description of dense matter would also be found in non-relativistic potential models. This is because the phase space of charged particles is similarly affected in both approaches by the presence of magnetic fields. The effects due to the anomalous magnetic moments would, however, enter linearly in a non-relativistic approach (see §4), and would thus be more dramatic in this case. It would be also be instructive to study the effects of magnetic fields including many-body correlations.
It would be useful to also consider cases in which strangeness-bearing hyperons, a Bose (pion or kaon) condensate or quarks, are present in dense matter. The covariant description of the anomalous magnetic moments developed in this work may be utilized to include hyperons, which are likely to be present in dense matter (Glendenning 1982, 1985; Weber & Weigel 1985; Kapusta & Olive 1990; Ellis, Kapusta & Olive 1991; Glendenning & Moszkowski 1991; Sumiyoshi & Toki 1994; Prakash et al. 1997 and references therein). The anomalous magnetic moments of hyperons are mostly known. The negatively charged hyperons, the neutral $`\mathrm{\Lambda }`$, and $`\mathrm{\Xi }^0`$ all have negative anomalous magnetic moments. $`\mathrm{\Sigma }^+`$ and $`\mathrm{\Sigma }^0`$ are the only hyperons with positive anomalous magnetic moments. The effects of Landau quantization on hyperons would be to soften the EOS relative to the case in which magnetic fields are absent. However, in the presence of strong magnetic fields, all of the hyperons will be spin polarized due to magnetic moment interactions with the field. This would cause their degeneracy pressures to increase compared to the field-free case. The resultant of these two opposing effects will depend on the relative concentrations of the various hyperons, which in turn depends sensitively on the hyperon-meson interactions for which only a modest amount of guidance is available (Glendenning & Moszkowski 1991, Knorren, Prakash & Ellis 1995, Schaffner & Mishustin 1996). For choices of $`\mathrm{\Sigma }^{}`$meson interactions that favor the appearence of $`\mathrm{\Sigma }^{}`$ hyperons at relatively low densities, the concentrations of the positively charged particles, $`p`$ and $`\mathrm{\Sigma }^+`$, may be expected to increase in the presence of strong magnetic fields. It would thus appear that the effects of including hyperons will not drastically alter the qualitative trends of increasing the concentrations of positively charged particles found in the case of $`npe\mu `$ matter. The main physical effects found in the absence of hyperons, namely increasing the stiffness of matter, and allowing the direct Urca process (Lattimer et al. 1991; Prakash et al. 1992) to occur, probably would not change, either. Feedback effects due to mass and energy shifts may, however, alter these expectations. Thus, detailed calculations are required to ascertain the influence of magnetic fields in multi-component matter. Work on this topic is currently in progress and will be reported separately.
It is intriguing that Bosons (pions and kaons), which have zero magnetic moment, do not feel the magnetic fields as fermions do. Similarly, quarks without sub-structure also have no anomalous magnetic moments. Thus, intense magnetic fields in the cores of stars containing a Bose condensate or quark matter might serve as a useful discriminant compared to those containing baryonic matter.
Work is in progress (Cardall et al. 1999) to complete a fully self-consistent calculation of neutron star structure including the combined effects of the direct effects of magnetic fields on the EOS, which we have developed in this paper, and general relativistic structure. The findings will help answer questions concerning the largest frozen-in magnetic field that a stationary neutron star can possess, and what the structure of stars with ultra-strong fields might be. It must be borne in mind, however, that for super-strong fields (much higher than $`B_c^p`$, which is the highest field considered in this work), the energy density in the field would be significantly higher than the baryon mass energy density. Under such conditions, the internal structure of the baryons will be affected and alternative descriptions for the EOS will become necessary.
We thank Hans Hansson for constructive suggestions concerning the covariant description of the anomalous magnetic moments. This work was supported in part by the NASA ATP Grant # NAG 52863, and by the USDOE grants DOE/DE-FG02-87ER-40317 & DOE/DE-FG02-88ER-40388.
## Appendix A SPINORS AND ENERGY SPECTRA FOR BARYONS WITH ANOMALOUS MAGNETIC MOMENTS
In this appendix, we derive relations for the spinors and energy spectra for baryons with anomalous magnetic moments. The Dirac equation is
$$\left[\stackrel{}{\alpha }\stackrel{}{}\left(\stackrel{}{p}q_b\stackrel{}{A}\right)+\beta m_b^{}+\beta \sigma _z\kappa _bB\right]\mathrm{\Psi }_b=E_{0,b}\mathrm{\Psi }_b,$$
(A1)
where the effective momentum is given by $`\stackrel{}{\pi }=\stackrel{}{p}q_b\stackrel{}{A}`$ and $`\kappa _b`$ denotes the baryon anomalous magnetic moment. The energy $`E_{0,b}`$ denotes the baryon energy eigenvalues when the meson fields are absent and are related to the neutron and proton energy spectra given in equations (39) and (50) by
$`E_{n,s}`$ $`=`$ $`E_{0,n}+g_{\omega _b}\omega ^0+{\displaystyle \frac{1}{2}}g_{\rho _b}\rho ^0`$ (A2)
$`E_{p,n,s}`$ $`=`$ $`E_{0,p}+g_{\omega _b}\omega ^0{\displaystyle \frac{1}{2}}g_{\rho _b}\rho ^0,`$ (A3)
repectively. Separating $`\mathrm{\Psi }_b`$ in to “big” and “small” components, we obtain
$`\left(E_{0,b}m_b^{}\kappa _bB\sigma _z\right)\varphi `$ $`=`$ $`\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\chi `$ (A4)
$`\left(E_{0,b}+m_b^{}+\kappa _bB\sigma _z\right)\chi `$ $`=`$ $`\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\varphi .`$ (A5)
Writing $`\chi `$ in terms of $`\varphi `$ (taking care to note that the terms on the left hand side of these equations are no longer proportional to the identity matrix because of the presence of the magnetic moments), equation (A4) becomes
$$\left(E_{0,b}m_b^{}\kappa _bB\sigma _z\right)\varphi =\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\frac{E_{0,b}+m_b^{}\kappa _bB\sigma _z}{\left(E_{0,b}+m_b^{}\right)^2\left(\kappa _bB\right)^2}\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\varphi .$$
(A6)
Note that the term with $`\sigma _z`$ does not commute with the momentum operators. Therefore,
$$\left(E_{0,b}m_b^{}\kappa _bB\sigma _z\right)\varphi =\frac{E_{0,b}+m_b^{}+\kappa _bB\sigma _z}{\left(E_{0,b}+m_b^{}\right)^2\left(\kappa _bB\right)^2}\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)^2\varphi \frac{2\kappa _bB\sigma _z}{\left(E_{0,b}+m_b^{}\right)^2\left(\kappa _bB\right)^2}\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\pi _z\varphi .$$
(A7)
This may be rewritten as
$$F_s\varphi =\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)^2\varphi a_s\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{\pi }\right)\pi _z\varphi ,$$
(A8)
where $`F_s`$ and $`a_s`$ are defined as
$`F_s`$ $`=`$ $`\left(E_{0,b}\kappa _bB\sigma _z\right)^2m_b^2`$
$`a_s`$ $`=`$ $`{\displaystyle \frac{2\kappa _bB\left(E_{0,b}+m_b^{}\kappa _bB\sigma _z\right)}{\left(E_{0,b}+m_b^{}\right)^2\left(\kappa _bB\right)^2}}.`$ (A9)
At this point it is necessary to consider individually the cases of the protons and neutrons.
#### Protons
The fact that $`[\pi _x,\pi _y]=i\mathrm{}(qB/c)`$ suggests the transformations
$$p_\xi =\sqrt{\frac{c}{q_pB}}\pi _x,\xi =\sqrt{\frac{c}{q_pB}}\pi _y.$$
(A10)
Using the identities
$$\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{a}\right)\left(\stackrel{}{\sigma }\stackrel{}{}\stackrel{}{b}\right)=\stackrel{}{a}\stackrel{}{}\stackrel{}{b}+i\stackrel{}{\sigma }\stackrel{}{}\left(\stackrel{}{a}\times \stackrel{}{b}\right),\stackrel{}{\pi }\times \stackrel{}{\pi }=iq_pB\widehat{z}$$
(A11)
and the above transformations, equation (A8) becomes
$$F_s\varphi =\left[q_pB\left(p_\xi ^2+\xi ^2\sigma _z\right)+p_z^2\right]\varphi a_s\left[\sqrt{\frac{c}{q_pB}}\left(\sigma _xp_xi\sigma _y\xi \right)+\sigma _zp_z\right]p_z\varphi .$$
(A12)
The similarities between equation (A12) and that for the leptons (see, for example, Itzykson & Zuber 1984), suggests the ansatz for the spin up spinor
$$\varphi _{+1}=e^{ik_zZ\frac{\xi ^2}{2}}\left(\begin{array}{c}H_n(\xi )\\ i\omega _{p,n,+1}H_{n+1}(\xi )\end{array}\right).$$
(A13)
The two coupled differential equations for the components of $`\varphi `$ (equations (A12)) reduce to two coupled algebraic equations for the eigenvalues of $`E_{0,p}`$ and $`\omega _{p,n,+1}`$. Explicitly,
$`F_{+1}`$ $`=`$ $`2(n+1)q_pB+(1a_{+1})k_z^2a_{+1}2(n+1)\sqrt{q_pB}k_z\omega _{p,n,+1}`$
$`F_1`$ $`=`$ $`2(n+1)q_pB+(1+a_1)k_z^2a_1k_z\sqrt{q_pB}\omega _{p,n,+1}^1.`$ (A14)
These may be solved to give
$$E_{0,p,+1}=\sqrt{k_z^2+\left(\sqrt{m_b^2+2(n+1)q_pB}+\kappa _bB\right)^2}.$$
(A15)
With this result, it is straightforward to solve for $`\omega _{p,n,+1}`$. Lacking a simple expression, we shall continue to refer to it as $`\omega _{p,n,+1}`$. Substituting this solution for $`\varphi _{+1}`$ into the expression for $`\chi `$ gives the Dirac spinor
$$\mathrm{\Psi }_{n,+1}^p=Ne^{ik_zz\frac{\xi ^2}{2}}\left(\begin{array}{c}H_n(\xi )\\ i\omega _{p,n,+1}H_{n+1}(\xi )\\ \frac{k_z+2(n+1)\omega _{p,n,+1}\sqrt{q_pB}}{E_{0,p,+1}+m_p^{}+\kappa _pB}H_{n1}(\xi )\\ \frac{i\sqrt{q_pB}+i\omega _{p,n,+1}k_z}{E_{0,p,+1}+m_p^{}\kappa _pB}H_n(\xi )\end{array}\right).$$
(A16)
A similar method may be employed to find an ansatz for the spin down spinor,
$$\varphi _1=e^{ik_zZ\frac{\xi ^2}{2}}\left(\begin{array}{c}i\omega _{p,n,1}H_{n1}(\xi )\\ H_n(\xi )\end{array}\right),$$
(A17)
with the energy eigenvalue
$$E_{0,p,1}=\sqrt{k_z^2+\left(\sqrt{m_b^2+2nq_pB}\kappa _bB\right)^2},$$
(A18)
and the Dirac spinor
$$\mathrm{\Psi }_{n,1}^p=Ne^{ik_zz\frac{\xi ^2}{2}}\left(\begin{array}{c}i\omega _{p,n,1}H_{n1}(\xi )\\ H_n(\xi )\\ \frac{2ni\sqrt{q_pB}i\omega _{p,n,1}k_z}{E_{0,p,1}+m_p^{}+\kappa _pB}H_{n1}(\xi )\\ \frac{k_z\omega _{p,n,1}\sqrt{q_pB}}{E_{0,p,1}+m_p^{}\kappa _pB}H_n(\xi )\end{array}\right).$$
(A19)
While all quantities in this work have been calculated in the zero temperature approximation, requiring only the postive energy spinors, for completeness the negative energy Dirac spinors are presented below. For the protons these may be determined in much the same manner as that employed for the positive energy spinors. Defining
$`F_s^{}`$ $`=`$ $`\left(E_{0,p}+\kappa _pB\sigma _z\right)^2m_p^2`$ (A20)
$`a_s^{}`$ $`=`$ $`{\displaystyle \frac{2\kappa _pB\left(E_{0,p}m_p^{}+\kappa _pB\sigma _z\right)}{\left(E_{0,p}m_p^{}\right)^2\left(\kappa _pB\right)^2}},`$ (A21)
the equation for $`\chi `$ takes the same form as equation (A8) where $`F_s^{}`$ and $`a_s^{}`$ replace $`F_s`$ and $`a_s`$ respectively. As a result, precisely the same formalism employed to determine the positive energy spinors may be used to determine the negative energy spinors. The Dirac spinor corresponding to the energy eigenvalue
$$E_{0,p,+1}^{}=\sqrt{k_z^2+\left(\sqrt{m_p^2+2\left(n+1\right)q_pB}+\kappa _pB\right)^2},$$
(A22)
is given by
$$\mathrm{\Psi }_{n,+1}^{p,}=Ne^{ik_zZ\frac{\xi ^2}{2}}\left(\begin{array}{c}\frac{k_z+2(n+1)\omega _{p,n,+1}^{}\sqrt{q_pB}}{E_{0,p,+1}^{}m_p^{}\kappa _pB}H_n(\xi )\\ \frac{i\sqrt{q_pB}+i\omega _{p,n,+1}^{}k_z}{E_{0,p,+1}^{}m_p^{}+\kappa _pB}H_{n+1}(\xi )\\ H_n(\xi )\\ i\omega _{p,n,+1}^{}H_{n+1}(\xi )\end{array}\right),$$
(A23)
where $`\omega _{p,n,+1}^{}`$ is defined by replacing $`F_s`$ and $`a_s`$ in equations (A14). Similarly, the Dirac spinor corresponding to the energy eigenvalue
$$E_{0,p,1}^{}=\sqrt{k_z^2+\left(\sqrt{m_p^2+2nq_pB}\kappa _pB\right)^2},$$
(A24)
is given by
$$\mathrm{\Psi }_{n,1}^{p,}=Ne^{ik_zZ\frac{\xi ^2}{2}}\left(\begin{array}{c}\frac{2in\sqrt{q_pB}i\omega _{p,n,1}^{}k_z}{E_{0,p,+1}^{}m_p^{}\kappa _pB}H_{n1}(\xi )\\ \frac{k_z\omega _{p,n,1}^{}\sqrt{q_pB}}{E_{0,p,+1}^{}m_p^{}+\kappa _pB}H_n(\xi )\\ i\omega _{p,n,1}^{}H_{n1}(\xi )\\ H_n(\xi )\end{array}\right).$$
(A25)
#### Neutrons
In this case, the trial wave function has the same form as the free particle solutions with unknown coefficients, which may be determined in a manner analougous to that employed for the protons. Define
$$G_s=F_sk^2+\sigma _za_sk_z^2.$$
(A26)
Then, equation (A8) becomes
$$G_s\varphi =a_s(\sigma _xk_x+\sigma _yk_y)k_z\varphi .$$
(A27)
Note that both $`G_s`$ and $`a_s`$ are diagonal and therefore the off-diagonal terms have been isolated on the right-hand side of equation (A27). The similarities with the case in which $`\kappa _n=0`$, namely the quadratic nature of the momentum operators, suggests the form
$$\varphi =e^{ik^\mu x_\mu }\left(\begin{array}{c}u\\ v\end{array}\right).$$
(A28)
Using equation (A28), we obtain the coupled algebraic equations
$`G_{+1}u`$ $`=`$ $`a_{+1}(k_xik_y)k_zv`$
$`G_1v`$ $`=`$ $`a_1(k_x+ik_y)k_zu.`$ (A29)
Combining these gives
$$G_{+1}G_1=a_{+1}a_1(k_x^2+k_y^2)k_z^2,$$
(A30)
which may be solved for the energy eigenvalue
$$E_{0,n,s}=\sqrt{k_z^2+\left(\sqrt{m_n^2+k_x^2+k_y^2}+s\kappa _nB\right)^2}.$$
(A31)
The eigenvectors may be determined, up to a normalization, by setting
$`u`$ $`=`$ $`1v={\displaystyle \frac{a_1}{G_1}}(k_x+ik_y)k_z`$ (A32)
$`v`$ $`=`$ $`1u={\displaystyle \frac{a_{+1}}{G_{+1}}}(k_x+ik_y)k_z`$ (A33)
in equation (A29). It is clear from direct substitution that the first gives the $`s=+1`$ and the second the $`s=1`$ spinors. Inserting these into equation (A28) and then into the expression for $`\chi `$ gives the neutron Dirac spinors
$$\mathrm{\Psi }_{+1}^n=Ne^{ik^\mu x_\mu }\left(\begin{array}{c}1\\ \frac{a_1}{G_1}\left(k_x+ik_y\right)k_z\\ \frac{\left[1\frac{a_1}{G_1}\left(k_x^2+k_y^2\right)\right]}{E_{0,n,+1}+M_n^{}+\kappa _nB}\\ \frac{\left[1\frac{a_{}}{b_{}}k_z^2\right]\left(k_x+ik_y\right)}{E_{0,n,+1}+M_n^{}\kappa _nB},\end{array}\right)$$
(A34)
$$\mathrm{\Psi }_1^n=Ne^{ik^\mu x_\mu }\left(\begin{array}{c}\frac{a_{+1}}{G_{+1}}\left(k_xik_y\right)k_z\\ 1\\ \frac{\left[1\frac{a_{+1}}{G_{+1}}k_z^2\right]\left(k_xik_y\right)}{E_{0,n,1}+M_n^{}+\kappa _nB}\\ \frac{\left[1+\frac{a_{+1}}{G_{+1}}\left(k_x^2+k_y^2\right)\right]k_x}{E_{0,n,1}+M_n^{}\kappa _nB}\end{array}\right).$$
(A35)
In order to determine the negative energy Dirac spinors for the neutron, an approach analogous to that employed in determining the negative energy Dirac spinors for the protons may be used. Define $`G_s^{}`$ by replacing $`F_s`$ and $`a_s`$ by $`F_s^{}`$ and $`a_s^{}`$, respectively, in equation (A26). As in the case of the protons, this produces an equation for $`\chi `$ which is of the same form as that employed for $`\varphi `$ in the derivation of the positive energy spinors. Proceeding in the same manner as before, one finds that the Dirac spinors corresponding to the energy eigenvalues
$$E_{0,n,s}^{}=\sqrt{k_z^2+\left(\sqrt{m_n^{}+k_x^2+k_y^2}+s\kappa _nB\right)^2},$$
(A36)
are given by
$`\mathrm{\Psi }_{+1}^{n,}`$ $`=`$ $`Ne^{ik^\mu x_\mu }\left(\begin{array}{c}\frac{\left[1\frac{a_1^{}}{G_1^{}}\left(k_x^2+k_y^2\right)\right]k_z}{E_{0,n,+1}m_n^{}\kappa _nB}\\ \frac{\left[1\frac{a_1^{}}{G_1^{}}k_z^2\right]\left(k_x+ik_y\right)}{E_{0,n,+1}m_n^{}+\kappa _nB}\\ 1\\ \frac{a_1^{}}{G_1^{}}\left(k_x+ik_y\right)k_z\end{array}\right)`$ (A41)
$`\mathrm{\Psi }_1^{n,}`$ $`=`$ $`Ne^{ik^\mu x_\mu }\left(\begin{array}{c}\frac{\left[1\frac{a_{+1}^{}}{G_{+1}^{}}k_z^2\right]\left(k_xik_y\right)}{E_{0,n,1}m_n^{}\kappa _nB}\\ \frac{\left[1+\frac{a_{+1}^{}}{G_{+1}^{}}\left(k_x^2+k_y^2\right)\right]k_z}{E_{0,n,1}m_n^{}+\kappa _nB}\\ \frac{a_{+1}^{}}{G_{+1}^{}}\left(k_xik_y\right)k_z\\ 1\end{array}\right).`$ (A46)
TABLE 1
NUCLEON-MESON COUPLING CONSTANTS
| Model | $`n_s`$ | $`B/A`$ | $`M^{}/M`$ | $`K_0`$ | $`a_{\text{sym}}`$ | $`g_{\sigma _N}/m_\sigma `$ | $`g_{\omega _N}/m_\omega `$ | $`g_{\rho _N}/m_\rho `$ | $`b`$ | $`c`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| HS81 | 0.148 | 15.75 | 0.54 | 545 | 35.0 | 3.974 | 3.477 | 2.069 | 0.0 | 0.0 |
| GM1 | 0.153 | 16.30 | 0.70 | 300 | 32.5 | 3.434 | 2.674 | 2.100 | 0.002947 | $`0.001070`$ |
| GM2 | 0.153 | 16.30 | 0.78 | 300 | 32.5 | 3.025 | 2.195 | 2.189 | 0.003478 | 0.01328 |
| GM3 | 0.153 | 16.30 | 0.78 | 240 | 32.5 | 3.151 | 2.195 | 2.189 | 0.008659 | $`0.002421`$ |
| ZM | 0.160 | 16.00 | 0.86 | 225 | 32.5 | 2.736 | 1.617 | 2.185 | 0.0 | 0.0 |
NOTE.– Coupling constants for the HS (Horowitz & Serot 1981), GM1-3 (Glendenning & Moszkowski 1991), and ZM (Zimanyi & Moszkowski 1990; 1992) models. The couplings are chosen to reproduce the binding energy $`B/A`$ (MeV), the nuclear saturation density $`n_s`$ ($`\text{fm}^3`$), the Dirac effective mass $`M^{}`$ in units of the baryon mass $`M`$, and the the symmetry energy $`a_{\text{sym}}(MeV)`$. The nuclear matter compression modulus $`K_0`$ (MeV) for the different models are also listed.
## FIGURE CAPTIONS
FIG. 1.– Matter pressure $`P_m`$, nucleon Dirac effective mass $`m_n^{}/m_n`$, and concentrations $`Y_i=n_i/n_b`$ as functions of the density $`u=n_b/n_s`$ (left panels; $`n_s=0.16\mathrm{fm}^3`$ is the fiducial nuclear saturation density) and magnetic field strength $`B^{}=B/B_e^c`$ (right panels; $`B_e^c=4.414\times 10^{13}`$ Gauss is the electron critical field), for the model GM3. The inset in the upper left panel shows $`P_m`$ as a function of the matter energy density $`\epsilon _m`$. The curve labeled $`P_f`$ in the upper right panel shows the $`B^2/8\pi `$ contribution to the total pressure. The inset in the lower left panel shows the effective mass as a function of $`B^{}`$. In the lower right panel, the electron and neutron concentrations have been suppressed for clarity ($`Y_e=Y_pY_\mu `$ and $`Y_n=1Y_p`$).
FIG. 2.– The ratio of the induced to applied magnetic field $`H/B`$ as functions of the density and magnetic field strength, for the model GM3. The insets show $`H/B`$ in expanded scales to highlight the effects of including several components.
FIG. 3.– Matter pressure $`P_m`$ and the nucleon Dirac effective mass $`m_n^{}/m_n`$ for the models shown in Table 1 (with the exception of model GM3, whose results are displayed in Figures 1 and 2), as functions of the density and magnetic field strength. The insets in the left panels show $`P_m`$ as a function of the matter energy density $`\epsilon _m`$.
FIG. 4.– Matter pressure $`P_m`$ and concentrations $`Y_i=n_i/n_b`$ as functions of the density and magnetic field strength for a charge neutral, beta-equilibrated, non-interacting $`npe\mu `$ gas with and without the inclusion of the nucleon anomalous magnetic moments $`\kappa _b`$. The curve labeled $`P_f`$ in the upper right panel shows the $`B^2/8\pi `$ contribution to the total pressure. The lower left panel shows the enhancement in the pressure, as a function of energy density, due to the presence of magnetic fields. In the lower right panel, the electron and neutron concentrations have been suppressed for clarity ($`Y_e=Y_pY_\mu `$ and $`Y_n=1Y_p`$).
FIG. 5.– Same as Figure 1, except that the nucleon anomalous magnetic moments are now included.
FIG. 6.– Same as Figure 2, except that the nucleon anomalous magnetic moments are now included.
FIG. 7.– Same as Figure 3, except that the nucleon anomalous magnetic moments are now included. |
warning/0001/astro-ph0001003.html | ar5iv | text | # Dark Halo and Disk Galaxy Scaling Laws in Hierarchical Universes
## 1. Introduction
The structural parameters of dark matter halos formed in hierarchically clustering universes are tightly related through simple scaling laws that reflect the cosmological context of their formation. These correlations result from the approximately scalefree process of assembly of collisionless dark matter into collapsed, virialized systems. Analytical studies and cosmological N-body simulations have been particularly successful at unraveling the relations between halo mass, size, and angular momentum, as well as the dependence of these correlations on the cosmological parameters. The picture that emerges is encouraging in its simplicity and in its potential applicability to the origin of scaling laws relating the structural properties of galaxy systems (see, e.g., Dalcanton, Spergel & Summers 1997, Mo, Mao & White 1998 and references therein).
One example is the relation between halo mass and size–a direct result of the finite age of the universe. The nature of this correlation and its dependence on the cosmological parameters is straightforward to compute using simple spherical collapse models of “top-hat” density perturbations that have been found to be in good agreement with the results of numerical experiments (see, e.g., Cole & Lacey 1996, Eke, Cole & Frenk 1996, Eke, Navarro & Frenk 1998, and references therein).
A second example concerns the angular momentum of dark halos, which is also linked to halo mass and size through simple scaling arguments. Expressed in nondimensional form, the angular momentum of dark matter halos is approximately independent of mass, environment, and cosmological parameters, a remarkable result likely due to the scalefree properties of the early tidal torques between neighboring systems responsible for the spin of individual halos (Peebles 1969, White 1984, Barnes & Efstathiou 1988, Steinmetz & Bartelmann 1995).
Finally, similarities in the halo formation process are also apparent in the internal structure of dark halos. A number of numerical studies have consistently shown that the shape of the density profiles of dark halos is approximately “universal”; i.e., it can be well approximated by a simple two-parameter function whose formulation is approximately independent of mass, redshift, and cosmology (Navarro, Frenk & White 1996, 1997, hereafter NFW96 and NFW97, respectively; Cole & Lacey 1996; Tormen, Bouchet & White 1996, Huss, Jain & Steinmetz 1999a,b).
How do these scaling properties relate to analogous correlations between structural parameters of disk galaxy systems? This paper is third in a series where this question is addressed through direct numerical simulation of galaxy formation in cold dark matter (CDM) dominated universes. In spirit, these studies are similar to those of Evrard, Summers & Davis 1994, Tissera, Lambas & Abadi 1997, Elizondo et al. 1999, although are conclusions differ in detail from those reached by them. The first paper in our series (Steinmetz & Navarro 1999, hereafter SN1) examined the origin of the Tully-Fisher relation under the assumption that star formation is dictated by the rate at which gas cools and collapses within dark halos. We were able to show that the velocity scaling of luminosity and angular momentum in spiral galaxies arise naturally in hierarchical galaxy formation models.
Large discrepancies, however, were observed in the zero point of these correlations, a result we ascribed to the early dissipative collapse of gas into the progenitor dark matter halos and to the subsequent assembly of the final system through a sequence of mergers (Navarro & Benz 1991, Navarro, Frenk & White 1995, Navarro & Steinmetz 1997, hereafter NS97). We concluded then, in agreement with a number of previous studies, that agreement between models and observations requires a large injection of energy (presumably “feedback” energy from evolving stars and supernovae) in order to prevent much of the gas from cooling and condensing into (proto)galaxies at early times, shifting the bulk of star formation to later times and alleviating the angular momentum losses associated with major mergers (White & Rees 1978, White & Frenk 1991, Kauffmann, White & Guiderdoni 1993, Cole et al. 1994).
This hypothesis remains to date the most attractive path to resolve the “disk angular momentum problem”. On the other hand, further analysis has shown that the zero-point offset between the model and observed Tully-Fisher relations is actually due to the high central concentrations of dark matter halos formed in currently popular versions of the Cold Dark Matter cosmogony, and therefore is unlikely to be reconciled through feedback effects (Navarro & Steinmetz 2000, hereafter NS2). We investigate these issues in detail in the present paper, which represents our first attempt at implementing a realistic numerical formulation of feedback within a full-fledged simulation of the formation of galaxies in CDM halos, and at gauging its effects on the origin of disk galaxy scaling laws.
The paper is organized as follows. Section 2 motivates our approach by reviewing the scaling laws linking the properties of dark matter halos and by comparing them with observational correlations. A brief description of the numerical procedure is included in §3; full details of the star formation, feedback prescription, and relevant tests will be presented in Steinmetz & Navarro (2000, hereafter SN3). Section 4 discusses our main results regarding the slope, scatter and zero-point of the Tully-Fisher relation and of the relation between disk angular momentum and rotation speed. We summarize and discuss our main results in §5, and conclude in §6.
## 2. Dark Halo and Disk Galaxy Scaling Laws
### 2.1. Mass, Radius, and Circular Velocity
The radial distribution of mass in dark halos has no obvious edge, so defining halo “sizes” is a somewhat arbitrary procedure. A plausible and useful choice is to associate the radius of a halo with the distance from the center at which mass shells are infalling for the first time. This “virial radius” is easily estimated from N-body simulations and imposes, by construction, a firm upper limit to the mass of the galaxy embedded inside each halo: baryons beyond the virial radius have not had time yet to reach the center of the halo and therefore cannot have been accreted by the central galaxy.
Numerical experiments show that virial radii depend sensitively on the enclosed mass of the system, in a way consistent with the predictions of the spherical top-hat collapse model. The “edge” of a halo occurs approximately where the mean inner density contrast, $`\mathrm{\Delta }`$, is of order a few hundred.<sup>1</sup><sup>1</sup>1We use the term “density contrast” to refer to densities expressed in units of the critical density for closure, $`\rho _{crit}=3H(z)^2/8\pi G`$, where $`H(z)`$ is the value of Hubble’s constant at redshift $`z`$. We parameterize the present value of $`H`$ as $`H_0=H(z=0)=100h`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. This result seems to apply equally well to halos of all masses, and depends only weakly on cosmology through the density parameter, $`\mathrm{\Omega }`$, and the cosmological constant, $`\mathrm{\Lambda }`$, (Eke et al. 1996, 1998),
$$\mathrm{\Delta }(\mathrm{\Omega },\mathrm{\Lambda })178\{\begin{array}{cc}\mathrm{\Omega }^{0.30},\hfill & \text{if }\mathrm{\Lambda }=0\text{;}\hfill \\ \mathrm{\Omega }^{0.45},\hfill & \text{if }\mathrm{\Omega }+\mathrm{\Lambda }=1\text{.}\hfill \end{array}$$
$`(1)`$
Fig. 1.—The I-band Tully-Fisher relation compared with the results of the numerical simulations. Dots correspond to the observational samples of Mathewson, Ford & Buchhorn, (1992), Giovanelli et al (1997), Willick et al (1997), and Han & Mould (1992). Error bars in the simulated magnitudes correspond to adopting a Salpeter or a Scalo IMF.
One consequence of this definition is that all halos identified at a given time have similar densities, from where it follows on dimensional grounds that the total mass of the halo and its circular velocity at the virial radius must be tightly related. In terms of the circular velocity, $`V_\mathrm{\Delta }`$, at the virial radius, $`r_\mathrm{\Delta }`$, halo masses are given by
$$M_\mathrm{\Delta }(V_\mathrm{\Delta },z)=\left(\frac{2}{\mathrm{\Delta }}\right)^{1/2}\frac{V_\mathrm{\Delta }^3}{GH(z)}=1.9\times 10^{12}\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}$$
$$\times \frac{H_0}{H(z)}\left(\frac{V_\mathrm{\Delta }}{200\mathrm{km}\mathrm{s}^1}\right)^3h^1M_{}.$$
$`(2)`$
This power-law dependence on velocity is similar to that of the I-band Tully-Fisher relation linking the luminosity and rotation speed of late-type spirals,
$$L_I2.0\times 10^{10}\left(\frac{V_{\mathrm{rot}}}{200\mathrm{km}\mathrm{s}^1}\right)^3h^2L_{},$$
$`(3)`$
(see solid line in Figure 1) a coincidence that suggests a direct cosmological origin for this scaling law. Eqs. 2 and 3 can be combined to read,
$$L_I=1.9\times 10^{12}\frac{1}{\mathrm{{\rm Y}}_I}\frac{M_{\mathrm{disk}}}{M_\mathrm{\Delta }}\frac{H_0}{H(z)}\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}\left(\frac{V_\mathrm{\Delta }}{V_{\mathrm{rot}}}\right)^3$$
$$\times \left(\frac{V_{\mathrm{rot}}}{200\mathrm{km}\mathrm{s}^1}\right)^3h^1L_{},$$
$`(4)`$
where we have introduced the parameter $`M_{\mathrm{disk}}`$ to represent the total mass associated with the galaxy disk and $`\mathrm{{\rm Y}}_I=M_{\mathrm{disk}}/L_I`$ is the disk mass-to-light ratio in solar units<sup>2</sup><sup>2</sup>2Note that this definition of the disk mass-to-light ratio includes stellar and gaseous mass. However, most of the baryons in the Tully-Fisher disks we consider here are actually in stars (this is true for observed as well as for simulated galaxies) so that the stellar and “baryonic” disk mass-to-light ratio differ very little. We shall not discriminate between them throughout this paper. Further, we assume $`M_I()=4.15`$ for all numerical values quoted in this paper.
Combining eqs. 2, 3, and 4, we find that the fraction of the total mass of the system in the galaxy disk is, at $`z=0`$,
$$f_{\mathrm{mdsk}}=\frac{M_{\mathrm{disk}}}{M_\mathrm{\Delta }}=8.5\times 10^3h^1\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}\mathrm{{\rm Y}}_I\left(\frac{V_{\mathrm{rot}}}{V_\mathrm{\Delta }}\right)^3.$$
$`(5)`$
This result reemphasizes our implicit assumption that within the “virial radius” galaxy systems are dominated by dark matter. Further insight can be gained by comparing $`M_{\mathrm{disk}}`$ to the total baryonic mass within $`r_\mathrm{\Delta }`$, $`M_{\mathrm{disk}}^{\mathrm{max}}=(\mathrm{\Omega }_b/\mathrm{\Omega }_0)M_\mathrm{\Delta }`$. Assuming that the baryon density parameter is $`\mathrm{\Omega }_b0.0125h^2`$, as suggested by Big Bang nucleosynthesis studies of the primordial abundance of the light elements (Schramm & Turner 1997), the fraction of baryons transformed into stars in disk galaxies is given by,
$$f_{\mathrm{bdsk}}=\frac{M_{\mathrm{disk}}}{M_{\mathrm{disk}}^{\mathrm{max}}}0.85\mathrm{\Omega }_0h\mathrm{{\rm Y}}_I\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}\left(\frac{V_{\mathrm{rot}}}{V_\mathrm{\Delta }}\right)^3.$$
$`(6)`$
Because by definition baryons outside the virial radius have yet to reach the galaxy, $`M_{\mathrm{disk}}^{\mathrm{max}}`$ is a firm upper bound to the baryonic fraction transformed into stars, implying that $`f_{\mathrm{bdsk}}<\mathrm{\hspace{0.17em}1}`$ (White et al. 1993).
The slope and zero point of the Tully-Fisher relation therefore implies that the fraction of the total mass (and of baryons) transformed into stars is a sensitive function of the cosmological parameters (through the product $`\mathrm{\Omega }_0h`$, and $`\mathrm{\Delta }`$), of the stellar mass-to-light ratio, and of the ratio between the rotation speed of the disk and the circular velocity of the surrounding halo.
#### 2.1.1 Constraints from the slope of the Tully-Fisher relation
As indicated by eq. 5, reproducing the observed slope of the Tully-Fisher relation entails a delicate balance between $`f_{\mathrm{mdsk}}`$, $`\mathrm{{\rm Y}}_I`$, and the ratio $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$. The simplest possibility is that the three parameters are approximately constant in all halos. This is the case argued by Mo et al. (1998), who suggest that $`f_{\mathrm{mdsk}}5\times 10^2`$, $`\mathrm{{\rm Y}}_I1.7h`$, and $`V_{\mathrm{rot}}/V_\mathrm{\Delta }1.5`$ are needed in order to reproduce observations of galaxy disks. Although plausible, this assumption is at odds with the results of the numerical experiments we present below (§3), so it is worthwhile considering a second possibility: that $`f_{\mathrm{mdsk}}`$, $`\mathrm{{\rm Y}}_I`$, and $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ are not constant from halo to halo but that their variations are strongly correlated, in the manner prescribed by eq. 5. We explore now how such correlation may emerge as a result of the dynamical response of the dark halo to the assembly of a disk galaxy at its center.
Fig. 2.—The disk mass fraction versus the ratio between disk rotation speed and halo circular velocity. The thick dashed and solid lines correspond to the constraint imposed on these two quantities by the Tully-Fisher relation (eq. 5) in the $`\mathrm{\Lambda }`$CDM and SCDM scenarios, respectively. Dotted lines correspond to the relation expected for galaxies assembled in NFW halos of constant “concentration” parameter, as labeled. Constant disk mass-to-light ratios are assumed throughout; $`\mathrm{{\rm Y}}_I=2`$ in the upper panel and $`\mathrm{{\rm Y}}_I=1`$ in the lower one, respectively.
In the interest of simplicity, we shall restrict our analysis to a case where the disk mass-to-light ratio, $`\mathrm{{\rm Y}}_I`$, is assumed to be constant, and we shall use the “adiabatic contraction” approximation to compute the dependence of the velocity ratio $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ on $`f_{\mathrm{mdsk}}`$. The resulting relation is sensitive to the detailed mass profile of the dark halo and to the assumed structure of the disk. Assuming that the model proposed by Navarro, Frenk & White (NFW96, NFW97) <sup>3</sup><sup>3</sup>3According to these authors, the density profile of a dark halo is well approximated by $`\rho (r)=\delta _c\rho _{crit}(r/r_s)^1(1+r/r_s)^2`$, where $`r_s`$ is a scale radius and $`\delta _c`$ is a characteristic density contrast. For halos of given mass, $`M_\mathrm{\Delta }`$, this formula has a single free parameter, which can be expressed as the “concentration”, $`c=r_\mathrm{\Delta }/r_s`$. The characteristic density contrast and the concentration are related by the simple formula, $`\delta _c=(\mathrm{\Delta }/3)c^3/(\mathrm{ln}(1+c)c/(1+c))`$. is a reasonable approximation to the structure of the halo, and adopting exponential disk models with radial scale lengths consistent with the assumption that halo and disk share the same specific angular momentum, it is possible to derive analytically the dependence of the velocity ratio on $`f_{\mathrm{mdsk}}`$ (see Mo et al 1998 for details). We show the results in Figure 2. The thick lines (solid and dashed) in this figure correspond to the constraint enunciated in eq. 5 for two different cosmological models (the “standard” Cold Dark Matter model, sCDM: with parameters $`\mathrm{\Omega }=1`$ and $`h=0.5`$, and a low-density, flat $`\mathrm{\Lambda }`$CDM model, with $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Lambda }_0=0.7`$, and $`h=0.7`$). The upper and lower panels adopt two different values for the disk mass-to-light ratio, $`\mathrm{{\rm Y}}_I=2`$, and $`1`$, respectively. The rightmost point in each of the thick lines corresponds to the maximum disk mass fraction allowed by the baryonic content of the halo, i.e., $`f_{\mathrm{mdsk}}^{\mathrm{max}}=\mathrm{\Omega }_b/\mathrm{\Omega }_0`$. Dotted lines in this figure are the “adiabatic contraction approximation” predictions for different values of the NFW concentration parameter, $`c`$.
Figure 2 illustrates a few interesting points. The first one is that the disk mass-to-light ratio and the cosmological parameters determine in practice the range of halo concentrations that are consistent with the zero-point of the Tully Fisher relation. Halos formed in the sCDM scenario must have $`c<\mathrm{\hspace{0.17em}3}`$ ($`5`$) if $`\mathrm{{\rm Y}}_I2`$ ($`1`$). This effectively rules out the sCDM scenario, since N-body simulations show that halos formed in this cosmology have much higher concentrations, typically $`c15`$-$`20`$ (NFW96). A similar conclusion was reached by van den Bosch (1999), who finds through similar considerations that the large $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ values expected in the sCDM scenario effectively rule out this cosmogony.
The low-density $`\mathrm{\Lambda }`$CDM model fares better, because the higher value of $`h`$ and the lower value of $`\mathrm{\Delta }`$ in this model imply smaller $`f_{\mathrm{mdsk}}`$ at a given value of the velocity ratio $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$. However, for $`\mathrm{{\rm Y}}_I=2`$ ($`1`$), concentrations lower than about $`5`$ ($`12`$) are needed. As discussed by NS2, high-resolution N-body simulations of halo formation in the $`\mathrm{\Lambda }`$CDM scenario yields concentrations of order $`20`$, in disagreement with these constraints unless $`\mathrm{{\rm Y}}_I1`$. Concentrations as high as this are similar to those found for sCDM, and are systematically higher than the values predicted by the approximate formula proposed by NFW97. This explains the apparent disagreement between the conclusions of NS2 and those of semi-analytic models that claim reasonable agreement with the observed Tully-Fisher relation (Mo et al 1998, van den Bosch 1999): the discrepancy can be fully traced to the lower halo central concentrations and lower disk mass-to-light ratios adopted by the latter authors. If concentrations are truly as high as reported by NS2, agreement with the Tully Fisher relation require $`\mathrm{{\rm Y}}_I1`$, in disagreement with estimates based on broad-band colors of Tully-Fisher disks (which suggest $`\mathrm{{\rm Y}}_I2`$) and with mass-to-light ratio estimates of the solar neighborhood (see NS2 for details).
The second important point that emerges from Figure 2 is that, in order to match the Tully-Fisher relation, halo concentrations must be an increasing function of $`f_{\mathrm{mdsk}}`$ (assuming $`\mathrm{{\rm Y}}_I`$constant). Indeed, as $`f_{\mathrm{mdsk}}`$ increases the thick solid and dashed lines cross (dotted) curves of increasing $`c`$. Since, according to NFW96 and NFW97, concentration depends directly on halo mass—low mass halos are systematically more centrally concentrated as a result of earlier collapse times—this is equivalent to requiring $`f_{\mathrm{mdsk}}`$ to be a function of halo mass. This is actually consistent with simple disk formation models where the mass of the disk is determined by gas cooling radiatively inside a dark halo, as first proposed by White & Rees (1978), and later worked out in detail by White & Frenk (1991). These authors show that, if disks form by gas cooling within an approximately isothermal halo, disk masses are expected to be roughly proportional to $`V_\mathrm{\Delta }^{3/2}`$ (White 1996), implying $`f_{\mathrm{mdsk}}V_\mathrm{\Delta }^{3/2}M_\mathrm{\Delta }^{1/2}`$. We show below (§4.1) that, although this dependence is stronger than found in our numerical experiments, the overall trend predicted is nicely reproduced in our numerical experiments.
Finally, Figure 2 illustrates that the structure and dynamical response of the halo to the assembly of the disk may be responsible for the small scatter in the Tully-Fisher relation. For illustration, consider two halos of the same mass, and therefore approximately similar concentration, where the fraction of baryons collected into the central galaxy, $`f_{\mathrm{mdsk}}`$, differs substantially. Provided that $`f_{\mathrm{mdsk}}>\mathrm{\hspace{0.17em}0.02}`$, where the “adiabatic contraction” dotted curves are approximately parallel to the observational constraint delineated by the thick lines, these two galaxies will lie approximately along the same Tully-Fisher relation. Even if the concentration of the two halos were to differ greatly its effect on the scatter of the Tully-Fisher relation would be relatively minor: at fixed $`f_{\mathrm{mdsk}}`$, $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ changes by only about $`20\%`$ when $`c`$ changes by a factor of two.
To summarize, assuming that the disk mass-to-light ratio is approximately constant for Tully-Fisher disks, the slope of the I-band Tully-Fisher relation may result from the combination of three effects: (i) the reduced efficiency of cooling in more massive halos, (ii) the mass dependence of halo concentrations, and (iii) the dynamical response of the halo to the disk assembly. This also offers a natural explanation for the small scatter in the observed Tully-Fisher relation. We shall use our numerical simulations to test the verisimilitude of this speculation below.
### 2.2. Circular Velocity and Angular Momentum
Another similarity between the properties of dark halos and galaxy disks concerns their angular momentum. N-body simulations show that, in terms of the dimensionless parameter, $`\lambda =J|E|^{1/2}/GM_\mathrm{\Delta }^{5/2}`$, the distribution of halo angular momenta is approximately independent of mass, redshift, and cosmological parameters, and may be approximated by a log-normal distribution peaked at around $`\lambda 0.05`$ (Cole & Lacey 1996 and references therein). ($`J`$ and $`E`$ are the total angular momentum and binding energy of the halo, respectively.)
Fig. 3.—Specific angular momentum vs circular velocity of model galaxies compared with observational data. Data correspond to the samples of Courteau (1997), Mathewson et al (1992), and the compilation of Navarro (1999). Specific angular momenta are computed from disk scalelengths and rotation speeds, assuming an exponential disk model with a flat rotation curve.
The binding energy depends on the internal structure of the halos but the structural similarity between dark halos established by NFW96 and NFW97 implies that $`E`$ is to good approximation roughly proportional to $`M_\mathrm{\Delta }V_\mathrm{\Delta }^2`$, with a very weak dependence on the characteristic density of the halo. The specific angular momentum of the halo then may be written as (see Mo et al 1998 for further details),
$$j_\mathrm{\Delta }2\frac{\lambda }{\mathrm{\Delta }^{1/2}}\frac{V_\mathrm{\Delta }^2}{H(z)}=2.8\times 10^3\frac{H_0}{H(z)}\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}$$
$$\times \left(\frac{V_\mathrm{\Delta }}{200\mathrm{km}\mathrm{s}^1}\right)^2\mathrm{km}\mathrm{s}^1h^1\mathrm{kpc},$$
$`(7)`$
where we have used the most probable value of $`\lambda =0.05`$ in the second equality (see dotted line in Figure 3). The simple velocity-squared scaling of this relation is identical to that illustrated in Figure 3 between the specific angular momentum of disks and their rotation speed,
$$j_{\mathrm{disk}}1.3\times 10^3\left(\frac{V_{\mathrm{rot}}}{200\mathrm{km}\mathrm{s}^1}\right)^2\mathrm{km}\mathrm{s}^1h^1\mathrm{kpc}$$
$`(8)`$
(solid line in Figure 3), suggestive, as in the case of the Tully-Fisher relation, of a cosmological origin for this scaling law.
Combining eqs. 7 and 8, we can express the ratio between disk and halo specific angular momenta at $`z=0`$ as,
$$f_j=\frac{j_{\mathrm{disk}}}{j_\mathrm{\Delta }}0.45\left(\frac{\mathrm{\Delta }}{200}\right)^{1/2}\left(\frac{V_{\mathrm{rot}}}{V_\mathrm{\Delta }}\right)^2.$$
$`(9)`$
If the rotation speeds of galaxy disks are approximately the same as the circular velocity of their surrounding halos, then disks must have retained about one-half of the available angular momentum during their assembly.
The velocity ratio may be eliminated using eq. 6 to obtain a relation between the fraction of baryons assembled into the disk and the angular momentum ratio,
$$f_j0.5\left(\frac{\mathrm{\Delta }}{200}\right)^{1/6}\left(\frac{f_{\mathrm{bdsk}}}{\mathrm{\Omega }_0h\mathrm{{\rm Y}}_I}\right)^{2/3}.$$
$`(10)`$
This combined constraint posed by the Tully-Fisher and the angular momentum-velocity relation is shown in Figure 4 for two different cosmological models. As in Figure 2, thick solid lines correspond to the “standard” cold dark matter model, sCDM, and thick dashed lines to the $`\mathrm{\Lambda }`$CDM model. Each curve is labelled by the value adopted for the disk mass-to-light ratio, $`\mathrm{{\rm Y}}_I`$. The precise values of $`f_{\mathrm{bdsk}}`$ and $`f_j`$ along each curve are determined by and the ratio $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$, and are shown by starred symbols for the case $`V_{\mathrm{rot}}=V_\mathrm{\Delta }`$ and $`\mathrm{{\rm Y}}_I=1`$.
Fig. 4.—The fraction of the baryons assembled into a disk galaxy ($`f_{\mathrm{bdsk}}`$) versus the ratio between the specific angular momenta of the disk and its surrounding halo ($`f_j`$). Thick solid and dashed lines correspond to the constraints imposed by the Tully-Fisher relation (Figure 1) and by the relation between rotation speed and angular momentum (Figure 3). The solid (dashed) thick line corresponds to the sCDM ($`\mathrm{\Lambda }`$CDM) scenario, shown for different values of $`\mathrm{{\rm Y}}_I`$, as labeled. Symbols correspond to simulated galaxy models as per the labels in Figure 1.
One important point illustrated by Figure 4 is that disk galaxies formed in a low-density universe, such as $`\mathrm{\Lambda }`$CDM, need only accrete a small fraction of the total baryonic mass to match the zero-point of the Tully-Fisher relation, but must draw a much larger fraction of the available angular momentum to be consistent with the spins of spiral galaxies. For example, if $`V_{\mathrm{rot}}=V_\mathrm{\Delta }`$ and $`\mathrm{{\rm Y}}_I=1`$, disk masses amount to only about $`30\%`$ of the total baryonic mass of the halo but contain about $`60\%`$ of the available angular momentum. This is intriguing and, at face value, counterintuitive. Angular momentum is typically concentrated in the outer regions of the system (see, e.g., Figure 9 in NS97), presumably the ones least likely to cool and be accreted into the disk, so it is puzzling that galaxies manage to tap a large fraction of the available angular momentum whilst collecting a small fraction of the total mass. The simulations in NS97, which include the presence of a strong photo-ionizing UV background, illustrate exactly this dilemma; the UV background suppresses the cooling of late-infalling, low-density, high-angular momentum gas and reduces the angular momentum of cold gaseous disks assembled at the center of dark matter halos.
The situation is less severe in high-density universes such as sCDM; we see from Figure 4 that disks are required to collect similar fractions of mass and of angular momentum in order to match simultaneously the Tully-Fisher and the spin-velocity relations. As a result, any difficulty matching the angular momentum of disk galaxies in sCDM will become only worse in a low-density $`\mathrm{\Lambda }`$CDM universe.
Note as well that the problem becomes more severe the lower the mass-to-light ratio of the disk. Indeed, from the point of view of this constraint, it would be desirable for disks formed in the $`\mathrm{\Lambda }`$CDM scenario to have $`\mathrm{{\rm Y}}_I>2`$; in this case $`f_jf_{\mathrm{bdsk}}`$ would be consistent with the constraint posed by observed scaling laws. However, this is the opposite of what was required to reconcile highly-concentrated halos with the zero-point of the Tully-Fisher relation. This conundrum illustrates the fact that accounting simultaneously for the mass and angular momentum of disk galaxies represents a serious challenge to hierarchical models of galaxy formation.
## 3. Numerical Experiments
### 3.1. The Code
The simulations were performed using GRAPESPH, a code that combines the hardware N-body integrator GRAPE with the Smooth Particle Hydrodynamics technique (Steinmetz 1996). GRAPESPH is fully Lagrangian and optimally suited to study the formation of highly non-linear systems in a cosmological context. The version used here includes the self-gravity of gas, stars, and dark matter components, a three-dimensional treatment of the hydrodynamics of the gas, Compton and radiative cooling, the effects of a photo-ionizing UV background (NS97), and a simple recipe for transforming gas into stars.
### 3.2. The Star Formation and Feedback Algorithm
The numerical recipe for star formation, feedback and metal enrichment is similar to that described in Steinmetz & Müller (1994, 1995, see also Katz 1992, Navarro & White 1993). Full details on our present implementation are presented in Steinmetz & Navarro (2000), together with validating and calibrating tests. A brief description follows.
“Star particles” are created in collapsing regions that are locally Jeans unstable at a rate controlled by the local cooling and dynamical timescales, $`\dot{\rho }_{}=c_{}\rho _{gas}/\mathrm{max}(\tau _{cool},\tau _{dyn})`$. The proportionality parameter, $`c_{}`$, effectively controls the depletion timescale of gas, which in high-density regions, where most star formation takes place and where $`\tau _{dyn}\tau _{cool}`$, is of order $`\tau _{dyn}/c_{}(4\pi G\rho _{gas})^{1/2}c_{}^1`$.
The equations of motion of star particles are only affected by gravitational forces, but newly formed stars devolve $`10^{49}`$ ergs (per solar mass of stars formed) about $`10^7`$ yrs after their formation. This energy input, a crude approximation to the energetic feedback from evolving massive stars and supernovae, is largely invested in raising the temperature of the surrounding gas. As discussed by Katz (1992) and Navarro & White (1993), this form of energetic feedback is rather inefficient; the high densities typical of star forming regions imply short cooling timescales that minimize the hydrodynamical effects of the feedback energy input. The net result is that the star formation history of an object simulated using this “minimal feedback” formulation traces closely the rate at which gas cools and collapses within dark matter halos (SN1).
We have generalized this formulation by assuming that certain fraction of the available energy, $`ϵ_v`$, is invested in modifying the kinetic energy of the surrounding gas. These motions can still be dissipated through shocks but on longer timescales, leading to overall reduced star formation efficiencies and longer effective timescales for the conversion of gas into stars.
We have determined plausible values for the two parameters, $`c_{}`$ and $`ϵ_v`$, by comparing the star forming properties of isolated disk galaxy models with the empirical “Schmidt-law” correlations between star formation rate and gas surface density reported by Kennicutt (1998, see SN3 for full details). In the case of “minimal feedback” ($`ϵ_v=0`$) low values of $`c_{}`$, typically $`0.05`$, are needed to prevent the rapid transformation of all of the gas in a typical galaxy disk into stars. Introducing a kinetic component to the feedback prolongs the depletion timescales somewhat, but in all cases, best results are obtained by choosing long star formation timescales, i.e., values of $`c_{}0.05`$. For such low values of $`c_{}`$, $`ϵ_v`$ must be less than $`0.2`$ in order to prevent slowing down star formation to rates significantly lower than observed in Kennicutt’s empirical correlations. Similar constraints on $`ϵ_v`$ can be derived from high resolution 1D simulations of supernova remnants (Thornton et al. 1998). We shall hereafter refer to the (rather extreme) choice of $`c_{}=0.05`$ and $`ϵ_v=0.2`$ as the “kinetic feedback” case. Adopting $`c_{}>0.05`$ and $`ϵ_v<0.2`$ produces results that are intermediate between the “minimal” and “kinetic” feedback cases we report here. Tests on isolated disks of varying circular velocity indicate that the “kinetic feedback” adopted here reduces substantially the efficiency of star formation in systems with circular velocity below $`100`$ km s<sup>-1</sup>, but has a more modest influence on the star formation rates in more massive systems, in rough agreement with current interpretation of observational data (Martin 1999).
### 3.3. The Initial Conditions
We investigate two variants of the Cold Dark Matter scenario. The first is the former “standard” CDM model, with cosmological parameters $`\mathrm{\Omega }=1`$, $`h=0.5`$, $`\mathrm{\Lambda }=0`$, normalized so that at $`z=0`$ the rms amplitude of mass fluctuations in $`8h^1`$ Mpc spheres is $`\sigma _8=0.63`$. Although this model fails to reproduce a number of key observations, such as the CMB fluctuations detected by COBE, it remains popular as a well-specified cosmological testbed and as a well studied example of a hierarchical clustering model of galaxy formation. The second is the currently popular low-density CDM model that includes a non-zero cosmological constant and which is normalized to match COBE constraints and current estimates of the Hubble constant, $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Lambda }=0.7`$, $`h=0.7`$, $`\sigma _8=1.1`$. Both models assume a value of $`\mathrm{\Omega }_b=0.0125h^2`$ for the baryon density parameter of the universe, consistent with constraints from Big-Bang nucleosynthesis of the light elements (Schramm & Turner 1997).
In both cosmologies, we simulate regions that evolve to form dark halos with circular velocities in the range (80, 350) km s<sup>-1</sup> at $`z=0`$. These regions are selected from cosmological simulations of large periodic boxes and are resimulated individually including the full tidal field of the original calculation. All resimulations have typically $`32,000`$ gas particles and the same amount of dark matter particles. The size of the resimulated region scales with the circular velocity of the selected halo, so that most systems have similar numbers of particles at $`z=0`$, regardless of circular velocity. This is important to ensure that numerical resolution is approximately uniform across all systems. Gas particle masses range from $`2.5\times 10^6h^1M_{}`$ to $`9\times 10^7h^1M_{}`$, depending on the system being simulated. Dark matter particle masses are a factor $`(\mathrm{\Omega }_0\mathrm{\Omega }_b)/\mathrm{\Omega }_b`$ more massive. Their mass is low enough to prevent artificial suppression of cooling due to collisional effects (Steinmetz & White 1997). Runs start at $`z=21`$ and use gravitational softenings that range between $`0.5`$ and $`1.0h^1`$ kpc.
We have concentrated our numerical efforts on the sCDM scenario: about $`60`$ galaxy models satisfy the conditions for analysis outlined below (§3.4). For comparison, only seven $`\mathrm{\Lambda }`$CDM galaxy models are considered here. This is because of our realization during the course of this study that the cosmological scenario has very modest effects on the scaling laws we discuss here. Indeed, regarding disk scaling laws, galaxy models formed in the sCDM and $`\mathrm{\Lambda }`$CDM scenarios differ mainly because of the adoption of different values of the Hubble constant. We discuss this in more detail below.
### 3.4. Identification and Analysis of Model Galaxies
Model galaxies are easily identified in our runs as star and gas “clumps” with very high density contrast. We retain for analysis only galaxies in halos represented with more than $`500`$ dark particles inside the virial radius; most of these systems have more than $`1,000`$ star particles in each galaxy. This list is culled to remove obvious ongoing mergers and satellites orbiting the main central galaxy of each halo. The properties of the luminous component are computed within a fiducial radius, $`r_{gal}=15(V_\mathrm{\Delta }/220`$ km s$`{}_{}{}^{1})h^1`$ kpc. This radius contains all of the baryonic material associated with the galaxy and is much larger than the spatial resolution of the simulations. The rotation speeds we quote are also measured at that radius, although in practice the circular velocity in the models is rather insensitive to the radius where it is measured: similar results are obtained using $`r_{gal}=3.5(V_\mathrm{\Delta }/220`$ km s$`{}_{}{}^{1})h^1`$ kpc (see Figure 1 of SN1).
Although the resolution of the modeled galaxies is adequate to compute reliably quantities such as the total mass or circular velocity as a function of radius, it is still insufficient to gain insight into the morphology of the galaxy. As a result, our sample is likely to contain a mixture of Hubble types, from S0 to Sd (most of the simulated galaxies retained for analysis are largely rotationally supported).
Galaxy luminosities are computed by simply adding the luminosities of each star particle, taking into account the time of creation of each particle and using the latest version of the spectrophotometric models of Bruzual & Charlot (G.Bruzual & S.Charlot 1996, unpublished), see Contardo, Steinmetz & Fritze-von Alvensleben (1998) for details. Corrections due to internal absorption and inclination are neglected. The IMF is assumed in all cases to be independent of time.
Fig. 5.—Correlations between disk mass, $`M_{\mathrm{disk}}`$, disk rotation speed (measured at $`r_{gal}`$), $`V_{\mathrm{rot}}`$, absolute I-band magnitude at $`z=0`$, $`M_I`$, and halo circular velocity, $`V_\mathrm{\Delta }`$, found in the numerical experiments. Symbols are as in Figure 1. Solid line in the lower-left panel is the best fit to the observational data shown in Figure 1. Solid and dashed lines in the upper left panel correspond to constant disk mass-to-light ratio, $`\mathrm{{\rm Y}}_I=2.5`$. Solid line in upper-right panel outline the loci of galaxies that have assembled all available baryons. Slopes quoted in all panels correspond to unweighted least-squares fits. The rms scatter values correspond to the $`y`$-axis quantity relative to the least-squares fit, except for the $`M_I`$-$`V_{\mathrm{rot}}`$ panel, where it corresponds to the scatter in $`M_I`$.
## 4. Numerical Scaling Laws
### 4.1. The numerical Tully-Fisher relation
The symbols with horizontal error bars in Figure 1 show the numerical Tully-Fisher relation obtained in our simulations. Solid squares and open circles denote the luminosities and rotation speeds (measured at $`r_{gal}`$) of galaxy models formed in the sCDM and $`\mathrm{\Lambda }`$CDM scenarios, respectively, under the “minimal feedback” assumption. Starred symbols correspond to the “kinetic feedback” case applied to the sCDM model. Error bars span the range in luminosities corresponding to assuming either a Salpeter or a Scalo stellar initial mass function.
A few points are clear from Figure 1. (i) The scatter and slope of the numerical Tully-Fisher relation are in good agreement with observation. (ii) The zero point of the numerical relation is offset by almost $`2`$ ($`1.25`$) magnitudes for sCDM ($`\mathrm{\Lambda }`$CDM) galaxy models. (iii) The effects of kinetic feedback on the numerical Tully-Fisher relation are quite modest. Indeed, only a slight dimming in galaxies with $`V_{\mathrm{rot}}<\mathrm{\hspace{0.17em}100}`$ km s<sup>-1</sup> is clearly noticeable in Figure 1. We analyze these three results in detail below.
#### 4.1.1 The Slope
As discussed in §2.1, agreement between the slopes of the numerical and observed $`I`$-band Tully-Fisher relation imply a tight relation between the fraction of the mass assembled into the galaxy, $`f_{\mathrm{mdsk}}`$, the disk mass-to-light ratio, $`\mathrm{{\rm Y}}_I`$, and the ratio between disk and halo circular velocities, $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ (see eq. 5). Correlations between these parameters obtained in our numerical simulations are shown in Figure 5.
The bottom left panel of Figure 5 is identical to Figure 1, but labels includes the (unweighted) best-fit slope and rms deviation (in magnitudes) of each numerical relation. Symbols are as in Figure 1. As mentioned above, the slopes of the sCDM relations are consistent with the observed one. The same applies to the $`\mathrm{\Lambda }`$CDM best-fit slope, which is virtually indistinguishable from the sCDM one. This result is only weakly affected by the narrow dynamic range covered by the $`\mathrm{\Lambda }`$CDM simulations. Indeed, extending our analysis to more poorly resolved clumps in the $`\mathrm{\Lambda }`$CDM runs (systems with $`V_{\mathrm{rot}}<200`$ km s<sup>-1</sup>, see dotted open circles in this panel) shows that there is little significant difference between the effective slopes obtained in the sCDM and $`\mathrm{\Lambda }`$CDM models (see also Figure 2 of NS2). The good agreement between observed and numerical slopes indicate that $`f_{\mathrm{mdsk}}`$, $`\mathrm{{\rm Y}}_I`$, and $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ approximately follow the constraint prescribed by eq. 5. The rest of the panels in Figure 5 illustrate the nature of this relation.
The upper left panel in Figure 5 shows that, to a large extent, the stellar mass-to-light ratios of galaxy models are approximately constant; the dotted and dashed lines in this panel correspond to $`\mathrm{{\rm Y}}_I=2.5=`$constant. On the other hand, the disk mass fraction is far from constant from system to system. As illustrated by the upper right panel in Figure 5, baryons in low circular velocity halos are much more effective at condensing into central galaxies than those in more massive systems. Indeed, as $`V_\mathrm{\Delta }`$ increases, the numerical results move away from the solid line representing galaxies which have assembled all of the available baryons. <sup>4</sup><sup>4</sup>4The symbols corresponding to the $`\mathrm{\Lambda }`$CDM model (open circles) have been rescaled downwards in the upper right panel of Figure 5 so that the $`M_{\mathrm{disk}}=M_{\mathrm{disk}}^{\mathrm{max}}`$ solid line is the same as in sCDM. Vertical deviations from this solid line thus indicate in both cases variations in the fraction of baryons assembled into the central galaxy. The fraction of baryons assembled into the central galaxy varies from $`70\%`$ in halos with $`V_\mathrm{\Delta }<100`$ km s<sup>-1</sup> to less than $`20\%`$ in $`V_\mathrm{\Delta }300`$ km s<sup>-1</sup> halos. This result is unlikely to be an artifact of limited numerical resolution since, as discussed in §3.3, the numerical resolution of the simulations do not depend systematically on halo circular velocity. Furthermore, the same trend was found in the convergence study of Navarro & Steinmetz (NS97, see their Table 1); the observed trend between $`f_{\mathrm{mdsk}}`$ and $`V_\mathrm{\Delta }`$ must therefore reflect the decreasing efficiency of gas cooling in halos of increasing mass discussed in §2.1.1.
In spite of the large systematic variations in the disk mass fraction put in evidence by Figure 5, the numerical slope of the Tully-Fisher relation is in good accord with observations because $`f_{\mathrm{mdsk}}`$ variations are compensated by corresponding changes in the $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ ratio. Low mass halos have higher $`f_{\mathrm{mdsk}}`$ values, but also higher $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ ratios than more massive ones (bottom right panel of Figure 5). For constant $`\mathrm{{\rm Y}}_I`$, agreement with the observed TF slope requires $`f_{\mathrm{mdsk}}=M_{\mathrm{disk}}/M_\mathrm{\Delta }M_{\mathrm{disk}}/V_\mathrm{\Delta }^3(V_{\mathrm{rot}}/V_\mathrm{\Delta })^3`$ (see eq. 5), so that if $`M_{\mathrm{disk}}V_\mathrm{\Delta }^\alpha `$, then $`V_{\mathrm{rot}}V_\mathrm{\Delta }^{\alpha /3}`$ must be satisfied. Inspection of the exponents (“slopes”) listed in the right-hand panels of Figure 5 demonstrate that this condition is approximately satisfied in the numerical experiments.
Fig. 6.—As in Figure 2, but including the results of the numerical simulations. The dotted line shows the proportionality $`f_{\mathrm{mdsk}}(V_{\mathrm{rot}}/V_\mathrm{\Delta })^3`$. The thick dashed line shows the “adiabatic contraction” prediction assuming that a $`f_{\mathrm{mdsk}}V_\mathrm{\Delta }^{0.7}`$, and that the NFW concentration parameter is given by $`c=20(V_\mathrm{\Delta }/100`$ km s$`{}_{}{}^{1})^{1/3}`$, as found in the numerical experiments. See text for further details.
Figure 6 illustrates this conclusion more explicitly: the numerical results follow very closely the $`f_{\mathrm{mdsk}}(V_{\mathrm{rot}}/V_\mathrm{\Delta })^3`$ proportionality outlined by the dotted line. As discussed in §2.1.1, this proportionality results from the combination of three effects: (i) the decreasing disk mass fraction in halos of increasing mass ($`f_{\mathrm{mdsk}}V_\mathrm{\Delta }^{0.7}`$ approximately, see upper-right panel of Figure 5), (ii) the decrease in concentration of halos of increasing mass ($`c20(V_\mathrm{\Delta }/100`$ km s$`{}_{}{}^{1})^{1/3}`$, see NFW96 and NS2), and (iii) the response of the dark halo to the assembly of the disk. The thick dashed line in Figure 6 shows the $`f_{\mathrm{mdsk}}`$ vs. $`(V_{\mathrm{rot}}/V_\mathrm{\Delta })`$ relation that results from these three premises using the adiabatic contraction approximation described in §2.1.1; it clearly reproduces the numerical results remarkably well.
We conclude that the agreement between the slopes of the numerical and observed Tully-Fisher relations lends support to the view that the Tully-Fisher slope is the (non-trivial) result of the combined effects of the halo structure and of its response to the assembly of the baryonic component of the galaxy.
#### 4.1.2 The Zero Point
In contrast with the good agreement found between models and observations for the Tully-Fisher slope, it is clear from Figure 1 that there is a serious mismatch in the zero-point of the numerical and observed Tully-Fisher relations. This is reflected in Figure 6 as a systematic offset between the numerical data points and the thick curves labeled “sCDM” and “$`\mathrm{\Lambda }`$CDM” (analogous to the lines in Figure 2, but drawn here for $`\mathrm{{\rm Y}}_I=2.5`$, as appropriate for our numerical results). At given $`f_{\mathrm{mdsk}}`$, the ratio $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ is much larger than required by observations, leading to galaxies which, at fixed rotation speed, are too faint to be consistent with observations, as clearly shown in Figure 1. Simulated galaxies formed in the sCDM and $`\mathrm{\Lambda }`$CDM scenarios are physically very similar, since they have similar disk mass-to-light ratios and share the same location in the $`f_{\mathrm{disk}}`$ vs. $`(V_{\mathrm{rot}}/V_\mathrm{\Delta })`$ plane. The reason why $`\mathrm{\Lambda }`$CDM models appear to be in better agreement with observations in Figure 1 is largely due to the adoption of a higher value of Hubble’s constant for this model. This issue is discussed in detail by NS2. From the analysis in that paper and the discussion in §2.1.1, it is clear that the zero-point discrepancy is testimony to the large central concentrations of dark halos formed in the two cosmologies we explore here.
#### 4.1.3 The Scatter
According to eq. 4, the scatter in the numerical Tully-Fisher relation is determined by the variance in the mass-to-light ratio, $`\mathrm{{\rm Y}}_I`$, the disk mass fraction, $`f_{\mathrm{mdsk}}`$, and the velocity ratio, $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$. These can be derived from the rms values quoted in the right-hand panels and upper-left panel of Figure 5, respectively, and would imply, taken at face value, a large dispersion ($`>\mathrm{\hspace{0.17em}0.6}`$ mag rms) for the numerical Tully-Fisher relation. This is actually about three times larger than measured in the simulations (see rms values quoted in bottom left panel of Figure 5). The reason behind the small scatter is once again linked to the dynamical response of the halo to the disk assembly, as discussed in §2.1.1. At fixed halo mass (i.e., fixed $`V_\mathrm{\Delta }`$), galaxies where fewer than average baryons collect into the central galaxy have lower than average $`V_{\mathrm{rot}}/V_\mathrm{\Delta }`$ ratios, and viceversa. This is shown in Figure 7, where we plot the residuals from the best-fit power laws to the data presented in the right-hand panels of Figure 5. As anticipated in §2.1.1, the residuals scale in such a way that variations in the fraction of baryons assembled into galaxies scatter along the observed Tully-Fisher relation, reducing substantially the resulting scatter. This helps to reconcile the large scatter predicted by analytical estimates (Eisenstein & Loeb 1996) with the results of our numerical simulations. Again, the halo structure and response to the assembly of the galaxy are crucial for explaining the observed properties of the Tully-Fisher relation.
Fig. 7.—Residuals from least squares fits to the numerical data presented in the right-hand panels of Figure 5. Solid line corresponds to the condition $`\delta \mathrm{log}M_{\mathrm{disk}}=3\delta \mathrm{log}V_\mathrm{\Delta }`$ required so that galaxy models scatter along the Tully-Fisher relation, reducing substantially its scatter.
### 4.2. The angular momentum of simulated disks
In agreement with prior work (see, e.g., NS97 and references therein), we find that the baryonic components of the simulated galaxy models are quite deficient in angular momentum. This is easily appreciated in Figure 3, which shows that, at a given $`V_{\mathrm{rot}}`$, the specific angular momentum of observed disks exceeds that of numerical models by more than one order of magnitude. This is due to the transfer of angular momentum from the baryons to the halo associated with merger events during the formation of the disk, as first suggested by Navarro & Benz (1991), and later confirmed by Navarro, Frenk & White (1995) and NS97. Indeed, as shown in Figure 4, the baryonic components of simulated galaxies have retained, on average, less than $`15\%`$ of the specific angular momentum of their surrounding halos, placing them in the ($`f_{\mathrm{bdsk}}`$,$`f_j`$) plane well outside of the constraints imposed by the observed scaling laws between luminosity, rotation speed, and disk size.
### 4.3. The effects of feedback
Our implementation of feedback seems to have only a very modest impact on the results discussed above, even for the rather extreme “kinetic feedback” case we tried. In the case of the Tully-Fisher relation only systems with $`V_{\mathrm{rot}}<\mathrm{\hspace{0.17em}100}`$ km s<sup>-1</sup> are affected (see, e.g., starred symbols in Figure 1), and the net effect is actually to make the zero-point disagreement worse: the rotation speeds of these low-mass galaxies are roughly unchanged, but feedback makes them fainter. This shows yet again that the zero-point problem afflicting the Tully-Fisher relation is not a result of processes that affect the baryonic component, such as feedback, but of the high central concentration of dark matter near the centers of dark halos and the relatively high disk mass-to-light ratios of our models.
Feedback also seems to have a relatively minor impact on the “angular momentum problem” as well. Disks simulated with our “minimum” and “kinetic” feedback implementations have angular momenta well below those of observed spirals (Figures 3 and 4). Feedback seems able to slow down the transformation of gas into stars in low-mass systems but fails to prevent most of the gas from collapsing into dense disks at early times. The final galaxy is thus built up as the outcome of a hierarchical sequence of merger events during which the gas transfers most of its angular momentum to its surrounding halo (NS97 and references therein).
We conclude that the feedback algorithm explored here is unable to bring the angular momenta of simulated galaxies into agreement with observations. Gas must be effectively prevented from collapsing into, or removed from, dense disks at early times, a requirement that requires a much more efficient transformation of supernova energy into gas bulk motions than our algorithm can accomplish (Weil, Eke & Efstathiou 1998). However, it is difficult to see how this can be attained without violating constraints posed by observed correlations between gas density and star formation rates in isolated disks (Kennicutt 1998). As discussed in §3.2, the “kinetic feedback” model parameters adopted here are already quite extreme; for example, increasing $`ϵ_v`$ beyond our choice of $`20\%`$ in an effort to expel more gas from star forming regions would lead to models that disagree substantially with Kennicutt’s findings (SN3). Although it cannot be discounted that feedback implementations different from the one adopted here may lead to improved results, we conclude that accounting simultaneously for the luminosity, velocity, and angular momentum of spiral galaxies in hierarchical models remains an unsolved problem for the CDM cosmogonies we explore here.
## 5. Summary and Discussion
We report here the results of numerical experiments designed to explore whether cosmologically induced correlations linking the structural parameters of cold dark matter halos may serve to elucidate the origin of disk galaxy scaling laws. The numerical experiments include gravity, pressure gradients, hydrodynamical shocks, radiative cooling, heating by a UV background, and a simple recipe for star formation. Feedback from evolving stars is also taken into account by injecting energy into gas that surrounds regions of recent star formation. A large fraction of the energy input is in the form of heat, and is radiated away quickly by the dense, cool, star-forming gas. The remainder is introduced through an extra acceleration term that affects directly the kinematics of the gas in the vicinity of star forming regions.
We find that the slope of the numerical Tully-Fisher relation is in good agreement with observation, although not, as proposed in previous work, as a direct result of the cosmological equivalence between mass and circular velocity of dark halos. Rather, the agreement results from a delicate balance between the dark halo structure, the fraction of baryons collected into each galaxy, and the dynamical response of the dark halo to the assembly of the luminous component.
Massive halos are significantly less efficient at assembling baryons into galaxies than their low-mass counterparts, a trend that agrees with expectations from theoretical models where the mass of the central galaxy is determined by the efficiency of gas cooling. As a result of the structural similarity of dark halos, systems that collect a large fraction of available baryons into a central galaxy have their rotation speeds increased substantially over and above the circular velocity of their surrounding halos. The combined effect leads, for Cold Dark Matter halos, to a direct scaling between disk mass and rotation speed, $`M_{\mathrm{disk}}V_{\mathrm{rot}}^3`$, which, for approximately constant stellar $`I`$-band mass-to-light ratio, is in very good agreement with the observed slope of the $`I`$-band Tully-Fisher relation.
The scatter in the numerical Tully-Fisher relation ($`0.25`$ mag rms, see lower left panel in Figure 5) is substantially smaller than observed, a somewhat surprising result given the sizeable dispersions in disk mass-to-light ratios and disk mass fractions, as well as in the ratios between disk rotation speed and halo circular velocity observed in our numerical simulations. The small scatter is mainly due to the fact that, although disk mass fractions may vary widely from galaxy to galaxy, these variations are strongly correlated with variations in the disk rotation speed. Model galaxies therefore scatter along the Tully-Fisher relation, minimizing the resulting scatter (Figure 7). This is again a non-trivial result stemming from the detailed dark halo structure and from its response to the assembly of the luminous component of the galaxy.
Does the agreement extend to passbands other than $`I`$? Since our galaxy models have approximately constant mass-to-light ratio in the $`I`$-band, the simulated slope will be approximately similar in all passbands where mass approximately traces light, such as in the infrared $`K`$ band. This is consistent with the results of Verheijen (1997), who find a slope of $`7.0`$ in the $`I`$-band and of $`7.9`$ in $`K`$, when his complete sample of galaxies is analyzed. The $`K`$-band slope is significantly steeper than the $`I`$-band’s ($`10.4`$ in $`K`$ versus $`8.7`$ in $`I`$) only when a restricted sample with favorable inclination, smooth morphology, steep HI profile edges, and free from bars is considered. The strong dependence on sample selection reflects the large uncertainties in the numerical values of the slope that result from the small dynamic range in velocity (less than a factor of $`3`$ typically) covered by existing Tully-Fisher samples. Taking this into account, and considering that the selection criteria of Verheijen’s complete sample is closer to those in our analysis (§3.4), we conclude that our models are generally consistent with the slopes of the $`I`$ and $`K`$ band Tully-Fisher relations.
The Tully-Fisher relation is also known to be significantly shallower in bluer passbands (see, e.g., Verheijen 1997) a result that can be traced to systematic variations in the mass-to-light ratios as a function of disk rotation speed. Qualitatively our simulations can reproduce the trend, although it is quantitatively much weaker than observed (SN1). This is because fewer stars are formed late in our models than implied by observations (SN3). This effect may have led us to underestimate the scatter in $`\mathrm{{\rm Y}}_I`$. However, considering that the measured scatter is much smaller than the observed one, we still consider our conclusion that slope and scatter are consistent with observations to be very solid.
In contrast with this success, we find that the zero-point of the numerical Tully-Fisher relation is offset by about $`2`$ ($`1.25`$) magnitudes relative to observations for galaxies formed in the sCDM ($`\mathrm{\Lambda }`$CDM) scenario. The offset can be traced to the high-central concentration of dark matter halos formed in these cosmogonies and, to a lesser extent, to the relatively high stellar mass-to-light ratios found in our numerical models ($`\mathrm{{\rm Y}}_I2.5`$). The “concentrations” of dark matter halos would have to be reduced by a factor of $`3`$-$`5`$ or, alternatively, $`\mathrm{{\rm Y}}_I`$ would have to be reduced to values as low as $`0.4`$-$`0.5`$ in order to restore agreement with the observed zero-point. As discussed by NS2, neither alternative is quite palatable, since they involve substantial modifications either to the underlying cosmological model, to our understanding of the spectrophotometric evolution of stars, or to our assumptions about the initial stellar mass function. Furthermore, mass-to-light ratios as low as that would exacerbate the problem posed by the small fraction of mass and the large fraction of angular momentum collected by the central disk in low-density universes (see §2.2 and Figure 4).
Including the energetic feedback from evolving stars and supernovae does not improve substantially the agreement between the properties of model galaxies and observation. This is because, in order to fit the empirical correlation linking star formation rate to gas surface density, our feedback model is only efficient in low-mass systems with circular velocities below $`100`$ km s<sup>-1</sup>. The progenitors of massive spirals typically exceed the threshold circular velocity at high redshift, leading to efficient gas cooling and to early onset of star formation. The effects of the kinetic feedback implementation on the Tully-Fisher relation are therefore minor, and only noticeable in systems with low circular velocities. The zero-point disagreement actually worsens at the faint end as star formation is slowed down, affecting the absolute magnitudes of model galaxies more than their rotation speeds. The limited impact of our feedback algorithm also implies that the luminous component of model galaxies is assembled through a sequence of mergers, accompanied by a substantial loss of its angular momentum to the surrounding dark matter. The angular momentum of model galaxies is, as a consequence, about one order of magnitude less that that of observed spirals, in agreement with previous work (cf. NS97).
## 6. Concluding Remarks
The results we discuss here illustrate the difficulties faced by hierarchical models that envision the formation of disk galaxies as the final outcome of a sequence of merger events. Agreement with observations appears to require two major modifications to our modeling. (i) Dark halos that are much less centrally concentrated than those formed in the two Cold Dark Matter scenarios we explore here. (ii) Feedback effects dramatically stronger than assumed here, affecting substantially the cooling, accretion, and star formation properties of dark halos perhaps as massive as $`V_\mathrm{\Delta }200`$-$`300`$ km s<sup>-1</sup>.
The extreme feedback mechanism apparently required to bring the mass and angular momentum of disk galaxies in accordance with hierarchical galaxy formation models should have major implications on further observational clues and traces of the galaxy formation process. Given the paucity of present-day examples of this process at work (Martin 1999), we are led to speculate that star formation (and therefore feedback) episodes were much more violent in the past than in the local universe, perhaps as a result of the lower angular momenta and increased surface density of star forming regions at high redshift. Large scale winds driven by early starbursts, perhaps associated with the formation of stellar spheroids, may rid sites of galaxy formation of early-accreting, low-angular momentum baryons altogether, allowing higher-than-average angular momentum material to collapse later into the modestly star-forming, extended disks prevalent in the local universe.
It remains to be seen whether such scenario for the assembly and transformation of baryons into galaxies withstands observational scrutiny, but so far the $`z3`$ evidence being gathered from “Ly-break” galaxies seems to point to a past where starbursts may have been the norm rather than the exception (Heckmann 1999). Resolving the puzzle created by disk galaxy scaling laws has thus the potential to unravel questions of fundamental importance to our understanding of the assembly and evolution of galaxies in a cosmological context.
This work has been supported by the National Aeronautics and Space Administration under NASA grant NAG 5-7151 and by NSERC Research Grant 203263-98. MS and JFN are supported in part by fellowships from the Alfred P. Sloan Foundation. MS is supported in part by fellowship from the David & Lucile Packard foundation. We thank Stephane Courteau, Riccardo Giovanelli and the MarkIII collaboration for making their data available in electronic form, and Cedric Lacey for making available his compilation of the Mathewson et al. data. JFN acknowledges enlightening discussions with Carlos Frenk and Simon White. |
warning/0001/physics0001009.html | ar5iv | text | # Abstract
### Abstract
Starting from a solution of the problem of a mechanical oscillator coupled to a scalar field inside a reflecting sphere of radius $`R`$, we study the behaviour of the system in free space as the limit of an arbitrarily large radius in the confined solution. From a mathematical point of view we show that this way of facing the problem is not equivalent to consider the system a priori embedded in infinite space. In particular, the matrix elements of the transformation turning the system to principal axis, do not tend to distributions in the limit of an arbitrarily large sphere as it should be the case if the two procedures were mathematically equivalent. Also, we introduce ”dressed” coordinates which allow an exact description of the oscillator radiation process for any value of the coupling, strong or weak. In the case of weak coupling, we recover from our exact expressions the well known decay formulas from perturbation theory.
## 1 Introduction
Since a long time ago the experimental and theoretical investigations on the polarization of atoms by optical pumping and the possibility of detecting changes in their polarization states has allowed the observation of resonant effects associated to the coupling of these atoms with strong radiofrequency fields . As remarked in , the theoretical understanding of these effects using perturbative methods requires the calculation of very high-order terms in perturbation theory, what makes the standard Feynman diagrams technique practically unreliable in those cases. The trials of treating non-perturbativelly such kind of systems consisting of an atom coupled to the electromagnetic field, have lead to the idea of ”dressed atom”, introduced in refs and . This approach consists in quantizing the electromagnetic field and analyzing the whole system consisting of the atom coupled to the electromagnetic field. Along the years since then, this concept has been extensively used to investigate several situations involving the interaction of atoms and electromagnetic fields. For instance, atoms embedded in a strong radiofrequency field background in refs. and , atoms in intense resonant laser beans in ref. or the study of photon correlations and quantum jumps. In this last situation, as showed in refs. , and , the statistical properties of the random sequence of outcoming pulses can be analyzed by a broadband photodetector and the dressed atom approach provides a convenient theoretical framework to perform this analysis.
Besides the idea of dressed atom in itself, another aspect that desserves attention is the non-linear character of the problem involved in realistic situations, which implies, as noted above, in very hard mathematical problems to be dealt with. An way to circunvect these mathematical difficulties, is to assume that under certain conditions the coupled atom-electromagnetic field system may be approximated by the system composed of an harmonic oscillator coupled linearly to the field trough some effective coupling constant $`g`$.
In this sense, in a slightly different context, recently a significative number of works has been spared to the study of cavity QED, in particular to the theoretical investigation of higher-generation Schrodinger cat-states in high-Q cavities, as has been done for instance in . Linear approximations of this type have been applied along the last years in quantum optics to study decoherence, by assuming a linear coupling between a cavity harmonic mode and a thermal bath of oscillators at zero temperature, as it has been done in and . To investigate decoherence of higher generation Schrodinger cat-states the cavity field reduced matrix for these states could be calculated either by evaluating the normal-ordering characteristic function, or by solving the evolution equation for the field-resevoir state using the normal mode expansion, generalizing the analysis of and .
In this paper we adopt a general physicist’s point of view, we do not intend to describe the specific features of a particular physical situation, instead we analyse a simplified linear version of the atom-field system and we try to extract the more detailed information we can from this model. We take a linear simplified model in order to try to have a clearer understanding of what we believe is one of the essential points, namely, the need of non-perturbative analytical treatments to coupled systems, which is the basic problem underlying the idea of dressed atom. Of course, such an approach to a realistic non-linear system is an extremelly hard task and here we make what we think is a good agreement between physical reality and mathematical reliability, with the hope that in future work our approach could be transposed to more realistic situations.
We consider a non relativistic system composed of a harmonic oscillator coupled linearly to a scalar field in ordinary Euclidean $`3`$-dimensional space. We start from an analysis of the same system confined in a reflecting sphere of radius $`R`$, and we assume that the free space solution to the radiating oscillator should be obtained taking a radius arbitrarily large in the $`R`$-dependent quantities. The limit of an arbitrarily large radius in the mathematics of the confined system is taken as a good description of the ordinary situation of the radiating oscillator in free space. We will see that this is not equivalent to the alternative continuous formulation in terms of distributions, which is the case when we consider a priori the system in unlimited space. The limiting procedure adopted here allows to avoid the inherent ambiguities present in the continuous formulation. From a physical point of view we give a non-perturbative treatment to the oscillator radiation introducing some coordinates that allow to divide the coupled system into two parts, the ”dressed” oscillator and the field, what makes unecessary to work directly with the concepts of ”bare” oscillator, field and interaction to study the radiation process. These are the main reasons why we study a simplified linear system instead of a more realistic model, to make evident some subtleties of the mathematics involved in the limiting process of taking a cavity arbitrarily large, and also to exhibit an exact solution valid for weak as well as for strong coupling. These aspects would be masked in the perturbative approach used to study non-linear couplings.
We start considering a harmonic oscillator $`q_0(t)`$ of frequency $`\omega _0`$ coupled linearly to a scalar field $`\varphi (𝐫,t)`$, the whole system being confined in a sphere of radius $`R`$ centered at the oscillator position. The equations of motion are,
$$\text{}_0(t)+\omega _0^2q_0(t)=2\pi \sqrt{gc}_0^Rd^3𝐫\varphi (𝐫,t)\delta (𝐫)$$
(1)
$$\frac{1}{c^2}\frac{^2\varphi }{t^2}^2\varphi (𝐫,t)=2\pi \sqrt{gc}q_0(t)\delta (𝐫)$$
(2)
which, using a basis of spherical Bessel functions defined in the domain $`<|𝐫|<R`$, may be written as a set of equations coupling the oscillator to the harmonic field modes,
$$\text{}_0(t)+\omega _0^2q_0(t)=\eta \underset{i=1}{\overset{\mathrm{}}{}}\omega _iq_i(t)$$
(3)
$$\text{}_i(t)+\omega _i^2q_i(t)=\eta \omega _iq_0(t).$$
(4)
In the above equations, $`g`$ is a coupling constant, $`\eta =\sqrt{2g\mathrm{\Delta }\omega }`$ and $`\mathrm{\Delta }\omega =\pi c/R`$ is the interval between two neighbouring field frequencies, $`\omega _{i+1}\omega _i=\mathrm{\Delta }\omega =\pi c/R`$.
## 2 The transformation to principal axis and the eigenfrequencies spectrum
2.1 - Coupled harmonic Oscillators
Let us consider for a moment the problem of a harmonic oscillator $`q_0`$ coupled to $`N`$ other oscillators. In the limit $`N\mathrm{}`$ we recover our original situation of the coupling oscillator-field after redefinition of divergent quantities, in a manner analogous as renormalization is done in field theories. In terms of the cutoff $`N`$ the coupled equations (3) and (4) are simply rewritten taking the upper limit $`N`$ instead of $`\mathrm{}`$ for the summation in the right hand side of Eq.(3) and the system of $`N+1`$ coupled oscillators $`q_0`$ $`\{q_i\}`$ corresponds to the Hamiltonian,
$$H=\frac{1}{2}\left[p_0^2+\omega _0^2q_0^2+\underset{k=1}{\overset{N}{}}p_k^2+\omega _k^2q_k^22\eta \omega _kq_0q_k\right].$$
(5)
The Hamiltonian (5) can be turned to principal axis by means of a point tranformation,
$$q_\mu =t_\mu ^rQ_r,p_\mu =t_\mu ^rP_r,$$
(6)
performed by an orthonormal matrix $`T=(t_\mu ^r)`$, $`\mu =(0,k)`$, $`k=1,2,\mathrm{}N`$, $`r=0,\mathrm{}N`$. The subscript $`0`$ and $`k`$ refer respectively to the oscillator and the harmonic modes of the field and $`r`$ refers to the normal modes. The transformed Hamiltonian in principal axis is
$$H=\frac{1}{2}\underset{r=0}{\overset{N}{}}(P_r^2+\mathrm{\Omega }_r^2Q_r^2),$$
(7)
where the $`\mathrm{\Omega }_r`$’s are the normal frequencies corresponding to the possible collective oscillation modes of the coupled system.
Using the coordinate transformation $`q_\mu =t_\mu ^rQ_r`$ in the equations of motion and explicitly making use of the normalization condition $`_{\mu =0}^N(t_\mu ^r)^2=1`$, we get,
$$t_k^r=\frac{\eta \omega _k}{\omega _k^2\mathrm{\Omega }_r^2}t_0^r,$$
(8)
$$t_0^r=\left[1+\underset{k=1}{\overset{N}{}}\frac{\eta ^2\omega _k^2}{(\omega _k^2\mathrm{\Omega }_r^2)^2}\right]^{\frac{1}{2}}$$
(9)
and
$$\omega _0^2\mathrm{\Omega }_r^2=\eta ^2\underset{k=1}{\overset{N}{}}\frac{\omega _k^2}{\omega _k^2\mathrm{\Omega }_r^2}.$$
(10)
There are $`N+1`$ solutions $`\mathrm{\Omega }_r`$ to Eq.(10), corresponding to the $`N+1`$ normal collective oscillation modes. To have some insight into these solutions, we take $`\mathrm{\Omega }_r=\mathrm{\Omega }`$ in Eq.(10) and transform the right hand term. After some manipulations we obtain
$$\omega _0^2N\eta ^2\mathrm{\Omega }^2=\eta ^2\underset{k=1}{\overset{N}{}}\frac{\mathrm{\Omega }^2}{\omega _k^2\mathrm{\Omega }^2}$$
(11)
It is easily seen that if $`\omega _0^2>N\eta ^2`$ Eq.(11) yelds only positive solutions for $`\mathrm{\Omega }^2`$, what means that the system oscillates harmonically in all its modes. Indeed, in this case the left hand term of Eq.(11) is positive for negative values of $`\mathrm{\Omega }^2`$. Conversely the right hand term is negative for those values of $`\mathrm{\Omega }^2`$. Thus there is no negative solution of that equation when $`\omega _0^2>N\eta ^2`$. On the other hand it can be shown that if $`\omega _0^2<N\eta ^2`$, Eq.(11) has a single negative solution $`\mathrm{\Omega }_{}^2`$. In order to prove it let us define the function
$$I(\mathrm{\Omega }^2)=(\omega _0^2N\eta ^2)\mathrm{\Omega }^2\eta ^2\underset{k=1}{\overset{N}{}}\frac{\mathrm{\Omega }^2}{\omega _k^2\mathrm{\Omega }^2}$$
(12)
Accordingly Eq.(11) can be rewritten as $`I(\mathrm{\Omega }^2)=0`$. It can be noticed that $`I(\mathrm{\Omega }^2)\mathrm{}`$ as $`\mathrm{\Omega }^2\mathrm{}`$ and
$$I(\mathrm{\Omega }^2=0)=\omega _0^2N\eta ^2<0$$
(13)
Furthermore $`I(\mathrm{\Omega }^2)`$ is a monotonically decreasing function in that interval. Consequently $`I(\mathrm{\Omega }^2)=0`$ has a single negative solution when $`\omega _0^2<N\eta ^2`$ as we have pointed out. This means that there is an oscillation mode whose amplitude varies exponentially and that does not allows stationary configurations. We will not care about this last situation. Thus we assume $`\omega _0^2>N\eta ^2`$ and define the renormalized oscillator frequency $`\overline{\omega }`$ ,
$$\overline{\omega }=\sqrt{\omega _0^2N\eta ^2}.$$
(14)
In terms of the renormalized frequency Eq.(10) becomes,
$$\overline{\omega }^2\mathrm{\Omega }_r^2=\eta ^2\underset{k=1}{\overset{N}{}}\frac{\mathrm{\Omega }_r^2}{\omega _k^2\mathrm{\Omega }_r^2}.$$
(15)
From Eqs. (8), (9) and (15), a straightforward calculation shows the orthonormality relations for the transformation matrix $`(t_\mu ^r)`$.
We get the transformation matrix elements for the oscillator-field system by taking the limit $`N\mathrm{}`$ in the above equations. Recalling the definition of $`\eta `$ from Eqs. (3) and (4), we obtain after some algebraic manipulations, from Eqs. (15), (8) and (9), the matrix elements in the limit $`N\mathrm{}`$,
$$t_0^r=\frac{\mathrm{\Omega }_r}{\sqrt{\frac{R}{2\pi gc}(\mathrm{\Omega }_r^2\overline{\omega }^2)^2+\frac{1}{2}(3\mathrm{\Omega }_r^2\overline{\omega })^2+\frac{\pi gR}{2c}\mathrm{\Omega }_r^2}}$$
(16)
and
$$t_k^r=\frac{\eta \omega _k}{\omega _k^2\mathrm{\Omega }_r^2}t_0^r.$$
(17)
2.2 - The eigenfrequencies spectrum
Let us now return to the coupling oscillator-field by taking the limit $`N\mathrm{}`$ in the relations of the preceeding subsection. In this limit it becomes clear the need for the frequency renormalization in Eq.(14). It is exactly the analogous of a mass renormalization in field theory, the infinite $`\omega _0`$ is chosen in such a way as to make the renormalized frequency $`\overline{\omega }`$ finite. Remembering Eq.(15) the solutions with respect to the variable $`\mathrm{\Omega }`$ of the equation
$$\overline{\omega }^2\mathrm{\Omega }^2=\frac{2\pi gc}{R}\underset{k=1}{\overset{\mathrm{}}{}}\frac{\mathrm{\Omega }^2}{\omega _k^2\mathrm{\Omega }^2},$$
(18)
give the collective modes frequencies. We remember $`\omega _k=k\frac{\pi c}{R}`$, $`k=1,2,\mathrm{}`$, and take a positive $`x`$ such that $`\mathrm{\Omega }=x\frac{\pi c}{R}`$. Then using the identity,
$$\underset{k=1}{\overset{\mathrm{}}{}}\frac{x^2}{k^2\mathrm{\Omega }^2}=\frac{1}{2}(1\pi x\mathrm{cot}\pi x),$$
(19)
Eq.(18) may be rewritten in the form,
$$cotg\pi x=\frac{c}{Rg}x+\frac{1}{\pi x}(1\frac{R\overline{\omega }^2}{\pi gc}).$$
(20)
The secant curve corresponding to the right hand side of the above equation cuts only once each branch of the cotangent in the left hand side. Thus we may label the solutions $`x_r`$ as $`x_r=r+ϵ_r`$, $`0<ϵ_r<1`$, $`r=0,1,2\mathrm{}`$, and the collective eigenfrequencies are,
$$\mathrm{\Omega }_r=(r+ϵ_r)\frac{\pi c}{R},$$
(21)
the $`ϵ`$’s satisfying the equation,
$$cot(\pi ϵ_r)=\frac{\mathrm{\Omega }_r^2\overline{\omega }^2}{\mathrm{\Omega }_r\pi g}+\frac{c}{\mathrm{\Omega }_rR}.$$
(22)
The field $`\varphi (𝐫,t)`$ can be expressed in terms of the normal modes. We start from its expansion in terms of spherical Bessel functions,
$$\varphi (𝐫,t)=c\underset{k=1}{\overset{\mathrm{}}{}}q_k(t)\varphi _k(𝐫),$$
(23)
where
$$\varphi _k(𝐫)=\frac{sin\frac{\omega _k}{c}|𝐫|}{𝐫\sqrt{2\pi R}}.$$
(24)
Using the principal axis transformation matrix together with the equations of motion we obtain an expansion for the field in terms of an orthonormal basis associated to the collective normal modes,
$$\varphi (𝐫,t)=c\underset{s=0}{\overset{\mathrm{}}{}}Q_s(t)\mathrm{\Phi }_s(𝐫),$$
(25)
where the normal collective Fourier modes
$$\mathrm{\Phi }_s(𝐫)=\underset{k}{}t_k^s\frac{sin\frac{\omega _k}{c}|𝐫|}{𝐫\sqrt{2\pi R}}$$
(26)
satisfy the equation
$$(\frac{\mathrm{\Omega }_s^2}{c^2}\mathrm{\Delta })\varphi _s(𝐫)=2\pi \sqrt{\frac{g}{c}}\delta (𝐫)t_0^s,$$
(27)
which has a solution of the form
$$\varphi (𝐫,t)=\sqrt{\frac{g}{c}}\frac{t_0^s}{2|𝐫|sin\delta _s}sin(\frac{\mathrm{\Omega }_s}{c}|𝐫|\delta _s).$$
(28)
To determine the phase $`\delta _s`$ we expand the right hand term of Eq.(28) and compare with the formal expansion (26). This imply the condition
$$sin(\frac{\mathrm{\Omega }_s}{c}R\delta _s)=0.$$
(29)
Remembering from Eq.(21) that there is $`0<ϵ_s<1`$ such that $`\mathrm{\Omega }_s=(s+ϵ_s)\frac{\pi }{R}`$, it is easy to show from the condition in Eq.(27) that the phase $`0<\delta _s<\pi `$ has the form
$$\delta _s=ϵ_s\pi .$$
(30)
Comparing Eqs.(24) and (26) and using the explicit form (16) of the matrix element $`t_0^s`$ we obtain the expansion for the field in terms of the normal collective modes,
$$\varphi (𝐫,t)=\frac{\sqrt{gc}}{2}\underset{s}{}\frac{Q_ssin(\frac{\mathrm{\Omega }_s}{c}|𝐫|\delta _s)}{|𝐫|\sqrt{sin^2\delta _s+(\frac{\eta R}{2c})^2(1\frac{sin\delta _scos\delta _s}{\mathrm{\Omega }_sR/c}})}$$
(31)
## 3 The limit $`R\mathrm{}`$ \- mathematical aspects
3.1 - Discussion of the mathematical problem
Unless explicitly stated, in the remaining of this paper the symbol $`R\mathrm{}`$ is to be understood as the situation of a cavity of fixed, arbitrarily large radius. In order to compare the behaviour of the system in a very large cavity to that it would be in free space, let us firstly consider the system embedded in an a priori infinite Euclidean space; in this case to compute the quantities describing the system means essentially to replace by integrals the discrete sums appearing in the confined problem, taking direcltly $`R=\mathrm{}`$. An alternative procedure is to compute the quantities describing the system confined in a sphere of radius $`R`$ and take the limit $`R\mathrm{}`$ afterwards. This last approach to describe the system in free space should keep in some way the ”memory” of the confined system. To be physically equivalent one should expect that the two approachs give the same results. We will see that at least from a mathematical point of view this is not exactly the case. We remark that a solution to the problem of a system composed of an oscillator coupled to a field in free space, is already known since a long time ago in the context of Bownian motion. This solution is quite different from ours, in the sense that it not concerns the system confined to a box and also that it is limited to the dipole term from the multipolar expansion to the field.
In the continuous formalism of free space the field normal modes Fourier components (analogous to the components $`\varphi _s`$ in Eq.(26)) are,
$$\varphi _\mathrm{\Omega }=h(\mathrm{\Omega })_0^{\mathrm{}}𝑑\omega \frac{\omega }{\omega ^2\mathrm{\Omega }^2}\frac{sin\frac{\omega }{c}|𝐫|}{|𝐫|},$$
(32)
where
$$h(\mathrm{\Omega })=\frac{2g\mathrm{\Omega }}{\sqrt{(\mathrm{\Omega }^2\overline{\omega }^2)^2+\pi g^2\mathrm{\Omega }^2}}$$
(33)
and where the we have taken the appropriate continuous form of Eqs.(16) and (17). Splitting $`\omega /(\omega ^2\mathrm{\Omega }^2)`$ into partial fractions we get
$$\varphi _\mathrm{\Omega }=h(\mathrm{\Omega })_{\mathrm{}}^+\mathrm{}𝑑\omega \frac{1}{\omega \mathrm{\Omega }}\frac{sin\frac{\omega }{c}|𝐫|}{|𝐫|}.$$
(34)
The pole at $`\omega =\mathrm{\Omega }`$ prevents the existence of the integral in Eq.(34). The usual way to circumvect this difficulty is to replace the integral by one of the quantities,
$$Lim_{ϵ0}_{\mathrm{}}^+\mathrm{}𝑑\omega \frac{1}{\omega (\mathrm{\Omega }\pm iϵ)}\frac{sin\frac{\omega }{c}|𝐫|}{|𝐫|}_{\mathrm{}}^+\mathrm{}𝑑\omega \delta _\pm (\omega \mathrm{\Omega })\frac{sin\frac{\omega }{c}|𝐫|}{|𝐫|},$$
(35)
where
$$\delta _\pm (\omega \mathrm{\Omega })=\frac{1}{\pi }P(\frac{1}{\omega \mathrm{\Omega }})\pm i\delta (\omega \mathrm{\Omega }),$$
(36)
with $`P`$ standing for principal value. In our case this redefinition of the normal modes Fourier components may be justified by the fact that both integrals in Eq.(35) are solutions of the equations of motion (1) and (2) for $`𝐫0`$, and so the solution should be a linear combination of them. The situation is different if we adopt the point of view of taking the limit $`R\mathrm{}`$ in the solution of the confined problem. In this case the Fourier component $`\varphi _\mathrm{\Omega }`$ is obtained by taking the limit $`R\mathrm{}`$ in the expression for the field, Eq(28), what allows to obtain an uniquely defined expression to the normal modes Fourier components, to each $`\varphi _\mathrm{\Omega }`$ corresponding a phase $`\delta _\mathrm{\Omega }`$ (the limit $`R\mathrm{}`$ of $`\delta _s`$ in Eq.(22) given by
$$cot\delta _\mathrm{\Omega }=\frac{1}{\pi g}\frac{\mathrm{\Omega }^2\overline{\omega }^2}{\mathrm{\Omega }}.$$
(37)
Also, comparing Eqs.(35), (36) and (26) we see that the adoption of the continuous formalism is equivalent to assume that in the limit $`R\mathrm{}`$ the elements $`t_i^s`$ of the transformation matrix should be replaced by $`\delta _+(\omega \mathrm{\Omega })`$ or by $`\delta _{}(\omega \mathrm{\Omega })`$. This procedure is, from a mathematical point of view, perfectly justified but at the price of loosing uniqueness in the definition of the field components.
If we take the solution of the confined problem and we compute the matrix elements $`t_i^s`$ for $`R`$ arbitrarily large, we will see in subsection 3.2 that these elements do not tend to distributions in this limit. As $`R`$ becomes larger and larger the set of non-vanishing elements $`t_i^s`$ concentrate for each $`i`$ in a small neighbourhood of $`\omega _i`$. In the limit $`R\mathrm{}`$ the whole set of the matrix elements $`t_i^s`$ contains an arbitrarily large number of elements quadratically summables . For the matrix elements $`t_0^s`$ we obtain a quadratically integrable expression.
In the continuous formulation the unit matrix, corresponding to the absence of coupling, has elements $`E_\omega ^\mathrm{\Omega }=\delta (\omega \mathrm{\Omega })`$, while if we start from the confined situation, it can be verified that in the limit $`g0`$, $`R\mathrm{}`$, the matrix $`T=(t_\mu ^s)`$ tends to the usual unit matrix of elements $`E_{\omega ,\mathrm{\Omega }}=\delta _{\omega ,\mathrm{\Omega }}`$.
The basic quantity describing the system, the transformation matrix $`T=(t_\mu ^s)`$ has, as we will see, different properties in free space, if we use the continuous formalism or if we adopt the procedure of taking the limit $`R\mathrm{}`$ from the matrix elements in the confined problem . In the first case we must define the matrix elements $`t_\omega ^\mathrm{\Omega }`$ linking free field modes to normal modes, as distributions. On the other side adopting the second procedure we will find that the limiting matrix elements $`Lim_R\mathrm{}t_i^s`$ are not distributions, but well defined finite quantities. The two procedures are not equivalent, the limit $`R\mathrm{}`$ does not commute with other operations. In this note we take as physically meaningfull the second procedure, we solve first the problem in the confined case (finite $`R`$) and take afterwards the limit of infinite (in the sense of arbitrarily large) radius of the cavity. In the next subsection we perform a detailed analysis of the limit $`R\mathrm{}`$ of the transformation matrix $`(t_\mu ^r)`$.
3.2 - The transformation matrix in the limit $`R\mathrm{}`$
From Eqs. (16) and (17) we obtain for $`R`$ arbitrarily large,
$$t_0^rLim_{\mathrm{\Delta }\mathrm{\Omega }0}t_{\overline{\omega }}^\mathrm{\Omega }\sqrt{\mathrm{\Delta }\mathrm{\Omega }}=Lim_{\mathrm{\Delta }\mathrm{\Omega }0}\frac{\sqrt{2g}\mathrm{\Omega }\sqrt{\mathrm{\Delta }\mathrm{\Omega }}}{\sqrt{(\mathrm{\Omega }^2\overline{\omega }^2)^2+\pi ^2g^2\mathrm{\Omega }^2}}.$$
(38)
and
$$t_k^r=\frac{2g\omega _k\mathrm{\Delta }\omega }{(\omega _k+\mathrm{\Omega }_r)(\omega _k\mathrm{\Omega }_r)}\frac{\mathrm{\Omega }_r}{\sqrt{(\mathrm{\Omega }_r^2\overline{\omega }^2)^2+\pi ^2g^2\mathrm{\Omega }_r^2}},$$
(39)
where we have used the fact that in this limit $`\mathrm{\Delta }\omega =\mathrm{\Delta }\mathrm{\Omega }=\frac{\pi c}{R}`$. The matrix elements $`t_{\overline{\omega }}^\mathrm{\Omega }`$ are quadratically integrable to one, $`(t_{\overline{\omega }}^\mathrm{\Omega })^2𝑑\mathrm{\Omega }=1`$, as may be seen using Cauchy theorem.
For $`R`$ arbitrarily large ($`\mathrm{\Delta }\omega =\frac{\pi c}{R}0`$), the only nonvanishing matrix elements $`t_i^r`$ are those for which $`\omega _i\mathrm{\Omega }_r\mathrm{\Delta }\omega `$. To get explicit formulas for these matrix elements in the limit $`R\mathrm{}`$ let us consider $`R`$ large enough such that we may take $`\mathrm{\Delta }\omega \mathrm{\Delta }\mathrm{\Omega }`$ and consider the points of the spectrum of eigenfrequencies $`\mathrm{\Omega }`$ inside and outside a neighbourhood $`\eta `$ (defined in Eqs.(3) and (4) of $`\omega _i`$. We note that $`R>\frac{2\pi c}{g}`$ implies $`\frac{\eta }{2}>\mathrm{\Delta }\omega `$, then we may consider $`R`$ such that the right (left) neighbourhood $`\frac{\eta }{2}`$ of $`\omega _i`$ contains an integer number, $`\kappa `$, of frequencies $`\mathrm{\Omega }_r`$,
$$\kappa \mathrm{\Delta }\omega =\frac{\eta }{2}=\sqrt{\frac{g\mathrm{\Delta }\omega }{2}}.$$
(40)
If $`R`$ is arbitrarily large we see from (40) that $`\frac{\eta }{2}`$ is arbitrarily small, but $`\kappa `$ grows at the same rate, what means firstly that the difference $`\omega _i\mathrm{\Omega }_r`$ for the $`\mathrm{\Omega }_r`$’s outside the neighbourhood $`\eta `$ of $`\omega _i`$ is abitrarily larger than $`\mathrm{\Delta }\omega `$, implying that the corresponding matrix elements $`t_i^r`$ tend to zero (see Eq.(39)). Secondly all frequencies $`\mathrm{\Omega }_r`$ inside the neighbourhood $`\eta `$ of $`\omega _i`$ are arbitrarily close to $`\omega _i`$, being in arbitrarily large number. Only the matrix elements $`t_i^r`$ corresponding to these frequencies $`\mathrm{\Omega }_r`$ inside the neighbourhood $`\eta `$ of $`\omega _i`$ are different from zero. For these we make the change of labels,
$$r=in(\omega _i\frac{\eta }{2}<\mathrm{\Omega }_r<\omega _i);r=i+n(\omega _i>\mathrm{\Omega }_r>\omega _i+\frac{\eta }{2}),$$
(41)
$`i=1,2,\mathrm{}`$. We get, from Eq.(39)
$$t_i^i=\frac{g\omega _i}{\sqrt{(\mathrm{\Omega }_r^2\overline{\omega }^2)^2+\pi ^2g^2\omega _i^2}}\frac{1}{ϵ_i}$$
(42)
and
$$t_i^{i\pm n}=\frac{g\omega _i}{\sqrt{(\mathrm{\Omega }_r^2\overline{\omega }^2)^2+\pi ^2g^2\omega _i^2}}\frac{1}{n\pm ϵ_i},$$
(43)
where $`ϵ_i`$ satisfies Eq.(22) in this case,
$$cot(\pi ϵ_i)=\frac{\omega _i^2\overline{\omega }^2}{\omega _i\pi g}.$$
(44)
Using the formula
$$\pi ^2cosec^2(\pi ϵ_i)=\frac{1}{ϵ_i}+\underset{n=1}{\overset{\mathrm{}}{}}\left[\frac{1}{(n+ϵ_i)^2}+\frac{1}{(nϵ_i)^2}\right],$$
(45)
it is easy to show the normalization condition for the matrix elements (42) and (43),
$$(t_i^i)^2+\underset{n=1}{\overset{\mathrm{}}{}}(t_i^{in})^2+(t_i^{i+n})^2=1$$
(46)
and also the orthogonality relation,
$$\underset{r}{}t_i^rt_k^r=0(ik)$$
(47)
in the limit $`R\mathrm{}`$.
3.3 - The transformation matrix in the limit $`g=0`$
From Eq. (16) we get for arbitrary $`R`$,
$$Lim_{g0}t_0^r=\{\begin{array}{cc}1,\hfill & \text{if }\mathrm{\Omega }_r=\overline{\omega }\text{;}\hfill \\ 0,\hfill & \text{otherwise.}\hfill \end{array}.$$
(48)
From Eqs.(42) and (43) we see that the matrix elements $`t_i^r`$ for $`ir`$ all vanish for $`g=0`$. Also, using Eqs.(21) we obtain for small $`g`$,
$$t_i^i\frac{2g\mathrm{\Omega }_i\omega _i}{(\mathrm{\Omega }_i^2\overline{\omega }^2)(\omega _i+\mathrm{\Omega }_i)}\frac{1}{ϵ_i},$$
(49)
or, expanding $`ϵ_i`$ for small $`g`$ from Eq.(44)
$$t_i^i(g=0)=1$$
(50)
We see from the above expressions that in the limit $`R\mathrm{}`$ the matrix $`(t_\mu ^r)`$ remains an orthonormal matrix in the usual sense as for finite $`R`$. With the choice of the procedure of taking the limit $`R\mathrm{}`$ from the confined solution, the matrix elements do not tend to distributions in the free space limit as it would be the case using the continuous formalism. All non- vanishing matrix elements $`t_i^r`$ are concentrated inside a neighbourhood $`\eta `$ of $`\omega _i`$, their set is a quadratically summable enumerable set. The elements $`(t_0^r)`$ tend to a quadratically integrable expression.
## 4 The Radiation Process
We start this section defining some coordinates $`q_0^{}`$, $`q_i^{}`$ associated to the ”dressed” mechanical oscillator and to the field. These coordinates will reveal themselves to be suitable to give an appealling non-perturbative description of the oscillator-field system. The general conditions that such coordinates must satisfy, taking into account that the system is rigorously described by the collective normal coordinates modes $`Q_r`$, are the following:
\- In reason of the linear character of our problem the coordinates $`q_0^{}`$, $`q_i^{}`$ should be linear functions of the collective coordinates $`Q_r`$
\- They should allow to construct ortogonal configurations corresponding to the separation of the system into two parts, the dressed oscillator and the field.
\- The set of these configurations should contain the ground state, $`\mathrm{\Gamma }_0`$.
The last of the above conditions restricts the transformation between the coordinates $`q_\mu ^{}`$, $`\mu =0,i=1,2,\mathrm{}`$ and the collective ones $`Q_r`$ to those leaving invariant the quadratic form,
$$\underset{r}{}\mathrm{\Omega }_rQ_r^2=\overline{\omega }(q_0^{})^2+\underset{i}{}\omega _i(q_i^{})^2$$
(51)
Our configurations will behave in a first approximation as independent states, but they will evolve as the time goes on, as if transitions among them were being in progress, while the basic configuration $`\mathrm{\Gamma }_0`$ represents a rigorous eigenstate of the system and does not change with time. The new coordinates $`q_\mu ^{}`$ describe dressed configurations of the oscillator and field quanta.
4.1 - The dressed coordinates $`q_\mu ^{}`$
The eigenstates of our system are represented by the normalized eigenfunctions,
$$\varphi _{n_0n_1n_2\mathrm{}}(Q,t)=\underset{s}{}\left[N_{n_s}H_{n_s}(\sqrt{\frac{\mathrm{\Omega }_s}{\mathrm{}}}Q_s)\right]\mathrm{\Gamma }_0e^{i_sn_s\mathrm{\Omega }_st},$$
(52)
where $`H_{n_s}`$ is the $`n_s`$-th Hermite polynomial, $`N_{n_s}`$ is a normalization coefficient,
$$N_{n_s}=(2^{n_s}n_s!)^{\frac{1}{2}}$$
(53)
and $`\mathrm{\Gamma }_0`$ is a normalized representation of the ground state,
$$\mathrm{\Gamma }_0=exp\left[\underset{s}{}\frac{\mathrm{\Omega }_sQ_s^2}{2\mathrm{}}\frac{1}{4}ln\frac{\mathrm{\Omega }_s}{\pi \mathrm{}}\right].$$
(54)
To describe the radiation process, having as initial condition that only the mechanical oscillator, $`q_0`$ be excited, the usual procedure is to consider the interaction term in the Hamiltonian written in terms of $`q_0`$, $`q_i`$ as a perturbation, which induces transitions among the eigenstates of the free Hamiltonian. In this way it is possible to treat approximatelly the problem having as initial condition that only the bare oscillator be excited. But as is well known this initial condition is physically not consistent due to the divergence of the bare oscillator frequency if there is interaction with the field. The traditional way to circumvect this difficulty is by the renormalization procedure, introducing perturbativelly order by order corrections to the oscillator frequency. Here we adopt an alternative procedure, we do not make explicit use of the concepts of interacting bare oscillator and field, described by the coordinates $`q_0`$ and $`\{q_i\}`$, we introduce ”dressed” coordinates $`q_0^{}`$ and $`\{q_i^{}\}`$ for, respectivelly the ”dressed” oscillator and the field, defined by,
$$\sqrt{\frac{\overline{\omega }_\mu }{\mathrm{}}}q_\mu ^{}=\underset{r}{}t_\mu ^r\sqrt{\frac{\mathrm{\Omega }_r}{\mathrm{}}}Q_r,$$
(55)
valid for arbitrary $`R`$, which satisfy the condition to leave invariant the quadratic form (51) and where $`\overline{\omega }_\mu =\overline{\omega },\{\omega _i\}`$. In terms of the bare coordinates the dressed coordinates are expressed as,
$$q_\mu ^{}=\underset{\nu }{}\alpha _{\mu \nu }q_\nu ,$$
(56)
where
$$\alpha _{\mu \nu }=\frac{1}{\sqrt{\overline{\omega }_\mu }}\underset{r}{}t_\mu ^rt_\nu ^r\sqrt{\mathrm{\Omega }_r}.$$
(57)
As $`R`$ becomes larger and larger we get for the various coefficients $`\alpha `$ in Eq.(57):
a) from Eq.(38),
$$Lim_R\mathrm{}\alpha _{00}=\frac{1}{\sqrt{\overline{\omega }}}_0^{\mathrm{}}\frac{2g\mathrm{\Omega }^2\sqrt{\mathrm{\Omega }}d\mathrm{\Omega }}{(\mathrm{\Omega }^2\overline{\omega }^2)^2+\pi ^2g^2\mathrm{\Omega }^2}A_{00}(\overline{\omega },g).$$
(58)
b) To evaluate $`\alpha _{0i}`$ and $`\alpha _{0i}`$ in the limit $`R\mathrm{}`$, we remember from the discussion in subsection 3.2 that in the the limit $`R\mathrm{}`$, for each $`i`$ the only non-vanishing matrix elements $`t_i^r`$ are those for which the corresponding eigenfrequencies $`\mathrm{\Omega }_r`$ are arbitrarily near the field frequency $`\omega _i`$. We obtain from Eqs. (38), (42) and (43),
$$Lim_R\mathrm{}\alpha _{i0}=Lim_{\mathrm{\Delta }\omega 0}\frac{1}{\sqrt{\omega _i}}\frac{(2g^2\omega _i^5\mathrm{\Delta }\omega )^{\frac{1}{2}}}{(\omega _i^2\overline{\omega }^2)^2+\pi ^2g^2\omega _i^2}(\underset{n=1}{\overset{\mathrm{}}{}}\frac{2ϵ_i}{n^2ϵ_i^2}\frac{1}{ϵ_i})$$
(59)
and
$$Lim_R\mathrm{}\alpha _{0i}=Lim_{\mathrm{\Delta }\omega 0}\frac{1}{\sqrt{\overline{\omega }}}\frac{(2g^2\omega _i^5\mathrm{\Delta }\omega )^{\frac{1}{2}}}{(\omega _i^2\overline{\omega }^2)^2+\pi ^2g^2\omega _i^2}(\underset{n=1}{\overset{\mathrm{}}{}}\frac{2ϵ_i}{n^2ϵ_i^2}\frac{1}{ϵ_i})$$
(60)
c) Since in the limit $`R\mathrm{}`$ the only non-zero matrix elements $`t_i^r`$ corresponds to $`\mathrm{\Omega }_r=\omega _i`$, the product $`t_i^rt_k^r`$ vanishes for $`\omega _i\omega _k`$. Then we obtain from Eqs.(57) and (46)
$$Lim_R\mathrm{}\alpha _{ik}=\delta _{ik}.$$
(61)
Thus, from Eqs.(56), (61), (59), (60) and (58) we can express the dressed coordinates $`q_\mu ^{}`$ in terms of the bare ones, $`q_\mu `$ in the limit $`R\mathrm{}`$,
$$q_0^{}=A_{00}(\overline{\omega },g)q_0,$$
(62)
$$q_i^{}=q_i.$$
(63)
It is interesting to compare Eqs.(56) with Eqs.(62), (63). In the case of Eqs.(56) for finite $`R`$, the coordinates $`q_0^{}`$ and $`\{q_i^{}\}`$ are all dressed, in the sense that they are all collective, both the field modes and the mechanical oscillator can not be separeted in this language. In the limit $`R\mathrm{}`$, Eqs.(62) and (63) tells us that the coordinate $`q_0^{}`$ describes the mechanical oscillator modified by the presence of the field in a indissoluble way, the mechanical oscillator is always dressed by the field. On the other side, the dressed harmonic modes of the field, described by the coordinates $`q_i^{}`$ are identical to the bare field modes, in other words, the field keeps in the limit $`R\mathrm{}`$ its proper identity, while the mechanical oscillator is always accompanied by a cloud of field quanta. Therefore we identify the coordinate $`q_0^{}`$ as the coordinate describing the mechanical oscillator dressed by its proper field, being the whole system divided into dressed oscillator and field, without appeal to the concept of interaction between them, the interaction being absorbed in the dressing cloud of the oscillator. In the next subsections we use the dressed coordinates to describe the radiation process.
4.2 - Dressed configurations and the radiation process
Let us define for a fixed instant the complete orthonormal set of functions,
$$\psi _{\kappa _0\kappa _1\mathrm{}}(q^{})=\underset{\mu }{}\left[N_{\kappa _\mu }H_{\kappa _\mu }(\sqrt{\frac{\overline{\omega }_\mu }{\mathrm{}}}q_\mu ^{})\right]\mathrm{\Gamma }_0,$$
(64)
where $`q_\mu ^{}=q_0^{},q_i^{}`$, $`\overline{\omega }_\mu =\overline{\omega },\omega _i`$ and $`N_{\kappa _\mu }`$ and $`\mathrm{\Gamma }_0`$ are as in Eq.(52). Using Eq.(55) the functions (64) can be expressed in terms of the normal coordinates $`Q_r`$. But since (52) is a complete set of orthonormal functions, the functions (64) may be written as linear combinations of the eigenfunctions of the coupled system (we take $`t=0`$ for the moment),
$$\psi _{\kappa _0\kappa _1\mathrm{}}(q^{})=\underset{n_0n_1\mathrm{}}{}T_{\kappa _0\kappa _1\mathrm{}}^{n_0n_1\mathrm{}}(0)\varphi _{n_0n_1n_2\mathrm{}}(Q,0),$$
(65)
where the coefficients are given by,
$$T_{\kappa _0\kappa _1\mathrm{}}^{n_0n_1\mathrm{}}(0)=𝑑Q\psi _{\kappa _0\kappa _1\mathrm{}}\varphi _{n_0n_1n_2\mathrm{}},$$
(66)
the integral extending over the whole $`Q`$-space.
We consider the particular configuration $`\psi `$ in which only one dressed oscillator $`q_\mu ^{}`$ is in its $`N`$-th excited state,
$$\psi _{0\mathrm{}N(\mu )0\mathrm{}}(q^{})=N_NH_N(\sqrt{\frac{\overline{\omega }_\mu }{\mathrm{}}}q_\mu ^{})\mathrm{\Gamma }_0.$$
(67)
The coefficients (66) can be calculated in this case using Eqs.(66), (64) and (55) with the help of the theorem ,
$$\frac{1}{m!}\left[\underset{r}{}(t_\mu ^r)^2\right]^{\frac{m}{2}}H_N(\frac{_rt_\mu ^r\sqrt{\frac{\mathrm{\Omega }_r}{\mathrm{}}}Q_r}{\sqrt{_r(t_\mu ^r)^2}})=\underset{m_0+m_1+\mathrm{}=N}{}\frac{(t_\mu ^0)^{m_0}(t_\mu ^1)^{m_1}\mathrm{}}{m_0!m_1!\mathrm{}}H_{m_0}(\sqrt{\frac{\mathrm{\Omega }_0}{\mathrm{}}}Q_0)H_{m_1}(\sqrt{\frac{\mathrm{\Omega }_1}{\mathrm{}}}Q_1)\mathrm{}$$
(68)
We get,
$$T_{0\mathrm{}N(\mu )0\mathrm{}}^{n_0n_1\mathrm{}}=(\frac{m!}{n_0!n_1!\mathrm{}})^{\frac{1}{2}}(t_\mu ^0)^{n_0}(t_\mu ^1)^{n_1}\mathrm{},$$
(69)
where the subscripts $`\mu =0,i`$ refer respectivelly to the dressed mechanical oscillator and the harmonic modes of the field and the quantum numbers are submited to the constraint $`n_0+n_1+\mathrm{}=N`$.
In the following we study the behaviour of the system with the initial condition that only the dressed mechanical oscillator $`q_0^{}`$ be in the $`N`$-th excited state. We will study in detail the particular cases $`N=1`$ and $`N=2`$, which will be enough to have a clear understanding of our approach.
\- $`N=1`$: Let us call $`\mathrm{\Gamma }_1^\mu `$ the configuration in which only the dressed oscillator $`q_\mu ^{}`$ is in the first excited level. The initial configuration in which the dressed mechanical oscillator is in the first excited level is $`\mathrm{\Gamma }_1^0`$. We have from Eq.(67), (65) (69) and (55) the following expression for the time evolution of the first-level excited dressed oscillator $`q_\mu ^{}`$,
$$\mathrm{\Gamma }_1^\mu =\underset{\nu }{}f^{\mu \nu }(t)\mathrm{\Gamma }_1^\nu (0),$$
(70)
where the coefficients $`f^{\mu \nu }(t)`$ are given by
$$f^{\mu \nu }(t)=\underset{s}{}t_\mu ^st_\nu ^se^{i\mathrm{\Omega }_st},$$
(71)
That is, the initially excited dressed oscillator naturally distributes its energy among itself and all others dressed oscillators, as time goes on. If the mechanical dressed oscillator is in its first excited state at $`t=0`$, its decay rate may evaluated from its time evolution equation,
$$\mathrm{\Gamma }_1^0=\underset{\nu }{}f^{0\nu }(t)\mathrm{\Gamma }_1^\nu (0).$$
(72)
In Eq.(72) the coefficients $`f^{0\nu }(t)`$ have a simple interpretation: remembering Eqs.(62) and (63), $`f^{00}(t)`$ and $`f^{0i}(t)`$ are respectivelly the probability amplitudes that at time $`t`$ the dressed mechanical oscillator still be excited or have radiated a field quantum of frequency $`\omega _i`$. We see that this formalism allows a quite natural description of the radiation process as a simple exact time evolution of the system. Let us for instance evaluate the oscillator decay probability in this language. From Eqs.(38) and (71) we get
$$f^{00}(t)=_0^{\mathrm{}}\frac{2g\mathrm{\Omega }^2e^{i\mathrm{\Omega }t}d\mathrm{\Omega }}{(\mathrm{\Omega }^2\omega ^2)^2+\pi ^2g^2\mathrm{\Omega }^2}.$$
(73)
The above integral can be evaluated by Cauchy theorem. For large $`t`$ ($`t>>\frac{1}{\overline{\omega }}`$), but arbitrary coupling $`g`$, we obtain for the oscillator decay probability, the result,
$$|f^{00}(t)|^2=e^{\pi gt}(1+\frac{\pi ^2g^2}{4\overline{\omega }^2})+e^{\pi gt}\frac{8\pi g}{\pi \overline{\omega }^4t^3}(sin\stackrel{~}{\overline{\omega }}t+\frac{\pi g}{2<\overline{\omega }>}cos\stackrel{~}{\overline{\omega }}t)+\frac{16\pi ^2g^2}{\pi ^2\overline{\omega }^8t^6},$$
(74)
where $`\stackrel{~}{\overline{\omega }}=\sqrt{\overline{\omega }^2\frac{\pi ^2g^2}{4}}`$. In the above expression the approximation $`t>>\frac{1}{\overline{\omega }}`$ plays a role only in the two last terms, due to the difficulties to evaluate exactly the integral in Eq. (73) along the imaginary axis. The first term comes from the residue at $`\mathrm{\Omega }=\stackrel{~}{\overline{\omega }}+i\frac{\pi g}{2}`$ and would be the same if we have done an exact calculation. If we consider the case of weak coupling, $`g<<\overline{\omega }`$, we obtain the well known perturbative exponential decay law for the harmonic oscillator,
$$|f^{00}(t)|^2e^{\pi gt},$$
(75)
but we emphasize that Eq.(74) is valid for all values of the coupling constant $`g`$, even large, it is an expression valid for weak as well as strong couplings.
\- $`N=2`$
Let us call $`\mathrm{\Gamma }_{11}^{\mu \nu }`$ the configuration in which the dressed oscillators $`q_\mu ^{}`$ and $`q_\nu ^{}`$ are at their first excited level and $`\mathrm{\Gamma }_2^\mu `$ the configuration in which $`q_\mu ^{}`$ is at its second excited level. Taking as initial condition that the dressed mechanical oscillator be at the second excited level, the time evolution of the state $`\mathrm{\Gamma }_2^0`$ may be obtained in an analogous way as in the preceeding case,
$$\mathrm{\Gamma }_2^0(t)=\underset{\mu }{}\left[f^{\mu \mu }(t)\right]^2\mathrm{\Gamma }_2^\mu +\frac{1}{\sqrt{2}}\underset{\mu \nu }{}f^{0\mu }(t)f^{0\nu }(t)\mathrm{\Gamma }_{11}^{\mu \nu },$$
(76)
where the coefficients $`f^{\mu \mu }`$ and $`f^{0\mu }`$ are given by (71). Then it easy to get the following probabilities:
Probability that the dressed oscillator still be excited at time $`t`$:
$$P_0(t)=|f^{00}(t)|^4,$$
(77)
probability that the dressed oscillator have decayed at time $`t`$ to the first level by emission of a field quantum:
$$P_1(t)=2|f^{00}(t)|^2(1|f^{00}(t)|^2)$$
(78)
and probability that the dressed oscillator have decayed at time $`t`$ to the ground state:
$$P_2(t)=12|f^{00}(t)|^2+|f^{00}(t)|^4.$$
(79)
Replacing Eq.(74) in the above expressions we get expressions for the probabilities decays valid for any value of the coupling constant. In the particular case of weak coupling we obtain the well known perturbative formulas for the oscillator decay ,
$$P_0(t)e^{2\pi gt},$$
(80)
$$P_1(t)2e^{\pi gt}(1e^{\pi gt})$$
(81)
and
$$P_2(t)12e^{\pi gt}+e^{2\pi gt}.$$
(82)
## 5 Concluding Remarks
In this paper we have analysed a symplified version of an atom-electromagnetic field system and we have tried to give the more exact and rigorous treatment we could to the problem. We have adopted a general physicist’ s point of view, in the sense that we have rennounced to approach very closely to the real behaviour of a complicated non-linear system, to study instead a simple linear model. As a counterpart, an exact solution has been possible. Our dressed coordinates give a description of the behaviour of the system that is exact and valid for weak as well as for strong coupling. If the coupling between the mechanical oscillator and the field is weak, we recover the well known behaviour from perturbation theory.
## 6 In Memoriam
This paper evolved from umpublished work we have done and discussions we have had, with Prof. Guido Beck when two of us (A.P.C.M. and N.P.A.) were his students at Instituto de Fisica Balseiro in Bariloche (Argentina), in the late sixties and the early seventies. We dedicate this article to his memory.
## 7 Acknowlegements
This paper was supported by Conselho Nacional de Desenvolvimento Cientifico e Tecnologico (CNPq) - Brazil. |
warning/0001/nlin0001069.html | ar5iv | text | # Asymptotic Dynamics of Ripples 11footnote 1 PACS # 03.40.Kf, 47.35.+i E-mail: manna@lpm.univ-montp2.fr
## 1 Introduction
Since the classical works of Boussinesq and Korteweg and de Vries, nonlinear evolution of long waves of small amplitude in shallow fluids has been widely studied. The asymptotic dynamics (in space and time) is by now well understood and represented by a large number of model equations, among which the different versions of the Boussinesq equations , the Korteweg-de Vries (KdV), the Benjamin-Bona-Mahoney-Peregrini (BBMP) , and the Camassa-Holm equations.
On the contrary, short waves have been studied very little and only a few results are known on their asymptotic dynamics. The main purpose of this paper is to study nonlinear short surface waves in fluids, the ripples, shown to build up as a result of superposition of two surface motions: an oscillatory flow and a laminar flow. The oscillatory flow corresponds to mechanical perturbations which propagate like a wave. The laminar flow may be created in many ways: by the action of an external wind, or in a two-layer liquid where the upper one displaces particles belonging to the upper surface of the lower fluid, or by an external electric field if the surface particles are charged, etc…
Here we study an ideal fluid (inviscid, incompressible and without surface tension) in which surface displacement can be achieved through the action of a steady external wind directed parallel to the water surface. A surface wind on a lake which produces surface flow is an ideal physical environment conducive to the above mentioned phenomenon.
## 2 Model equation for ripples propagation
Let’s consider $`u(x,t)`$ which represents at time $`t`$ an unidirectional surface wave propagation in the $`x`$ direction of a fluid medium which is involved in a flow on a large scale. This large scale flow is a superposed motion of the surface which moves under the action of the wind with a velocity $`c_0`$ in relation to the bulk. Using perturbative methods we will show that ripples can propagate and obey the nonlinear equation
$$u_{xt}=\frac{3g}{hc_0}uuu_{xx}+(u_x)^2$$
(2.1)
where subscripts denote partial derivatives, $`h`$ is the unperturbed initial depth and $`g`$ is the acceleration of gravity. The equation (2.1) has several types of interesting solutions. The peakon solution is (using the formula $`(1_{zz})e^{|z|}=2\delta (z)`$, )
$$u=\alpha \lambda ^2\mathrm{exp}(|\frac{x+\alpha \lambda ^2t}{\lambda }|),\alpha =\frac{3g}{hc_0},$$
(2.2)
where the width $`\lambda `$ is a free parameter. Unlike the peakon of Camassa-Holm equation, the amplitude $`(\alpha \lambda ^2)`$, the velocity $`(\alpha \lambda ^2)`$ and the width $`(\lambda )`$ are interrelated.
The static compacton solution is
$$u=8\alpha \lambda ^2\mathrm{cos}^2(\frac{x}{4\lambda }),|\frac{x}{\lambda }|2\pi ,$$
(2.3)
and $`u=0`$ otherwise. Unlike the compacton solution recently introduced and investigated by Roseneau and Hyman , this solution presents a dependence between width and amplitude. Solutions (2.2) is a coherent struture analogous to the solitary wave of KdV. Moreover, contrary to KdV, (2.1) possess a plane monochromatic wave solution of arbitrary amplitude $`A`$
$$u=A\mathrm{exp}i(kx\mathrm{\Omega }t)$$
(2.4)
with the dispersion relation $`\mathrm{\Omega }(k)=3g/khc_0`$ identical to the one of the linearized system . Moreover this dispersion relation is a function of $`k`$ which has a good behavior in the short wave limit.
Despite the fact that solution (2.3) is unnatural in the physics under consideration, it shows that (2.1) is an adequate mathematical tool to modelize statics patterns in nature.
The underlying mechanism responsible for structural stability of solutions of the Camassa-Holm equation or the KdV equation with nonlinear dispersion is the balance between nonlinear dispersion, nonlinear convection and nonlinearity. Indeed these equations are nonlinear evolution equations without linear dispersion (the plane wave is not a solution of the linear associated evolution equations). However the exibited solutions of (2.1) come from the balance between linear dispersion and two nonlinear terms (in the case of (2.2)) and only nonlinearity (in the case of (2.3)). An essential point, which remains to be proved by numerical or analytical methods is: under what conditions do the solution (2.2) dominates the initial value problem of equation (2.1)?. Another important open problem is the structural stability of (2.2) and (2.3) . Here we only show that the plane wave (2.4) is unstable.
## 3 Physical context: modified Green-Naghdi system of equations.
We consider an inviscid, incompressible, homogeneous fluid with density $`\sigma `$. Let the particles of this continuum medium be identified by a fixed rectangular Cartesian system of center $`O`$ and axes $`(x_1,x_2,x_3)=(x,z,y)`$ with $`Oz`$ the upward vertical direction. We assume symmetry in $`y`$ and we will only consider a sheet of fluid in the $`xz`$ plane. This fluid sheet is moving in a domain with a rigid botton at $`z=0`$ and an upper free surface at $`z=\varphi (x,t)`$. The vector velocity is $`\stackrel{}{v}=(v_1,v_2)=(u,w)`$ for $`0z<\varphi (x,t)`$. Thanks to homogenity and incompressibility the continuity equation reduces to
$$\stackrel{}{}\stackrel{}{v}=u_x(x,z,t)+w_z(x,z,t)=0,\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}0}z\varphi (x,t),$$
(3.1)
where $`\stackrel{}{}=(_x,_z)`$. The Euler equations of motions (law of conservation of momentun) of a fluid under gravity $`g`$ and for $`0z\varphi (x,t)`$ are
$$\sigma \dot{u}(x,z,t)=p_x^{}(x,z,t),$$
(3.2)
$$\sigma \dot{w}(x,z,t)=p_z^{}(x,z,t)g\sigma ,$$
(3.3)
where $`p^{}(x,z,t)`$ is the pressure and a superposed dot denotes the material derivative for $`x,z`$ fixed.
We complete now the fundamental equations of continuity and momentun conservation with appropiate kinematic and dynamics boundaries conditions. Let $`S(x,z,t)`$ be the interface between the inviscid fluid sheet and the air (external medium). We represented $`S(x,z,t)`$ by the (classical) equation
$$S(x,z,t)=\varphi (x,t)z=0.$$
(3.4)
The kinematic condition is that the normal velocity of the surface $`S(x,z,t)`$ must be equal to the velocity of the fluid sheet normal to the surface. The normal velocity of the surface is
$$\frac{S_t}{\stackrel{}{}S},$$
(3.5)
while the velocity on the surface $`z=\varphi (x,t)`$ is
$$\stackrel{}{v}=(u+c_0,w),z=\varphi (x,t),$$
(3.6)
whose normal component is
$$\stackrel{}{v}\frac{\stackrel{}{}S}{\stackrel{}{}S}.$$
(3.7)
From (3.5), (3.6) and (3.7) the kinematic conditions read
$$\varphi _t+u\varphi _xw+c_0\varphi _x=0,z=\varphi .$$
(3.8)
The equation (3.6) leading to (3.8) lies at the heart of our approach. It points out that an external agent - wind in our case - drives the particles belonging to $`S(x,z,t)`$. The motion is uniform, of value $`c_0`$ and in the $`x`$ direction only. So $`S(x,z,t)`$ is a surface of discontinuity for the $`u`$ component of $`\stackrel{}{v}`$, which experiences finite jumps of value $`c_0`$ in $`z=\varphi `$. This wind induced motion of the fluid’s surface is found in particular in the hydrodynamics of lakes . From a theoretical point of view this phenomenon was studied in the ray tracing theory by Lighthill in where an equation analogous to (3.8) was derived.
At the surface of the sheet $`z=\varphi `$, there is a constant normal pressure $`p_0`$. At the bed $`z=0`$, there is an unknown pressure $`p^{}(x,0,t)`$ and the normal fluid velocity is zero: $`w=0`$. Several types of approximations can be used in order to solve this wave problem. Here we adopt an approximation in the velocity field introduced by Green, Laws and Naghdi in and based on the monumental work of Naghdi . We assume that $`u`$ is independent of $`z`$. This is equivalent to considering the vertival component $`w`$ as a linear function of $`z`$. This simple and realistic assumption enables us to satisfy exactly the equation of incompressibility and the boundary condition at the bed. Hence $`u=u(x,t)`$ and from (3.1) we have
$$w=z\xi (x,t),\xi (x,t)=u_x(x,t).$$
(3.9)
Now we integrate (3.2) in the variable $`z`$ from $`z=0`$ to $`z=\varphi `$ (the integration is granted by the Riemann’s condition of integrability). The result is
$$\sigma (u_t+uu_x)\varphi =p_x,$$
(3.10)
where we use $`\dot{u}=u_t+uu_x`$ and
$$p(x,t)=_0^{\varphi (x,t)}p^{}(x,z,t)𝑑zp_0\varphi (x,t).$$
(3.11)
Next, we multiply equation (3.3) by $`z`$ and integrate from $`z=0`$ to $`z=\varphi `$ which yields
$$\sigma (\xi ^2+\dot{\xi })\frac{\varphi ^3}{3}+\sigma g\frac{\varphi ^2}{2}=p.$$
(3.12)
The pressure $`p`$ can be eliminated using (3.12) in (3.10) and, with the help of (3.9) we eventually obtain
$`u_t+uu_x+g\varphi _x=`$ $`\varphi \varphi _x(u_{xt}+uu_{xx}(u_x)^2)+`$ (3.13)
$`{\displaystyle \frac{\varphi ^3}{3}}(u_{xt}+uu_{xx}(u_x)^2)_x`$
The remaining upper boundary condition on $`\varphi `$ reads using (3.9)
$$\varphi _t+(\varphi u)_x+c_0\varphi _x=0.$$
(3.14)
With $`c_0=0`$, the system (3.13) (3.14) is the Green-Naghdi system of equations . The inclusion of the term $`c_0\varphi _x`$ drastically changes its dynamics: $`c_0`$ cannot be eliminated neither by a Galilean transformation nor by a rescaling of $`u`$ or $`\varphi `$. These extended Green-Naghdi equations represent the nonlinear interaction between two separate forms of motion: a wave motion associated with the elastic response of the fluid to a perturbation, a surface (uniform) motion generated by an external agent.
## 4 Nonlinear dynamics of ripples.
Let us consider now asymptotic nonlinear dynamics of ripples in (3.13) and (3.14). Ripples are, from a geometrical point of view, nearly local objects and we are looking for their large time behavior. So we need to introduce two appropiate variables: one space variable $`\zeta `$ describing a local and asymptotic pattern and a time variable $`\tau `$ measuring asymptotic time dynamics. These asymptotic are worked out by means of the change of variables
$$\zeta =\frac{1}{ϵ}x,\tau =ϵt.$$
(4.1)
where the small parameter $`ϵ`$ is related to the size of the wavelength: $`\mathrm{}=2\pi /kϵ`$.
Such variables cannot always be defined as shown in , but exist if the linear dispersion relation $`\mathrm{\Omega }(k)`$ can be expanded as a Laurent series with a simple pole for $`k\mathrm{}`$, that is
$$\mathrm{\Omega }(k)=Ak+\frac{B}{k}+\frac{C}{k^3}+\mathrm{}.k\mathrm{}.$$
where $`A`$, $`B`$, $`C`$, $`\mathrm{},`$ are constants. Consequently the phase and group velocities remain bounded in the short wave limit $`k\mathrm{}`$. This is the case here because the dispersion relation for (3.13) (3.14), obtained for $`\varphi (x,t)=h+a\mathrm{exp}i[kx\mathrm{\Omega }(k)t]`$ and $`u(x,t)=b\mathrm{exp}i[kx\mathrm{\Omega }(k)t]`$, has the asymptotic expansion
$$\mathrm{\Omega }(k)=(\frac{3g}{c_0h})\frac{1}{k}+\frac{3}{c_0h^2}(\frac{3g}{h}+\frac{3g^2}{c_0^2})\frac{1}{k^3}+𝒪(\frac{1}{k^5}).$$
(4.2)
Such an expansion justifies the change of variables (4.1) and it is essential in revealing in (3.13) (3.14) the asymptotic short wave dynamics . We are looking now for nonlinear dynamics of ripples of small amplitude represented by the expansions
$`\varphi =h+ϵ^2(H_0+ϵ^2H_2+ϵ^4H_4+\mathrm{}),`$ (4.3)
$`u=ϵ^2(U_0+ϵ^2U_2+ϵ^4U_4+\mathrm{}),`$ (4.4)
which from (3.13) and (3.14) give at lower order in $`ϵ`$:
$`h(U_0)_\zeta +c_0(H_0)_\zeta =0,`$ (4.6)
$`g(H_0)_\zeta ={\displaystyle \frac{h^2}{3}}\{(U_0)_{\tau \zeta }+U_0(U_0)_{\zeta \zeta }(U_0)_\zeta ^2\}_\zeta .`$
Using (4.6) in (4.6), going back to the laboratory fields ($`\varphi ,u`$) and co-ordinates ($`x,t`$) and integrating once we eventually arrive at equation (2.1).
## 5 Benjamin-Feir instability of the plane wave.
For a plane wave, the nonlinears terms in (2.1) give rise to harmonics of the fundamental. Assume that a disturbance is present consisting of modes with sideband frequences and wavenumbers close to the fundamental. We can have interaction between harmonics and these sideband modes. This interaction is likely to produce a resonant phenomenon manifesting itself by the modulation of the plane wave solution. The exponencial growth in time of the modulation, originating from synchronous resonance between harmonics and sideband modes, leads to the Benjamin-Feir instability. A formal solutions can be given via an asymptotic expansion conducing to the nonlinear Schodinger equation (NLS). The particular interest of NLS is the existence of a general and simple criterion enable to detect stability or unstability of the monochromatic wave train. Let us seek for a solution of (2.1) under the form of a Fourier expansion in harmonics of the fundamental $`\mathrm{exp}i(kx\mathrm{\Omega }t)`$ and where the Fourier components are developed in a Taylor serie in powers of a small parameter $`\gamma `$ mesuring the amplitude of the fundamental
$$u=\underset{l=p}{\overset{l=p}{}}\underset{p=1}{\overset{\mathrm{}}{}}\mathrm{exp}il(kx\omega t)\gamma ^pu_l^p(\xi ,\tau )$$
(5.1)
In (5.1), $`u_l^p=u_l^p`$ (”star” denotes complex conjugation) and $`\xi `$ and $`\tau `$ are slow variables introduced through the stretching $`\xi =\gamma (xvt)`$ and $`\tau =\gamma ^2t`$ and where $`v`$ will be determined as a solvability condition. The expansions (5.1) includes fast local oscillations through the dependence on the harmonics and slow variation (modulation) in amplitude taken into account by the $`\xi `$, $`\tau `$ dependence of $`u_l^p`$. Introducing now this expansion and the slow variables in (2.1) we may proceed to collect and solve different order $`\gamma `$ and $`l`$. We obtain: $`u_0^1=u_2^2=u_2^3=u_3^3=0`$, $`v=\gamma /k^2`$, $`u_1^1=A(\xi ,\tau )0`$ and $`u_0^2=4k^3/\alpha `$. At order $`\gamma =3`$, $`l=1`$ we obtain NLS for $`A(\xi ,\tau )`$
$$iA_\tau +\frac{\alpha }{k^3}A_{\xi \xi }+\frac{4k^3}{\alpha }A|A|^2=0.$$
(5.2)
The nature of solutions of NLS depends drastically of the sign of the product between the coeficient of $`A_{\xi \xi }`$ and that of $`A|A|^2`$. In this case the product of $`\alpha /k^3`$ by $`4k^3/\alpha `$ is positive, and according to a well known stability criterion (see for exemple ) the plane wave solution (2.4) is unstable.
## 6 Conclusion and final coments on short waves.
We have shown that ripples result from balance between linear dispersion and nonlinearity as in the case of long waves. However the physical model under consideration represents an ideal fluid because viscosity has been neglected, and, since dissipative phenomena take place at small scales, viscosity must affect asymptotic dynamics of ripples.
The same approach can be used to study a flow of multiple layers. Such a flow occurs in many practical applications (e.g multilayer coating) in which the short wave behavior is of primordial importance.
Note that, while the original Green-Naghdi system is Galilean invariant, this invariance is lost in the Camassa-Holm equation derived by Hamiltonian methods. In our case the perturbative theory conserves Galilean invariance and consequently equation (2.1) inherits this property from modified Green-Naghdi.
In nonlinear dispersive systems, short wave dynamics does not occur as naturally as long wave dynamics. It is interesting to point out that some intermediate long wave shallow water models behave well, paradoxically, in the short wave limit. In these intermediate models a second asymptotic limit is always possible, in general of long wave type, leading then to the ubiquitous KdV equation. In some cases however, there exists a second asymptotic limit of short wave type, leading to new nonlinear evolution equations. This is the case for BBMP, the integrable Camassa-Holm equation, and one of the Boussinesq system. The short wave limit of BBMP reads : $`u_{xt}=u3u^2`$, where $`u`$ has the same meaning as in (2.1). Its solitary wave solution comes from the balance between dispersion and nonlinearity as in KdV. Explicit solutions, blow-up and nonlinear instabilities are studied in . The Camassa-Holm equation has as short wave limit, an integrable equation which belongs to the Harry-Dym hierarchy ($`\kappa `$ is a constant): $`u_{xt}=\kappa u\frac{1}{2}(u_x)^2uu_{xx}`$. Finally the short wave limit of the Boussinesq system results as : $`u_{xtt}=u_tuu_x`$. The linear limit for all these cases is the equation $`u_{xt}=au`$ ($`a`$ constant) or else
$$u_ta_{\mathrm{}}^xu𝑑x=0.$$
(6.1)
This linear nonlocal equation for $`u`$ corresponds to the linear local wave equation $`u_tau_x=0`$ that appears in the long wave case. Note that the short wave limits of BBMP, Camassa-Holm and Boussinesq system are not evolution equations stricto sensu, rather being integro-differential equations. Thus the nonlocality of nonlinear evolution equations for short waves appears to reflect the basic non locality of (6.1).
### Aknowledgements
The author wish to thank J. Leon and G. Mennessier for many helpful and stimulating discussion. |
warning/0001/gr-qc0001095.html | ar5iv | text | # A Metric Theory of Gravity with Condensed Matter Interpretation
## 1 Introduction
General relativity is a very beautiful and successful theory of gravity. Nonetheless, the consideration of alternative theories of gravity remains to be legitimate part of science. Even if the purpose is only “to play devil’s advocate”, as in the case of some interesting theories of gravity (Lightman and Lee , Ni ) or to answer Rosen’s question “whether one can set up a theory of gravitation which will give agreement with observation without permitting black holes”.
The current research was motivated by the conceptual conflict between general relativity and quantum theory. The basic idea was that the Einstein equations may appear in a completely different metaphysical framework, which is better compatible with quantum principles than relativistic spacetime. This idea is in itself in full agreement with “the present educated view on the standard model, and of general relativity, … that these are leading terms in effective field theories” . The simplest choice would be a classical Newtonian framework with absolute time. This is a known way to solve the “problem of time” in quantum gravity, usually rejected for metaphysical reasons: “… in quantum gravity, one response to the problem of time is to ‘blame’ it on general relativity’s allowing arbitrary foliations of spacetime; and then to postulate a preferred frame of spacetime with respect to which quantum theory should be written. Most general relativists feel this response is too radical to countenance: they regard foliation-independence as an undeniable insight of relativity.” .
Nonetheless, following this “too radical to countenance” way, we have found a surprisingly simple and beautiful scheme which allows to derive a variant of the Einstein equations based not only on the classical Newtonian framework, but also on classical condensed matter theory – in other words, an ideal realization of the last century “ether” concept. In this derivation, we do not need any conspiracy to explain the Einstein equivalence principle. All we need are classical conservation laws and their connection with the Lagrange formalism. The point is the combination of the symmetry of the Lagrange formalism (self-adjoint equations) with the special character of the conservation laws (their relation with the preferred coordinates).
The mere existence of a viable theory of gravity with preferred frame is of great importance for other foundational problems. A preferred frame is, for example, required for compatibility of the EPR criterion of reality with the violation of Bell’s inequality . Bell himself concludes : “the cheapest resolution is something like going back to relativity as it was before Einstein, when people like Lorentz and Poincare thought that there was an aether — a preferred frame of reference — but that our measuring instruments were distorted by motion in such a way that we could no detect motion through the aether.” The theory presented here is strong support for this “cheapest resolution”. A closely related question is the extension of Bohmian mechanics into the domain of relativistic gravity which requires a preferred frame too.
The resulting theory differs from general relativity in an interesting way. It contains additional terms which depend on the preferred frame. These additional terms allow the definition of local energy and momentum densities of the gravitational field. But they don’t violate the Einstein equivalence principle – the theory remains to be a metric theory of gravity. They influence only the gravitational field itself, similar to dark matter.
The close analogy between condensed matter theory and gravity is well-known. It has been recognized that “effective gravity, as a low-frequency phenomenon, arises in many condensed matter systems” . This has been used to study Hawking radiation and the Unruh effect and vacuum energy for condensed matter examples. Wilczek mentions the general exchange of ideas with high energy physics, which “includes global and local spontaneous symmetry breaking, the renormalization group, effective field theory, solitons, instantons, and fractional charge and statistics”. This analogy has suggested the idea of a “Planck ether” . Our theory fits very well into this general context, and suggests interesting modifications: the critical length should not be Planck length.
## 2 The Theory
Our theory describes a classical medium in a Newtonian framework – Euclidean space and absolute time. The medium is described by steps of freedom typical for condensed matter theory. The gravitational field is defined by a positive density $`\rho `$, a velocity $`v^i`$, and a negative-definite symmetrical tensor field $`p^{ij}`$ which we name “pressure”. The effective metric $`g_{\mu \nu }`$ is defined algebraically by
$`\widehat{g}^{00}=g^{00}\sqrt{g}`$ $`=`$ $`\rho `$
$`\widehat{g}^{i0}=g^{i0}\sqrt{g}`$ $`=`$ $`\rho v^i`$
$`\widehat{g}^{ij}=g^{ij}\sqrt{g}`$ $`=`$ $`\rho v^iv^j+p^{ij}`$
This decomposition of $`g^{\mu \nu }`$ into $`\rho `$, $`v^i`$ and $`p^{ij}`$ is a variant of the ADM decomposition. The signature of $`g^{\mu \nu }`$ follows from $`\rho >0`$ and negative definiteness of $`p^{ij}`$.
The theory does not specify all properties of the medium, but only a few general properties – the conservation laws and their relation to the Lagrange formalism. The “material properties” of the medium, denoted by $`\phi ^m`$, remain unspecified. They become the matter fields. The complete specification of the medium – which includes the material laws of the medium – gives the theory of everything. The few general properties fixed here define a theory of gravity similar to GR. While it leaves the matter steps of freedom and the matter Lagrangian unspecified, it derives the Einstein equivalence principle.
For the derivation of the Lagrange formalism we prefer a formalism where the non-covariant terms are disguised as covariant, with the preferred coordinates considered formally as scalar fields $`X^\mu (x)`$. <sup>1</sup><sup>1</sup>1It is well-known that every physical theory may be described in a covariant way. But usually this is done in another way (for example, by Fock for SR): a flat background metric $`\gamma _{\mu \nu }`$ is described by vanishing curvature $`R_{\nu \kappa \lambda }^\mu =0`$. It is easy to transform a non-covariant Lagrangian $`L=L(T_{\mathrm{}}^{\mathrm{}},_\mu T_{\mathrm{}}^{\mathrm{}})`$ into a (formally) covariant form $`L=L(T_{\mathrm{}}^{\mathrm{}},_\mu T_{\mathrm{}}^{\mathrm{}},X_{,\nu }^\mu )`$. For example, the non-covariant component $`a^0`$ of a vector $`a^\mu `$ will be written as $`a^\nu X_{,\nu }^0`$. Here $`X^0=T`$ is no longer a spatial index, but enumerates one of the four “scalar fields” $`T,X,Y,Z`$.
This formalism is interesting in itself. Especially the question how to distinguish “truly covariant” theories from theories made covariant using this formalism is a very interesting one. The most interesting point of the formalism is the relation between conservation laws and the preferred coordinates. The conservation laws may be defined as the Euler-Lagrange equations for the preferred coordinates. The related energy-momentum tensor
$$T_\mu ^\nu =\frac{L}{X_{,\nu }^\mu }$$
is not the same as in Noether’s theorem, but is equivalent: if the Lagrangian does not depend on the $`X^\mu `$ them-self, we obtain immediately conservation laws in the form
$$_\nu T_\mu ^\nu =0$$
Now, the main postulate of the theory is that these conservation laws are identified with the classical conservation laws we know from condensed matter theory. First, the Euler-Lagrange equation for the preferred time T we identify with the classical continuity equation for the medium:
$$_t\rho +_i(\rho v^i)=0$$
(1)
The equations for the preferred spatial coordinates $`X^i`$ we identify with the Euler equation:
$$_t(\rho v^j)+_i(\rho v^iv^j+p^{ij})=0$$
(2)
Note that the Euler equation contains an important physical assumption: there is no momentum exchange with other materials, because there are no other materials. We have only one, universal, medium. All usual “matter fields” $`\phi ^m`$ are material properties of this universal medium.
The four conservation laws transform into the harmonic condition for the metric $`g_{\mu \nu }`$. Thus, they really look like equations for the preferred coordinates:
$$\mathrm{}X^\nu =_\mu (g^{\mu \nu }\sqrt{g})=0$$
Therefore, the main postulate transforms into the following relation between the the Euler-Lagrange equations for $`S=L`$ and the preferred coordinates $`X^\mu `$:
$$\frac{\delta S}{\delta X^\mu }(4\pi G)^1\gamma _{\mu \nu }\mathrm{}X^\nu $$
We have introduced here a constant diagonal matrix $`\gamma _{\mu \nu }`$ and a common factor $`(4\pi G)^1`$ to obtain appropriate units below. Euclidean symmetry gives $`\gamma _{11}=\gamma _{22}=\gamma _{33}`$. Thus, we have two coefficients $`\gamma _{00}=\mathrm{{\rm Y}},\gamma _{ii}=\mathrm{\Xi }`$. Now, we can derive the general form of the Lagrangian. First, we have the particular Lagrangian
$$L_0=(8\pi G)^1\gamma _{\mu \nu }X_{,\alpha }^\mu X_{,\beta }^\nu g^{\alpha \beta }\sqrt{g}$$
which fulfils this property. For the difference $`LL_0`$ we obtain
$$\frac{\delta (LL_0)}{\delta X^\mu }0$$
Thus, the remaining part is not only covariant in the weak, formal sense enforced by our decision to handle the preferred coordinates as fields. It does not depend on the preferred coordinates $`X^\mu `$. But this is the original “strong” covariance, the classical requirement for the Lagrangian of general relativity. Thus, we can identify the difference $`LL_0`$ with the most general classical Lagrangian of general relativity. In the preferred coordinates we obtain
$$L=(8\pi G)^1\gamma _{\mu \nu }g^{\mu \nu }\sqrt{g}+L_{GR}(g_{\mu \nu })+L_{matter}(g_{\mu \nu },\phi ^m).$$
Note that this Lagrangian fulfils the Einstein equivalence principle in its full beauty. That means, we have derived this principle starting with few general assumptions about a medium in a classical Newtonian framework.
This derivation of exact relativistic symmetry in the context of a classical condensed matter theory is the main result of this paper. To improve our understanding, let’s consider how this has happened, and what has been really used to derive the EEP. The derivation is extremely simple, but given in an unusual formalism.
But how relativistic symmetry appears may be explained without reference to this formalism. There are three principles involved: first, the inherent symmetry of the Lagrange formalism – the equations should be self-adjoint, or “action equals reaction”. Next, there is the relation between preferred coordinates and conservation laws in the Lagrange formalism well-known from Noether’s theorem. And, last not least, we have the independence of the conservation laws from the material properties of the medium enforced by our choice of variables. As a consequence of these principles, the Euler-Lagrange equations for the material properties do not depend on the preferred coordinates. But this is already the EEP:
$$\frac{\delta }{\delta X^\mu }\frac{\delta S}{\delta \varphi ^m}=\frac{\delta }{\delta \varphi ^m}\frac{\delta S}{\delta X^\mu }=\frac{\delta }{\delta \varphi ^m}\text{[cons. laws]}=0$$
The following heuristic, informal picture may be useful for the understanding: The medium is universal. All usual matter fields (gauge fields, fermions) are material properties of this medium, something like defect densities in a crystal. In this picture human beings consist of crystal defects and interact only with other crystal defects. It seems quite obvious that such beings have only restricted observational possibilities. This restriction of observational possibilities leads to relativistic symmetry – we cannot distinguish by observation states which are really different. Thus, there is nothing strange in the appearance of relativistic symmetry in usual condensed matter. It is a natural consequence of the special nature of matter fields in this theory – they are all material properties of a single universal medium.
### 2.1 Equations and Energy-Momentum Tensor
After the derivation of the theory, the “covariant formalism” has done its job, and we can return to a form more appropriate for the comparison with other theories of gravity. We obtain the following equations:
$$G_\nu ^\mu =8\pi G(T_m)_\nu ^\mu +(\mathrm{\Lambda }+\gamma _{\kappa \lambda }g^{\kappa \lambda })\delta _\nu ^\mu 2g^{\mu \kappa }\gamma _{\kappa \nu }.$$
The harmonic condition
$$_\mu (g^{\mu \kappa }\sqrt{g})=0$$
is a consequence of these equations and one form of energy-momentum conservation in the theory. Remarkably, there is also another form – the basic equation may be simply considered as a decomposition of the full energy-momentum tensor $`g^{\mu \kappa }\sqrt{g}`$ into a part which depends on matter fields and a part which depends on the gravitational field:
$$(T_g)_\nu ^\mu =(8\pi G)^1\left(\delta _\nu ^\mu (\mathrm{\Lambda }+\gamma _{\kappa \lambda }g^{\kappa \lambda })G_\nu ^\mu \right)\sqrt{g}$$
Thus, instead of no local conservation law in GR we obtain even two equivalent forms of local conservation laws. The first is equivalent to classical conservation laws from condensed matter theory and, in our variables, to the harmonic condition. The other is equivalent to the conservation law in Noether’s theorem, and splits into a part which depends on the “matter fields” $`\phi ^m`$ and a purely gravitational part.
## 3 Predictions
Using small enough values $`\mathrm{\Xi },\mathrm{{\rm Y}}0`$ leads to the classical Einstein equations. Therefore it is not problematic to fit observation. It is much more problematic to find a way to distinguish our theory from GR by observation.
### 3.1 A dark matter candidate
Let’s consider the influence of the new terms on the expansion of the universe. In our theory a homogeneous universe should be flat. Solutions with non-zero curvature may be solutions of our theory too, but they cannot be homogeneous. The the usual ansatz $`ds^2=d\tau ^2a^2(\tau )(dx^2+dy^2+dz^2)`$ gives
$`3(\dot{a}/a)^2`$ $`=`$ $`\mathrm{{\rm Y}}/a^6+3\mathrm{\Xi }/a^2+\mathrm{\Lambda }+\epsilon `$
$`2(\ddot{a}/a)+(\dot{a}/a)^2`$ $`=`$ $`+\mathrm{{\rm Y}}/a^6+\mathrm{\Xi }/a^2+\mathrm{\Lambda }p`$
We see that $`\mathrm{\Xi }`$ influences the expansion of the universe similar to homogeneous (hot) dark matter with $`p=\frac{1}{3}\epsilon `$.
### 3.2 Big bounce instead of big bang singularity
$`\mathrm{{\rm Y}}`$ becomes important only in the very early universe. But for $`\mathrm{{\rm Y}}>0`$, we obtain a qualitatively different picture. We obtain a lower bound $`a_0`$ for $`a(\tau )`$ defined by
$$\mathrm{{\rm Y}}/a_0^6=3\mathrm{\Xi }/a_0^2+\mathrm{\Lambda }+\epsilon $$
The solution becomes symmetrical in time, with a big crash followed by a big bang. For example, if $`\epsilon =\mathrm{\Xi }=0,\mathrm{{\rm Y}}>0,\mathrm{\Lambda }>0`$ we have the solution
$$a(\tau )=a_0\mathrm{cosh}^{1/3}(\sqrt{3\mathrm{\Lambda }}\tau )$$
In time-symmetrical solutions of this type the horizon is, if not infinite, at least big enough to solve the cosmological horizon problem (cf. ) without inflation.
### 3.3 Frozen stars instead of black holes
The choice $`\mathrm{{\rm Y}}>0`$ influences also another physically interesting solution – the gravitational collapse. There are stable “frozen star” solutions with radius slightly greater than their Schwarzschild radius. The collapse does not lead to horizon formation, but to a bounce from the Schwarzschild radius. Let’s consider an example. The general stable spherically symmetric harmonic metric depends on one step of freedom m(r) and has the form
$$ds^2=(1\frac{m}{r}\frac{m}{r})(\frac{rm}{r+m}dt^2\frac{r+m}{rm}dr^2)(r+m)^2d\mathrm{\Omega }^2$$
Let’s consider the ansatz $`m(r)=(1\mathrm{\Delta })r`$. We obtain
$`ds^2`$ $`=`$ $`\mathrm{\Delta }^2dt^2(2\mathrm{\Delta })^2(dr^2+r^2d\mathrm{\Omega }^2)`$
$`0`$ $`=`$ $`\mathrm{{\rm Y}}\mathrm{\Delta }^2+3\mathrm{\Xi }(2\mathrm{\Delta })^2+\mathrm{\Lambda }+\epsilon `$
$`0`$ $`=`$ $`+\mathrm{{\rm Y}}\mathrm{\Delta }^2+\mathrm{\Xi }(2\mathrm{\Delta })^2+\mathrm{\Lambda }p`$
Now, for very small $`\mathrm{\Delta }`$ even a very small $`\mathrm{{\rm Y}}`$ becomes important, and we obtain a non-trivial stable solution for $`p=\epsilon =\mathrm{{\rm Y}}g^{00}`$. Thus, the surface remains visible, with time dilation $`\sqrt{\epsilon /\mathrm{{\rm Y}}}M^1`$.
## 4 Relativity Principle, Realism and Bohmian mechanics
As we have shown, we have relativistic symmetry for all observable effects (relativity principle for observables). On the other hand, reality itself does not have this relativistic symmetry (no relativity principle for reality).
But this distinction is meaningful only if we have a realistic theory. Here we define realism in the following classical way: Assume we have an experiment described by observables $`X`$ with the observable probability distribution $`\rho _X(X,x)dX`$, which depends on a set of control parameters $`x`$ (the decisions of experimenters). In a realistic theory this is described by some reality $`\lambda \mathrm{\Lambda }`$ with observer-independent probability distribution $`\rho _\lambda (\lambda )d\lambda `$ which explains the observations X with a function $`X(x,\lambda )`$ so that for a test function f we have:
$$f(X)\rho _X(X,x)𝑑X=f(X(x,\lambda ))\rho (\lambda )𝑑\lambda $$
This definition is appropriate to define causality in a natural way: The decision of the experimenter x has a causal influence on X if this function $`X(x,\lambda )`$ depends on x.
These definitions are sufficient to proof Bell’s inequality following . Therefore, the violation of Bell’s inequality defines a contradiction between realism, causality and the relativity principle for reality. If we want to preserve the full relativity principle, we have to reject realism or causality, and the usual decision is the rejection of realism. But in our theory the relativity principle is automatically restricted to observable effects. Therefore, no contradiction appears. We cannot prove Bell’s inequality for space-like separated events and, therefore, have no conflict with Aspect’s experiment. This solution of the puzzle has been preferred by Bell : “the cheapest resolution is something like going back to relativity as it was before Einstein, when people like Lorentz and Poincare thought that there was an aether — a preferred frame of reference — but that our measuring instruments were distorted by motion in such a way that we could not detect motion through the aether. Now, in that way you can imagine that there is a preferred frame of reference, and in this preferred frame of reference things go faster than light.”
In this context, Bohmian mechanics (BM) is very important. BM is a realistic, even deterministic, theory which makes in a “quantum equilibrium” the same predictions as quantum theory. In the relativistic context it requires a preferred frame. This is usually considered as a decisive argument against BM. But in the context of our theory this argument no longer holds: our theory shows the way how to extend BM into the domain of relativistic gravity.
Let’s add a few additional arguments in favour of realism: In our definition, realism is not related with a particular spacetime theory, it does not even depend on the notion of spacetime. This makes realism more fundamental than spacetime theory, and, moreover, more fundamental than a particular spacetime theory like relativity. In the case of conflict, the natural decision is to reject the less fundamental theory. In our case, it is the relativity principle for reality, and not realism. Moreover, there is no independent evidence against realism – this is proven by the existence of BM for all quantum effects. Instead, the problems related with quantum theory, especially the problem of time , may be interpreted as independent evidence against relativity. And, last not least, we loose essentially nothing if we reduce the relativity principle to observable effects. If we reject realism, the relativity principle essentially reduces to observable effects too.
As we see, the relation between our theory, realism, and BM is mutual support. On one hand, compatibility with realism and BM are strong arguments in favour of our theory. On the other hand, our theory weakens the most serious arguments against realism and BM – incompatibility with relativistic principles.
## 5 Quantization
Most workers would agree that “at the root of most of the conceptual problems of quantum gravity” is the idea that “a theory of quantum gravity must have something to say about the quantum nature of space and time” . These problems obviously disappear for a theory of gravity with fixed Newtonian background. Especially this holds for the “problem of time”: “… in quantum gravity, one response to the problem of time is to ‘blame’ it on general relativity’s allowing arbitrary foliations of spacetime; and then to postulate a preferred frame of spacetime with respect to which quantum theory should be written.” . In this way, in our theory the problem of time simply disappears. Now, “most general relativists feel this response is too radical to countenance: they regard foliation-independence as an undeniable insight of relativity.” . That means, this rejection is based on metaphysical preference for the GR spacetime concept only.
It should be noted that together with black holes another quantization problem disappears – the information loss problem. The frozen stars are stable and do not evaporate. Most problems related with energy and momentum conservation disappear too – the Hamiltonian is no longer a constraint.
To solve the ultraviolet problems, we have to make an additional assumption, but a very natural one for a condensed matter theory: an “atomic hypothesis”. After this, the theory is in ideal agreement with “the present educated view on the standard model, and of general relativity, … that these are leading terms in effective field theories” – an idea introduced by Sakharov . An interpretation of $`\rho `$ as the number of “atoms” per volume leads to an interesting prediction for the cutoff:
$$\rho (x)V_{cutoff}=1.$$
It is non-covariant. For a homogeneous “expanding” universe, it seems to expand together with the universe. Thus, the cutoff differs in a principal way from the usual expectation that the cutoff is the Planck length $`a_P10^{33}cm`$ (cf. , ).
## 6 Comparison with other theories of gravity
Because of the simplicity of the additional terms it is no wonder that they have been already considered. Two other theories have a similar Lagrangian for appropriate signs of the cosmological constants: the “relativistic theory of gravity” proposed by Logunov et al. and classical GR with some additional scalar “dark matter” fields. Nonetheless, equations are not all. There are other physical important things which makes the theories different as physical theories, like global restrictions, boundary conditions, causality restrictions, quantization concepts which are closely related with the underlying “metaphysical” assumptions.
### 6.1 Comparison with RTG
The “relativistic theory of gravity” (RTG) proposed by Logunov et al. has Minkowski background metric $`\gamma _{\mu \nu }`$. The Lagrangian of RTG is
$$L=L_{Rosen}+L_{matter}(g_{\mu \nu },\psi ^m)m_g^2(\frac{1}{2}\gamma _{\mu \nu }g^{\mu \nu }\sqrt{g}\sqrt{g}\sqrt{\gamma })$$
which de facto coincides with our theory for $`\mathrm{\Lambda }=m_g^2<0`$, $`\mathrm{\Xi }=\gamma ^{11}m_g^2>0`$, $`\mathrm{{\rm Y}}=\gamma ^{00}m_g^2>0`$.
The metaphysical concept of RTG is completely different. It is a special-relativistic theory, therefore incompatible with classical realism and Bohmian mechanics because of the violation of Bell’s inequality. Another difference is the causality condition: In RTG, only solutions where the light cone of $`g_{ij}`$ is inside the light cone of $`\gamma _{ij}`$ are allowed. A comparable but weaker condition exists in our theory too: $`T(x)`$ should be a time-like function, or, $`\rho (X,T)>0`$.
The metaphysical differences become physical if we consider quantization. Indeed, RTG suggests quantization following standard QFT schemes. Instead, our theory suggests to quantize an atomic model. These ways are conceptually incompatible. Indeed, the prediction for the cutoff length $`l_{cutoff}`$, which is based on the interpretation of $`g^{00}\sqrt{g}`$ as the atomic density, is not Lorentz-covariant and therefore incompatible with RTG.
### 6.2 Comparison with GR plus scalar fields
In the formalism where the preferred coordinates are handled as scalar fields, the Lagrangian of our theory looks equivalent to GR with some dark matter – four scalar fields $`X^\mu `$. Similar “clock fields” in GR have been considered by Kuchar . Usual energy conditions require $`\mathrm{\Xi }>0,\mathrm{{\rm Y}}<0`$. To obtain the most interesting effects (no black holes, no big bang singularity) we have to choose $`\mathrm{{\rm Y}}>0`$, which violates the usual GR energy conditions.
Nonetheless, even if the Lagrangian seems to be the same, the theories are completely different as physical theories. In a typical solution of GR with scalar fields the fields $`X^\mu (x)`$ cannot be used as global coordinates. Especially this holds for all solutions with non-trivial topology. Even if they may be used as global coordinates, the field $`T(x)`$ may be not time-like. All these solutions are forbidden in our theory. <sup>2</sup><sup>2</sup>2That means, their observation falsifies our theory, but does not falsify GR with scalar fields. According to Popper’s criterion of empirical content this means higher empirical content for our theory.
But the remaining solutions – that means, solutions of our theory which may be interpreted as solutions of GR with scalar fields too – are very unnatural from point of this theory. They have very strange boundary values – the fields $`X^\mu (x)`$ are unbounded. Thus, if we consider boundary conditions of type $`|X^\mu (x)|<C`$ as part of GR with scalar fields, then we have no common solutions for above theories.
As we see, to handle preferred coordinates like scalar fields is justified only in a very restricted domain. It does not make our condensed matter theory on a preferred background equivalent to a general-relativistic theory, even if the Lagrangian has the same form. Preferred coordinates and scalar fields remain to be very different things.
Of course, above theories differ also in their metaphysical principles and their quantization concepts. Especially it is incompatible with classical realism and Bohmian mechanics. Similarly, we have a completely different quantization approach, the cutoff length $`l_{cutoff}`$ is not Lorentz-covariant and therefore incompatible with GR metaphysics. In our theory the $`X^\mu `$ are not fields, but fixed background coordinates and therefore should not be quantized, while the “fields” $`X^\mu (x)`$ in GR with scalar fields should be quantized.
## 7 Conclusions
We have started with postulates for a medium in a classical Newtonian world: classical conservation laws and their connection with the Lagrange formalism. We have obtained a viable theory of gravity which, in a certain limit $`\mathrm{\Xi },\mathrm{{\rm Y}}0`$, leads to the classical Einstein equations. We have derived the Einstein equivalence principle from these first principles.
The resulting theory is compatible with classical realism and Bohmian mechanics. This is not only an argument in favour of this theory, but removes serious arguments against realism and Bohmian mechanics, making them viable in the domain of relativistic gravity.
The theory has a lot of other interesting advantages in comparison with general relativity: well-defined local energy and momentum densities, a classical Hamilton formalism, no black hole and big bang singularities, no cosmological horizon problem, a natural dark matter candidate, no problem of time in quantum gravity, no information loss problem, no problems with non-trivial topologies, a natural “atomic ether” quantization concept compatible with modern effective field theory.
Are there serious disadvantages? The relativity principle is restricted to observable effects, but the rejection of realism related with relativistic quantum theory has a similar effect in relativity. SF authors probably don’t like that non-trivial topologies and causal loops are forbidden. Nonetheless, it is an advantage according to Popper’s criterion of empirical content. Is our theory less beautiful than GR? That’s, of course, a matter of taste. But many beautiful aspects of GR appear in our theory too, and some very beautiful concepts often used but of no fundamental importance in GR (ADM decomposition, harmonic gauge) play a fundamental role in our theory. The situation with conservation laws is certainly more beautiful in our theory.
But, even if you nonetheless decide to prefer relativistic theory – the mere existence of a theory which, based on first principles, predicts a variant of the Einstein equations is an interesting fact. As a consequence, the strong empirical evidence in favour of the Einstein equations them-self (Solar system observations, binary pulsars) are no longer support for general-relativistic spacetime concepts. The choice between relativistic spacetime and classical ether should be justified in a different way. |
warning/0001/hep-th0001046.html | ar5iv | text | # The Rest-Frame Darwin Potential from the Lienard-Wiechert Solution in the Radiation Gauge
## I <br>Introduction
Recently the rest-frame Wigner-covariant instant form of dynamics has been developed in Ref. for isolated systems in Minkowski spacetime $`M^4`$ starting from the case of $`N`$ scalar charged particles plus the electromagnetic field. The charges of the particles are described by bilinears in Grassmann variables following the scheme (called pseudo-classical mechanics) which uses a semiclassical approximation to quantum operators with a finite discrete spectrum like spin, that otherwise would have no strict classical limit; Grassmann variables give fermionic oscillators after quantization. The extension of this scheme to the electric charge is based on the experimental fact that all measurable charges are multiples of $`\pm e`$, the electron and positron charges. Therefore, even if it is not clear whether the electric charge has to be considered as a quantum operator in the standard sense (except in the case of the existence of magnetic charges; in this case there is the Dirac quantization rule for the product of the electric and magnetic charges), one can consider it as a two-level system \[which becomes a six-level system ($`\pm e`$, $`\pm \frac{1}{3}e`$, $`\pm \frac{2}{3}e`$) at the quark-lepton level\] described by an operator with quantum $`e`$ instead of $`\mathrm{}`$. Then one can define a semiclassical approximation with Grassmann variables like in the case of spin. As shown in Ref., this semiclassical approximation automatically implies the regularization of the Coulomb self-energies (the $`ij`$ rule). Therefore, this semiclassical approximation may be considered as an alternative to the extended electron models, which were introduced for regularization aims.
The idea leading to the rest-frame instant form is to consider an arbitrary 3+1 splitting of Minkowski spacetime by means of a foliation with spacelike hypersurfaces $`\mathrm{\Sigma }(\tau )`$ diffeomorphic to $`R^3`$. The parameter $`\tau `$ labelling the leaves is used as a Lorentz scalar mathematical time parameter. For each $`\tau `$ the leaf $`\mathrm{\Sigma }(\tau )`$ is defined through the embedding $`R^3\mathrm{\Sigma }(\tau )M^4`$, $`(\tau ,\stackrel{}{\sigma })z^\mu (\tau ,\stackrel{}{\sigma })`$, where $`\stackrel{}{\sigma }`$ are curvilinear coordinates on $`R^3`$. Then one considers the Lagrangian describing the coupling of the given isolated system to an external gravitational field and replaces the 4-metric with the induced metric on $`\mathrm{\Sigma }(\tau )`$, which is a functional of $`z^\mu (\tau ,\stackrel{}{\sigma })`$. In this way one gets the Lagrangian for the description of the isolated system on arbitrary spacelike hypersurfaces (i.e. in arbitrary accelerated reference frames in Minkowski spacetime) with the embedding functions $`z^\mu (\tau ,\stackrel{}{\sigma })`$ as extra configuration variables describing the hypersurface. However, there are four first class constraints at each point implying the independence of the description from the chosen 3+1 splitting. Thus the $`z^\mu (\tau ,\stackrel{}{\sigma })`$’s are gauge variables. Therefore, one can restrict the description of the isolated system to spacelike hyperplanes $`\mathrm{\Sigma }_H(\tau )`$, $`z^\mu (\tau ,\stackrel{}{\sigma })=x_s^\mu (\tau )+b_{\stackrel{ˇ}{r}}^\mu (\tau )\sigma ^{\stackrel{ˇ}{r}}`$ (inertial reference frames in Minkowski spacetime; $`x_s^\mu (\tau )`$ is an arbitrary origin).
Then, if one selects all the configurations of the isolated system with total timelike 4-momentum (they are dense in the space of all configurations), one finds that each timelike configuration identifies a privileged family of hyperplanes: those orthogonal to its total 4-momentum \[Wigner hyperplanes $`\mathrm{\Sigma }_W(\tau )`$\]. At this stage one has obtained the analogue of the nonrelativistic center-of-mass separation and the definition of a new instant form of dynamics , the rest-frame one . There is a decoupled point $`\stackrel{~}{x}_s^\mu (\tau )`$ on each Wigner hyperplane describing the “external” center of mass of the isolated system (thus serving as a decoupled point particle clock) with conjugate momentum $`p_s^\mu .`$ $`\stackrel{~}{x}_s^\mu (\tau )`$ is a canonical variable, but it is not covariant like the Newton-Wigner position operator: it has only covariance under the little group of timelike Poincaré orbits.
After the restriction to the Wigner hyperplane only four first class constraints are left:
i) three of them say that the total 3-momentum of the isolated system vanishes (rest-frame condition): the natural gauge fixing for these constraints is the requirement that the “internal” 3-center of mass of the isolated system inside the Wigner hyperplanes coincides with the origin $`x_s^\mu (\tau )`$ of the coordinates $`\stackrel{}{\sigma }`$ in it (see Refs. , where the group-theoretical results of Ref. are used). In this way only “internal” relative variables describe the isolated system: they are either Lorentz scalars or Wigner spin 1 3-vectors.
ii) the fourth one identifies the invariant mass of the isolated system as the Hamiltonian of the evolution in $`\tau `$ when, with a gauge fixing, $`\tau `$ is made to coincide with the Lorentz scalar rest-frame time $`T_s=p_s\stackrel{~}{x}_s/\sqrt{p_s^2}=p_sx_s/\sqrt{p_s^2}`$ of the decoupled external center of mass.
In this description the standard manifestly covariant fields like the Klein-Gordon field $`\stackrel{~}{\varphi }(z^\mu )`$ are replaced by the new fields $`\varphi (\tau ,\stackrel{}{\sigma })=\stackrel{~}{\varphi }(z^\mu (\tau ,\stackrel{}{\sigma }))`$, which know the embedding and have the non-local information about the equal time hypersurfaces $`\mathrm{\Sigma }(\tau )`$ built-in. In the case of gauge theories, one can make a canonical reduction to a canonical basis of Dirac’s observables in the radiation gauge (or Coulomb or generalized Coulomb; the literature is ambiguous about the terminology to be used), with the only universal breaking of manifest covariance connected with the external center of mass, since the relative motions are Wigner covariant.
See Ref. for a complete review of the research program aiming to give a unified description of the four interactions in terms of Dirac-Bergmann’s observables in the framework of the rest-frame instant form of dynamics, which is the classical background of the Tomonaga-Schwinger formulation of quantum field theory.
The description of scalar (or spinning ) particles on arbitrary spacelike hypersurfaces requires the choice of the sign of the energy of the particle. This happens because the position of the particle on $`\mathrm{\Sigma }(\tau )`$ is identified by 3 numbers $`\stackrel{}{\sigma }=\stackrel{}{\eta }(\tau )`$ and not by 4: this implies that the mass-shell first class constraints of the standard manifestly covariant approach have been solved and that one of the two disjoint branches of the mass spectrum has been chosen. In this way, one gets a different Lagrangian for each branch of the mass spectrum of an isolated system of particles (in the standard manifestly covariant description the branches are topologically disjoint for free particles): there is no possibility of crossing of the branches (the classical background of pair production) when interactions are present inside the isolated system (for instance charged particles plus the electromagnetic field) as happens in the manifestly covariant approach. While there is no problem in the coupling to magnetic fields of these particles with a definite sign of the energy, the minimal coupling in the Lagrangian will miss those couplings to the electric fields which are at the basis of the non-diagonalizability of both the Feshbach-Villars description of the Klein-Gordon field and of the Dirac equation through the Foldy-Wouthuysen transformation (in the case of spinning particles these couplings will have to be extracted from the iterative diagonalization of these theories and added non-minimally). However, this description of particles with, say, positive energy (whose quantization requires pseudodifferential operators ) seems suited for the description of the asymptotic Tomonaga-Schwinger states and will be used to introduce a notion of particle in a future quantization of classical fields on Wigner hyperplanes in the rest-frame instant form. These asymptotic states will replace the Fock ones of the standard manifestly covariant theory, which are the main source of problems in the theory of relativistic quantum bound states (the spurious solutions of the Bethe-Salpeter equation; see Ref.). This framework should allow the introduction of bound states among the asymptotic states.
Coming back to the isolated system formed by $`N`$ scalar charged particles of positive energy plus the electromagnetic field, we recall that in Ref., after canonical reduction to the radiation gauge, the final invariant mass of the reduced system is a function only of Dirac’s observables (gauge invariant particles dressed with a Coulomb cloud and transverse radiation field) and contains:
i) the kinetic energy of the radiation field;
ii) the kinetic energy for the particles with minimal coupling to the radiation field;
iii) the instantaneous action-at-a-distance Coulomb potential among the charges with the Coulomb self-energies regularized (the $`ij`$ rule, $`Q_iQ_j0`$) due to the Grassmann character of the electric charges $`Q_i`$ , i.e. $`Q_i^2=0`$ \[at this semiclassical level we have $`Q_i=e\theta _i^{}\theta _i`$; this does not imply the vanishing of the fine structure constant $`\alpha =e^2/4\pi 1/137`$, since it gets contributions from $`Q_iQ_j`$; therefore, we are retaining effects of order $`\alpha `$ but not of higher order, because, as it will be shown, we do not have many-body forces\].
Then, in Ref. there was the evaluation of the retarded Lienard-Wiechert potentials of the charged particles in the radiation gauge in the rest-frame instant form. Since the Lienard-Wiechert potentials and fields are linearly dependent on the charges $`Q_i`$, the semiclassical regularization $`Q_i^2=0`$ eliminates the radiation coming from a single particle (the electromagnetic energy-momentum tensor contains only terms in $`Q_iQ_j`$, $`ij`$) and the causal problems of the Abraham-Lorentz-Dirac equation of each charged particle \[since the radiation reaction term has the coefficient $`\tau _o=2Q^2/3mc^2`$ ($`\tau _o`$ is proportional to the time needed for light to travel across a classical electron radius), which vanishes if $`Q^2=0`$, there are neither the Schott term nor the Larmor one but only the Lorentz forces produced by the other particles\]. In Refs one may find an extended discussion about this equation and in Ref. a recent review on its derivation and its causal problems (preacceleration, runaway solutions). See Ref. for modern attempts to extract the subset of causal solutions of this equation with the requirement of selecting only its solutions which admit a smooth limit for $`\tau _o0`$ (the runaway solutions are singular in this limit) and to find an effective second order equation for this subset of solutions.
Even if at the semiclassical level the “single charged particle” has no acausal behaviour, because, notwithstanding it produces a Grassmann-valued vector potential, it does not irradiate, we can recover the asymptotic Larmor formula for a system of charges (considered as external sources of the electromagnetic field) due to the interference radiation from $`Q_iQ_j`$, $`ij`$, terms (this result is in accord with macroscopic experimental facts). See Ref. for the extension to spinning particles.
Being in the radiation gauge, at each $`\tau `$ the retarded Lienard-Wiechert potentials evaluated in Ref. contain also a non-local (in $`\stackrel{}{\sigma }`$)term (coming from the transverse projector; see Eqs.(280), (294) in Section V): since this term involves all the points of $`\mathrm{\Sigma }(\tau )`$, the Lienard-Wiechert potential receives contributions from “all” the retarded times before $`\tau `$, namely from the whole past history of the particles. In the absence of incoming radiation one could put these retarded Lienard-Wiechert potentials inside the Lagrangian given in Eq.(74) of Ref. for the description of the isolated system in the rest-frame radiation gauge on the Wigner hyperplanes: one would get a Fokker-like action in the radiation gauge, replacing the standard Fokker-Tetrode one in the Lorentz gauge of the manifestly covariant description, and would have to face the problem of how to find a Hamiltonian description when there are integro-differential equations of motion with delay. The existing attempts are based on the idea of replacing retarded particle coordinates and velocities with instantaneous coordinates and accelerations of all the orders (see for instance Refs.). In this way one replaces integro-defferential equations of motion with an infinite set of coupled differential ones. Since it is unknown how to formulate the Cauchy problem for these integro-differential equations (see Ref.; exceptions are the 1-dimensional case or the time-asymmetric case ), there have been complicated attempts to find conditions for extracting a set of effective second order differential equations from the infinite set (see for instance Refs.). In Ref. there was an attempt to study the Dirac constraints originating from actions depending on accelerations of all orders, following previous attempt of Kerner of defining a Hamiltonian approach. See also the recent approach of Ref..
Moreover, one should face a problem similar to the one raised in Ref., that only with symmetric Green’s functions like $`\frac{1}{2}(retarded+advanced)`$ can the Fokker-Tetrode action corresponding to the Lorentz gauge give rise to a variational principle whose extremals are equivalent to the subspace of extremals of the original action defined by the symmetric Lienard-Wiechert solutions without incoming radiation, i.e. the adjunct Lienard-Wiechert fields (otherwise there are problems with the boundary terms). This is compatible with the Feynman-Wheeler starting point for their theory of the absorbers (see Ref. for the definition of radiation in this theory). A noncovariant justification of the results of Ref. is given in Ref.; by ignoring the self-interactions and assuming a Lagrangian for two charged particles at equal times in which each particle interacts with the retarded Lienard-Wiechert potential of the other one, one obtains in the equations of motion Lorentz forces which correspond to $`\frac{1}{2}(retarded+advanced)`$ interactions, because in this Lagrangian the transition from $`retarded`$ to $`\frac{1}{2}(retarded+advanced)`$ interactions is a total time derivative. However, self-reaction is ignored in these calculations and it is not clear how to arrive at a covariant formulation of these results. Let us remark that from the point of view of quantum field theory its regularization and renormalization require the use of the complex Feynman Green function (which does not vanish outside the lightcone; the solutions with the retarded Green function cannot be regularized at the distributional level): while its imaginary part is connected to absorption in other channels, its real part is just the $`\frac{1}{2}(retarded+advanced)`$ Green function like in Feynman-Wheeler \[ see Ref. for the extraction of a Fokker-Tetrode action with $`\frac{1}{2}(retarded+advanced)`$ kernel from the particle limit of QED (it does not work in QCD)\].
However, both in the Dirac derivation of the Abraham-Lorentz-Dirac equation through the evaluation of the near zone self-field (with the same results obtainable with balance equations using the far zone fields; see Refs.) and in the Feynman-Wheeler approach with the assumption of complete absorbers the radiation is determined by the radiative Green function $`\frac{1}{2}(retardedadvanced)`$ \[at the conceptual level this introduces the acausal advanced Green function and interpretational problems\]. Indeed, the regularization of the self-energy divergence due to radiation reaction is done by rewriting $`retarded=\frac{1}{2}(retarded+advanced)+\frac{1}{2}(retardedadvanced)`$, by noting that the non radiative Coulomb piece of the fields (which does not influence the motion of the particle only giving a divergent electromagnetic contribution to the mass) is in $`\frac{1}{2}(retarded+advanced)`$ and by discarding this term as a regularization.
However, till now all the calculations have been done in the Lorentz gauge and it is not clear whether the previous statements are “gauge invariant”.
Since we now have the results of Refs. in the radiation gauge and a new type of regularization with the semiclassical approximation, it is interesting to revisit all these problems.
The aim of this paper is to show that, starting from the rest-frame instant form description of $`N`$ charged scalar particles plus the electromagnetic field with Grassmann-valued electric charges $`Q_i`$, the semiclassical regularization $`Q_i^2=0`$ allows one to transform the subspace of Lienard-Wiechert solutions (without or with incoming radiation) into a symplectic submanifold of the space of all solutions. This takes place since all the higher accelerations coming from an equal time development of the delay decouple, being of order $`Q_i^2`$ on the solution of the particle equations of motion. As a consequence, at this semiclassical level the retarded, advanced and symmetric Lienard-Wiechert solutions coincide by using the equations of motion, so that there is only one sector of semiclassical Lienard-Wiechert solutions (modulo the incoming radiation). The semiclassical Lienard-Wiechert potential and fields can be expressed as phase space functions and it is possible to eliminate the electromagnetic degrees of freedom by means of second class constraints added to the reduced phase space of the radiation gauge in the rest frame. Having gone to Dirac brackets with respect to these constraints, we get a reduced phase space containing only particles with mutual instantaneous action-at-a-distance interactions. We can find new canonical variables for the particles corresponding to a Darboux basis for these brackets. We can evaluate the final Hamiltonian, showing that besides the Coulomb potential there are vector potentials under the particle kinetic energy square roots (coming from the minimal coupling to the radiation field) and a scalar potential outside them (coming from the energy of the radiation field): due to $`Q_i^2=0`$ one can extract the vector potentials from under the square roots and write a unique effective scalar potential added to the Coulomb one. This can be done both in the original (no longer) canonical variables and in the final canonical basis. The effective scalar potential is the complete Darwin potential: at the lowest order in $`1/c^2`$ we obtain the known form of the Darwin potential.
We find that in the framework of Maxwell (not Feynman-Wheeler) theory, the semiclassical regularization $`Q_i^2=0,Q_iQ_j0`$ extracts automatically the instantaneous action-at-a-distance potential hidden in the delay, which as mentioned above turns out to be the same in the retarded and advanced solutions, because the difference is in the $`Q_i^2`$ terms which depend on the higher accelerations (in QFT these effects are hidden in the radiative corrections coming out from the regularization of the ultraviolet divergences, and this is possible only with the Feynman Green function).
This means that at this semiclassical level, the elimination of the electromagnetic degrees of freedom produces a system of particles with instantaneous action at a distance given by the Coulomb and Darwin potentials. Since $`\frac{1}{2}(retardedadvanced)=0`$, all the effects now come from the regularized $`retarded=advanced=\frac{1}{2}(retarded+advanced)`$ solution and there is no mass renormalization. Even if the transverse projector implies contributions from the whole past (or future) history of the particles, in the semiclassical approximation only the instantaneous action-at-a-distance effects on $`\mathrm{\Sigma }_W(\tau )`$ survive. Each particle feels only the action of the other $`N1`$ \[thus giving us an effective Abraham-Lorentz-Dirac equation with no self-reaction and with the Lorentz forces of the other particles replaced by action-at-a-distance interactions\]. Like in the Feynman-Wheeler approach , we can now speak of radiation only as the effect of the other $`N1`$ particles on the one chosen as a detector of radiation when it is far away from the other particles: equations of motion), one has accord with the Larmor formula coming from the $`Q_iQ_j`$ interference terms; (the Larmor formula gives zero at the semiclassical level due to the particle equations of motion).
An important and new feature of this formalism is the possibility to find the final canonical variables for the particles after the introduction of the Dirac brackets: their use introduce new higher order contributions to the Darwin potential coming from the kinetic energy square roots. These contributions lead to a substantial cancellation with corresponding terms coming from what began as the electromagnetic energy integral. Our final generalized Darwin interaction naturally divides into two portions. The first portion has the same form either as the original lowest order correction derived originally by Darwin in the retarded case or as the lowest order $`1/c^2`$ well known form of the Darwin potential in the case of symmetric \[$`\frac{1}{2}(retarded+advanced)`$ \] Lienard-Wiechert potentials but with the masses replaced by the kinetic energies, $`m_i\sqrt{m_i^2+\kappa _i^2}`$, and this is a new result (strictly speaking this is a higher order correction, but it shows that we are using the correct relativistic kinematics without $`1/c^2`$ expansions). The second portion, a double infinite series, is, like that of the generalization of the $`1/c^{2\text{ }}`$Darwin potential, is also new. It is of higher order in $`1/c^2`$ than the more familiar first portion. Furthermore, our generalized Darwin interaction for $`N`$ bodies is equal to the pairwise sum of two body pieces
For the restricted case of two bodies considerable simplifications result when one evaluates the series in the center-of-mass rest frame. It then can be written in closed form.
The Darwin potential we obtain can be regarded as the classical analogue of the full effects (complete transverse as well as longitudinal) of the single photon exchange in the Bethe-Salpeter equation since it is the same order in the coupling constants. The effect of the semiclassical regularization $`Q_i^2=0`$ is to truncate out the classical analogue of the numerous higher order ladder and cross ladder diagrams. To the extent that the Darwin potential we obtain has a low order ($`1/c^2`$) portion that agrees with the standard result, it would be expected to contribute correctly to the spectral results in a quantized formalism.
For two particles there is the problem of the comparison of the semiclassical Lienard-Wiechert sector with the 2-particle models defined by two first class constraints with covariant instantaneous action-at-a-distance interactions (phenomenological approximations of the Bethe-Salpeter equation). Several authors beginning with and continued by have developed pairs of commuting generalized mass shell conditions, first class constraints with instantaneous potentials in the center-of-mass sytem, whose quantization gives coupled Klein-Gordon equations for two spin zero particles \[see Refs. for similar equations for Dirac particles deriving from pairs of first class constraints for spinning particles\]. Some of these models were generated as approximations to the Bethe-Salpeter equation, by reducing it in a covariant instantaneous approximations to a 3-dimensional equation (with the elimination of the spurious abnormal sectors of relative energy excitations) of the Lippmann-Schwinger type and then to the equation of the quasipotential approach \[see the bibliography of the quoted references\], which Todorov reformulated as a pair of first class constraints at the classical level. In Ref. it is directly shown how the normal sectors of the Bethe-Salpeter equation are connected with the quantization of pairs of first class constraints with instantaneous (in general nonlocal, but approximable with local) potentials like in Todorov’s examples. Ref. shows how to derive the Todorov potential for the electromagnetic and world scalar case from Tetrode-Fokker-Feynman-Wheeler dynamics with scalar and vector potentials \[this theory is connected with $`\frac{1}{2}(retarded+advanced)`$ solutions with no incoming radiation (adjunct Lienard-Wiechert fields) of Maxwell equations with particle currents in the Lorentz gauge\]; besides the Coulomb potential, at the order $`1/c^2`$ one gets the standard Darwin potential (becoming the Breit one at the quantum level when spin is added), which is known to be phenomenologically correct.
What are the connections between the center-of-mass rest frame form of the two-body interaction Hamiltonian that we develope in this paper to all orders in $`1/c^2`$ and that obtained in the above references? The Darwin potential becomes a common overlap of the two approaches and thus an important testing ground for the approach we develope in this paper. Furthermore, when in future papers pseudoclassical spin is introduced (extending as in to interacting systems of particles and fields) the types of tests we perform in this paper will be relevant (note, however, there are difficulties with other categories of Darwin type of interactions brought on by the introduction of transverse spin-dependent electric field effects which unlike magnetic fields cannot be diagnonalized by a Foldy-Wouthuysen transformation ).
In Section II we give a review of parametrized Minkowski theories on arbitrary spacelike hypersurfaces.
In Section III we apply this formalism to the isolated system of $`N`$ charged scalar particles, with Grassmann-valued electric charges, plus the electromagnetic field and we arrive at its rest-frame instant form on the Wigner hyperplanes. We also make the canonical reduction to the radiation gauge.
In Section IV we study the “internal” Poincaré algebra and the “internal” center of mass on the Wigner hyperplanes and we derive the Hamiltonian and Lagrange equations of motion for fields and particles. Also the energy-momentum tensor is evaluated.
In Section V we evaluate the Lienard-Wiechert potentials, we show that at the semiclassical level they depend only on particle coordinates and velocities (since at this level $`retarded=advanced`$) and we find their phase space expression. Then we eliminate the electromagnetic degrees of freedom by means of second class constraints which force them to coincide with the Lienard-Wiechert solution. We introduce the associated Dirac brackets and we find their canonical Darboux basis.
In Section VI we derive the physical Hamiltonian and the effective Darwin potential to all orders in $`1/c^2`$ in terms of the old (noncanonical) and of the new (canonical) variables. In the two-particle case we get a closed form of this potential using the rest-frame condition. Also the final form of the energy-momentum tensor is given.
In the Conclusions there are some general considerations and hints for future developments.
Appendix A gives an explicit summation for the vector potential presented in Section V from the Lienard-Wiechert series in the rest-frame radiation gauge.
In Appendix B we evaluate the field energy and momentum integrals used in Section VI when the electric and magnetic fields are expressed in terms of the Lienard-Wiechert series for the vector potential.
We derive in Appendix C a general formula for a certain quantity, which is important for obtaining the closed form expression of the Darwin potential of Section VI for $`N`$=2 in the rest frame.
In Appendix D we use a technique similar to that developed by Kerner ( and applied in to obtain the Todorov quasipotential from the Wheeler-Feynman action) in the transformation of the Lagrangian expression for the invariant mass to $`M`$. In this proof, (done to all orders) we must use Dirac brackets since we have used the Lienard-Wiechert constraints as a strong condition on the dynamical variables. This necessitates the explicit expression for the field momentum integrals developed in the earlier appendix from the Lienard-Wiechert solutions.
In Appendix E we obtain a special solution of Hamilton’s equations for the two-body problem that is analogous to Schild’s solution for circular orbits .
## II Parametrized Minkowski Theory for N Free Scalar Particles.
In this Section we shall review the description of $`N`$ scalar free particles on arbitrary spacelike hypersurfaces , leaves of the foliation of Minkowski spacetime associated with one of its 3+1 splittings following the suggestions of Refs.. The scalar parameter $`\tau `$ labelling the hypersurfaces $`\mathrm{\Sigma }(\tau )`$ allows one to introduce a covariant concept of “equal time”, which will be useful in the description of an isolated system of interacting particles and fields. It will be shown that in these parametrized Minkowski theories there are first class constraints implying the independence of the description of the isolated system from the chosen 3+1 splitting.
As said in the Introduction this requires the addition to the theory of an infinite number of new configuration variables $`z^\mu (\tau ,\stackrel{}{\sigma })`$ describing the spacelike hypersurfaces as an embedding of $`R^3`$ into Minkowski spacetime. We use the notation $`\sigma ^{\stackrel{ˇ}{A}}=(\sigma ^\tau =\tau ,\sigma ^{\stackrel{ˇ}{r}})`$, i.e. $`\stackrel{ˇ}{A}=(\tau ,\stackrel{ˇ}{r})`$ \[the notation $`A=(\tau ,r)`$ will be used for the Wigner indices on the Wigner hyperplane, see Section III\].
In the manifestly Lorentz covariant approach the worldlines of scalar particles are described by 4-vector coordinates $`x_i^\mu (\tau _i)`$, where the $`\tau _i`$’s are affine parameters (often they are restricted to be the proper times of the particles). Even if one uses a unique affine parameter $`\tau `$ for all the particles, $`x_i^\mu (\tau )`$, there is the problem that the particles times $`x_i^0(\tau )`$ are gauge variables due to the presence of the first class mass-shell constraints $`p_i^2m_i^20`$. For each free particle the constraint manifold is the union of two disjoint submanifolds $`p^0=\pm \sqrt{m_i^2+\stackrel{}{p}_i^2}`$. The gauge nature of the $`x_i^0`$’s is connected with: i) the arbitrariness in the choice of the center-of-mass time; ii) the arbitrariness in the choice of how to trigger the $`N`$ particles (at equal times or with any conceivable mutual delay; this is the gauge freedom of relative times). Given the foliation of Minkowski spacetime with leaves $`\mathrm{\Sigma }(\tau )`$ we can give a covariant description of the particles at “equal times” (covariant zero relative times condition) by parametrizing the worldlines as $`x_i^\mu (\tau )=z^\mu (\eta _i^A(\tau ))=z^\mu (\tau ,\stackrel{}{\eta }_i(\tau ))`$, with with the Lorentz scalar coordinates $`\eta _i^{\stackrel{ˇ}{A}}(\tau )=(\tau ,\eta _i^{\stackrel{ˇ}{r}}(\tau ))`$. Only 3 Lorentz scalar coordinates $`\stackrel{}{\eta }_i(\tau )`$ identify the intersection of the particle worldline with $`\mathrm{\Sigma }(\tau )`$. This implies that in this description there are no mass-shell constraints, namely that we are describing particles with a well defined sign of the energy: $`\eta _i=signp_i^0=\pm 1`$. In this paper we shall consider only positive energy particles, so that $`\eta _i=1`$ for every $`i`$. There will be a conjugate momentum $`\stackrel{}{\kappa }_i(\tau )`$ for each particle and the standard momentum $`p_i^\mu (\tau )`$ will be a derived quantity which satisfies $`p_i^2=m_i^2`$.
The metric induced on $`\mathrm{\Sigma }(\tau )`$ from the Minkowski metric $`\eta ^{\mu \nu }=(+)`$ is
$$g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(z(\tau ,\stackrel{}{\sigma })):=\eta _{\mu \nu }z_{\stackrel{ˇ}{A}}^\mu z_{\stackrel{ˇ}{B}}^\nu ,$$
(1)
where we have used the notation
$$z_{\stackrel{ˇ}{A}}^\mu (\tau ,\stackrel{}{\sigma })=\frac{z^\mu }{\sigma ^{\stackrel{ˇ}{A}}}:=_{\stackrel{ˇ}{A}}z^\mu .$$
(2)
If we define the quantity $`z_\mu ^{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })`$ by means of
$$z_\mu ^{\stackrel{ˇ}{A}}z_{\stackrel{ˇ}{B}}^\mu =\delta _{\stackrel{ˇ}{B}}^{\stackrel{ˇ}{A}},$$
(3)
they satisfy
$$g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}z_\mu ^{\stackrel{ˇ}{A}}z_\nu ^{\stackrel{ˇ}{B}}=\eta _{\lambda \kappa }z_{\stackrel{ˇ}{A}}^\lambda z_{\stackrel{ˇ}{B}}^\kappa z_\mu ^{\stackrel{ˇ}{A}}z_\nu ^{\stackrel{ˇ}{B}}=\eta _{\mu \nu }.$$
(4)
Therefore, the $`z_\mu ^{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })`$ are a set of vierbeins, with the $`z_{\stackrel{ˇ}{A}}^\mu (\tau ,\stackrel{}{\sigma })`$ the inverse vierbeins.
Since we require $`g_{\tau \tau }>0`$ as a condition on the embedding $`z^\mu (\tau ,\stackrel{}{\sigma })`$, $`z_\tau ^\mu `$ is a time-like 4-vector and the $`z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })`$’s are spacelike 4-vectors tangent to $`\mathrm{\Sigma }(\tau )`$.
The determinant of the metric is
$$g=detg_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}=(detz_{\stackrel{ˇ}{A}}^\mu )^2.$$
(5)
The spatial part of the metric has an associated determinant which is defined by
$$\gamma =detg_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}};\mathrm{\Gamma }=\sqrt{\gamma }.$$
(6)
We next define the inverse metric $`g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}`$ by
$$g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}g_{\stackrel{ˇ}{B}\stackrel{ˇ}{C}}=\delta _{\stackrel{ˇ}{C}}^{\stackrel{ˇ}{A}}.$$
(7)
This implies
$$g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}=4=g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}z_{\stackrel{ˇ}{A}}^\mu z_{\stackrel{ˇ}{B}}^\nu \eta _{\mu \nu }=\eta ^{\mu \nu }\eta _{\mu \nu },$$
(8)
so that
$$\eta ^{\mu \nu }=g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}z_{\stackrel{ˇ}{A}}^\mu z_{\stackrel{ˇ}{B}}^\nu =g^{\tau \tau }z_\tau ^\mu z_\tau ^\nu +2g^{\tau \stackrel{ˇ}{r}}z_\tau ^\mu z_{\stackrel{ˇ}{r}}^\nu +g^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}z_{\stackrel{ˇ}{r}}^\mu z_{\stackrel{ˇ}{s}}^\nu .$$
(9)
By definition of the element of an inverse matrix
$$g^{\tau \tau }=\frac{\mathrm{\Gamma }^2}{g}=\frac{\gamma }{g},$$
(10)
while the inverse of the spatial $`g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}`$ is defined as $`\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}`$, that is,
$$\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}g_{\stackrel{ˇ}{s}\stackrel{ˇ}{t}}=\delta _{\stackrel{ˇ}{t}}^{\stackrel{ˇ}{r}}.$$
(11)
To find $`g^{\tau \stackrel{ˇ}{u}}`$ use the fact that
$$g^{\tau \stackrel{ˇ}{r}}g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}+g^{\tau \tau }g_{\tau \stackrel{ˇ}{s}}=\delta _{\stackrel{ˇ}{s}}^\tau =0,$$
(12)
and therefore
$$g^{\tau \stackrel{ˇ}{r}}g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}=g^{\tau \tau }g_{\tau \stackrel{ˇ}{s}}.$$
(13)
Multiplying by $`\gamma ^{\stackrel{ˇ}{s}\stackrel{ˇ}{u}}`$ leaves us with
$$g^{\tau \stackrel{ˇ}{u}}=\frac{\mathrm{\Gamma }^2}{g}g_{\tau \stackrel{ˇ}{s}}\gamma ^{\stackrel{ˇ}{s}\stackrel{ˇ}{u}}.$$
(14)
Now using this consider
$$g^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}g_{\stackrel{ˇ}{t}\stackrel{ˇ}{s}}+g^{\stackrel{ˇ}{r}\tau }g_{\stackrel{ˇ}{t}\tau }=\delta _{\stackrel{ˇ}{t}}^{\stackrel{ˇ}{r}}.$$
(15)
Multiply both sides by $`\gamma ^{\stackrel{ˇ}{t}\stackrel{ˇ}{u}}`$, use the definition of $`\gamma ^{\stackrel{ˇ}{t}\stackrel{ˇ}{u}}`$ and the above expression for $`g^{\stackrel{ˇ}{r}\tau }`$, and we obtain
$$g^{\stackrel{ˇ}{r}\stackrel{ˇ}{u}}=\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{u}}+\frac{\mathrm{\Gamma }^2}{g}g_{\tau \stackrel{ˇ}{s}}g_{\tau \stackrel{ˇ}{v}}\gamma ^{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}\gamma ^{\stackrel{ˇ}{v}\stackrel{ˇ}{u}}.$$
(16)
Thus in summary we have expressed the inverse metric in terms of the metric and the inverse of its spatial parts
$`g^{\tau \tau }`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2}{g}},`$ (17)
$`g^{\tau \stackrel{ˇ}{r}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }^2}{g}}g_{\tau \stackrel{ˇ}{s}}\gamma ^{\stackrel{ˇ}{s}\stackrel{ˇ}{r}},`$ (18)
$`g^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}`$ $`=`$ $`\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}+{\displaystyle \frac{\mathrm{\Gamma }^2}{g}}g_{\tau \stackrel{ˇ}{u}}g_{\tau \stackrel{ˇ}{v}}\gamma ^{ur}\gamma ^{vs}.`$ (19)
Moreover, we have
$$\eta ^{\mu \nu }z_\mu ^{\stackrel{ˇ}{A}}z_\nu ^{\stackrel{ˇ}{B}}=g^{\stackrel{ˇ}{C}\stackrel{ˇ}{D}}z_{\stackrel{ˇ}{C}}^\mu z_{\stackrel{ˇ}{D}}^\nu z_\mu ^{\stackrel{ˇ}{A}}z_\nu ^{\stackrel{ˇ}{B}}=g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}.$$
(20)
The normal to $`\mathrm{\Sigma }(\tau )`$ at the point $`z^\mu (\tau ,\stackrel{}{\sigma })`$ is the Lorentz four vector
$$l^\mu (\tau ,\stackrel{}{\sigma })=\frac{1}{\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })}ϵ^{\mu \alpha \beta \gamma }z_{\stackrel{ˇ}{1}\alpha }(\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{2}\beta }(\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{3}\gamma }(\tau ,\stackrel{}{\sigma }).$$
(21)
with the normalization $`l^2(\tau ,\stackrel{}{\sigma })=1`$. By construction we have $`l_\mu (\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })=0`$.
The evolution 4-vector $`z_\tau ^\mu (\tau ,\stackrel{}{\sigma })`$ can be decomposed on the mutually orthogonal four vectors $`l^\mu `$ and $`z_s^\mu `$:
$$z_\tau ^\mu (\tau ,\stackrel{}{\sigma })=N(\tau ,\stackrel{}{\sigma })l^\mu (\tau ,\stackrel{}{\sigma })+N^{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{s}}^\mu (\tau ,\stackrel{}{\sigma }).$$
(22)
where $`N(\tau ,\stackrel{}{\sigma })`$ is the lapse and $`N^{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })`$ the shift functions in the terminology of ADM general relativity. To determine the vector $`N^{\stackrel{ˇ}{r}}`$ we use the orthogonality of $`z_{\mu \stackrel{ˇ}{r}}`$ and $`l^\mu `$. That is multiplying both sides of the above equation by $`z_{\mu \stackrel{ˇ}{r}}`$ and using the expression for $`g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}`$ and multiplying by $`\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{\stackrel{ˇ}{u}}}`$ we obtain.
$$N^{\stackrel{ˇ}{u}}(\tau ,\stackrel{}{\sigma })=g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{u}}(\tau ,\stackrel{}{\sigma }).$$
(23)
Using this expression for $`N^{\stackrel{ˇ}{r}}`$ we determine the scalar $`N`$ by first multiplying the previous equation by $`z_{\mu \tau }`$. Then we use the definition of $`l^\mu `$ and the determinant, multiply the result by $`g^{\tau \tau }`$. Then use the definition of $`\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}`$ , and we obtain $`N=\mathrm{\Gamma }/\sqrt{\gamma }`$. Hence
$$z_\tau ^\mu (\tau ,\stackrel{}{\sigma })=\left[\frac{\mathrm{\Gamma }}{\sqrt{\gamma }}l^\mu +g_{\tau \stackrel{ˇ}{r}}\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}z_{\stackrel{ˇ}{s}}^\mu \right](\tau ,\stackrel{}{\sigma }).$$
(24)
Substituting this expression for $`z_\tau ^\mu `$ into
$$\eta ^{\mu \nu }=g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}z_{\stackrel{ˇ}{A}}^\mu z_{\stackrel{ˇ}{B}}^\nu =g^{\tau \tau }z_\tau ^\mu z_\tau ^\nu +g^{\tau \stackrel{ˇ}{r}}(z_\tau ^\mu z_{\stackrel{ˇ}{r}}^\nu +z_\tau ^\nu z_{\stackrel{ˇ}{r}}^\mu )+g^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}z_{\stackrel{ˇ}{r}}^\mu z_{\stackrel{ˇ}{s}}^\nu ,$$
(25)
together with those for $`g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}`$, we obtain after some algebra the following decomposition of the Minkowski metric
$$\eta ^{\mu \nu }=l^\mu (\tau ,\stackrel{}{\sigma })l^\nu (\tau ,\stackrel{}{\sigma })+\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{s}}^\nu (\tau ,\stackrel{}{\sigma }).$$
(26)
Coming back to the scalar particles, the relation between their world line velocities in the two descriptions is
$$\dot{x_i^\mu }(\tau )=z_\tau ^\mu (\tau ,\stackrel{}{\eta }_i(\tau ))+z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}(\tau ).$$
(27)
Noting that
$`\dot{x_i}^2`$ $`=`$ $`\dot{x_i^\mu }\dot{x_i^\nu }\eta _{\mu \nu }=z_\tau ^\mu z_\tau ^\nu \eta _{\mu \nu }+2z_\tau ^\mu z_{\stackrel{ˇ}{r}}^\nu \eta _{\mu \nu }\dot{\eta _i^{\stackrel{ˇ}{r}}}+z_{\stackrel{ˇ}{s}}^\mu z_{\stackrel{ˇ}{r}}^\nu \eta _{\mu \nu }\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}}=`$ (28)
$`=`$ $`g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}},`$ (29)
the standard action for $`N`$ free scalar particles becomes
$$S=𝑑\tau \underset{i}{\overset{N}{}}\left[m_i\sqrt{\dot{x_i}^2}\right]=𝑑\tau L(\tau )=𝑑\tau d^3\sigma (\tau ,\stackrel{}{\sigma }),$$
(30)
with the Lagrangian density
$$(\tau ,\stackrel{}{\sigma })=\underset{i=1}{\overset{N}{}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))m_i\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\sigma })+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )},$$
(31)
and the Lagrangian
$$L(\tau )=\underset{i=1}{\overset{N}{}}m_i\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}.$$
(32)
The above action is invariant under separate $`\tau `$ and $`\stackrel{}{\sigma }`$ reparametrization. This leads naturally to constraints.
The canonical momenta are determined from the dependence on the “velocity” $`z_\tau ^\mu (\tau ,\stackrel{}{\sigma })`$ associated with the hypersurface and the particle velocities $`\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau ).`$
$`\rho _\mu (\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{(\tau ,\stackrel{}{\sigma })}{z_\tau ^\mu (\tau ,\stackrel{}{\sigma })}}={\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))`$ (34)
$`m_i{\displaystyle \frac{z_{\tau \mu }(\tau ,\stackrel{}{\sigma })+z_{\stackrel{ˇ}{r}\mu }(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )}{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\sigma })+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}}}=`$
$`=`$ $`[(\rho _\nu l^\nu )l_\mu +(\rho _\nu z_{\stackrel{ˇ}{r}}^\nu )\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}z_{\stackrel{ˇ}{s}\mu }](\tau ,\stackrel{}{\sigma }),`$ (35)
$`\kappa _{i\stackrel{ˇ}{r}}(\tau )`$ $`=`$ $`{\displaystyle \frac{L(\tau )}{\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )}}=`$ (37)
$`=`$ $`m_i{\displaystyle \frac{g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}}},`$ (38)
Using the above we have
$`m_i^2\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}\kappa _{i\stackrel{ˇ}{r}}\kappa _{i\stackrel{ˇ}{s}}=m_i^2[1{\displaystyle \frac{\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(g_{\tau \stackrel{ˇ}{r}}+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{t}}\dot{\eta _i^{\stackrel{ˇ}{t}}})(g_{\tau \stackrel{ˇ}{s}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{u}}\dot{\eta _i^u})}{g_{\tau \tau }+2g_{\tau \stackrel{ˇ}{r}}\dot{\eta _i^{\stackrel{ˇ}{r}}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}}}}]=`$ (39)
$`=({\displaystyle \frac{m_i}{\sqrt{g_{\tau \tau }+2g_{\tau r}\dot{\eta _i^{\stackrel{ˇ}{r}}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}}}}})^2(g_{\tau \tau }\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}g_{\tau \stackrel{ˇ}{r}}g_{\tau \stackrel{ˇ}{s}}).`$ (40)
Use the following two forms
$`\rho _\mu l^\mu `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\delta ^3(\stackrel{}{\sigma }\eta _i(\tau ))m_iz_{\tau \mu }l^\mu }{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}}}}},`$ (41)
$`z_\tau l`$ $`=`$ $`{\displaystyle \frac{1}{\gamma }}ϵ^{\mu \alpha \beta \mathrm{\Gamma }}z_{\tau \mu }z_{1\alpha }z_{2\beta }z_{3\gamma }={\displaystyle \frac{\sqrt{g}}{\mathrm{\Gamma }}},`$ (42)
together with the square root of the above relation and $`g^{\tau \stackrel{ˇ}{s}}=\frac{\gamma }{g}g_{\tau \stackrel{ˇ}{r}}\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}`$and
$$g^{\tau \stackrel{ˇ}{s}}g_{\tau \stackrel{ˇ}{s}}=g^{\tau A}g_{\tau A}g^{\tau \tau }g_{\tau \tau }=1\frac{\gamma }{g}g_{\tau \tau },$$
(43)
to obtain
$$\rho _\mu l^\mu =\underset{i=1}{\overset{N}{}}\delta ^3(\stackrel{}{\sigma }\eta _i(\tau ))\sqrt{m_i^2\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}\kappa _{i\stackrel{ˇ}{r}}\kappa _{i\stackrel{ˇ}{s}}}.$$
(44)
Using finally
$`\rho _\mu z_{\stackrel{ˇ}{r}}^\mu `$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\delta ^3(\stackrel{}{\sigma }\eta _i(\tau ))m_i(z_{\stackrel{ˇ}{r}}^\mu z_{\tau \mu }+z_{\stackrel{ˇ}{r}}^\mu z_{\stackrel{ˇ}{s}\mu }\dot{\eta _i^{\stackrel{ˇ}{s}}})}{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau r}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}+g_{\stackrel{ˇ}{s}\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta _i^{\stackrel{ˇ}{r}}}\dot{\eta _i^{\stackrel{ˇ}{s}}}}}}=`$ (45)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\eta _i(\tau ))\kappa _{ri},`$ (46)
we obtain the form of the four primary first class constraints $`_\mu `$ following from $`\tau `$ and $`\stackrel{}{\sigma }`$ reparametrization invariance:
$`_\mu (\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\rho _\mu (\tau ,\stackrel{}{\sigma })l_\mu (\tau ,\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\sqrt{m_i^2\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\kappa _{i\stackrel{ˇ}{r}}(\tau )\kappa _{i\stackrel{ˇ}{s}}(\tau )}`$ (47)
$``$ $`z_{\stackrel{ˇ}{r}\mu }(\tau ,\stackrel{}{\sigma })\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\kappa _{i\stackrel{ˇ}{s}}0.`$ (48)
Assuming the following Poisson brackets
$`\{z^\mu (\tau ,\stackrel{}{\sigma }),\rho _\nu (\tau ,\stackrel{}{\sigma }^{^{}})\}`$ $`=`$ $`\eta _\nu ^\mu \delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{^{}}),`$ (49)
$`\{\eta _i^{\stackrel{ˇ}{r}}(\tau ),\kappa _{j\stackrel{ˇ}{s}}(\tau )\}`$ $`=`$ $`\delta _{ij}\delta _{\stackrel{ˇ}{s}}^{\stackrel{ˇ}{r}}.`$ (50)
one can show the first class nature of the above constraints . Since these constraints are solved in terms of one of the independent momenta \[$`\rho _\mu (\tau ,\stackrel{}{\sigma })`$\], their Poisson brackets are exactly zero:
$$\{_\mu (\tau ,\stackrel{}{\sigma }),_\nu (\tau ,\stackrel{}{\sigma }^{^{}})\}=0.$$
(51)
These constraints imply that the description of the system is independent from the chosen 3+1 splitting of Minkowski spacetime.
The standard particle momenta $`p_i^\mu `$ are reconstructed as $`p_i^\mu =(\sqrt{m_i^2\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}\kappa _{i\stackrel{ˇ}{r}}\kappa _{i\stackrel{ˇ}{s}}};\stackrel{}{\kappa }_i)`$ and satisfy $`p_i^2=m_i^2`$.
In the next Section we shall add Grassmann-valued electric charges to the scalar particles, we shall find the new action and constraints, eventually arriving at the rest-frame instant form.
## III N Charged Scalar Particles and the Electromagnetic Field.
In this Section we will extend the formalism of the previous Section to the case of an isolated system of $`N`$ charged scalar particles plus the electromagnetic field. By using Lorentz scalar electromagnetic potentials and field strengths, which employ the covariant “equal time” concept associated with the spacelike hypersurfaces $`\mathrm{\Sigma }(\tau )`$, we define the action for the combined field and particle system, that, like the free system, is separately invariant under $`\tau `$ and $`\stackrel{}{\sigma }`$ reparametrizations. As in the free particle case we obtain four primary first class constraints $`^\mu (\tau ,\stackrel{}{\sigma })0`$ at each point $`\stackrel{}{\sigma }`$ on each of the space-like surfaces $`\mathrm{\Sigma }(\tau )`$, implying the independence from the chosen 3+1 splitting. Moreover, there are the two additional first class constraints describing the electromagnetic gauge invariance of the theory. By using a gauge fixing condition which restricts the $`\mathrm{\Sigma }(\tau )`$ ’s to hyperplanes $`\mathrm{\Sigma }_H(\tau )`$ and by using Dirac brackets, the embedding variables $`z^\mu (\tau ,\stackrel{}{\sigma })`$ are reduced to only ten ones. The original constraints $`^\mu (\tau ,\stackrel{}{\sigma })0`$ are reduced to just 10 first class global constraints on each of the hyperplanes. We then specialize to hyperplanes orthogonal to the total timelike four-momentum of the system with 6 new gauge fixings depending on the standard Wigner boost for timelike Poincaré orbits. These hyperplanes, defined by the system configuration, are called the Wigner hyperplanes $`\mathrm{\Sigma }_W(\tau )`$. After having defined the new Dirac brackets, we remain with a decoupled “external” canonical non-covariant center of mass, with Wigner-covariant particle and field degrees of freedom on the Wigner hyperplane and with only four first class contraints. One of these four constraints identifies the invariant mass of the system as the effective Hamiltonian, while the other three define the rest-frame condition $`\stackrel{}{}_p(\tau )0`$ (vanishing of the total 3-momentum inside the Wigner hyperplane) for the combined system of particles and fields. Finally, we make the canonical reduction to eliminate the electromagnetic gauge degrees of freedom, placing the formalism (including the Poincaré generators) in the Wigner covariant rest-frame radiation gauge. Also we give the energy-momentum tensor of the full isolated system.
### A The Action, The Constraints and the Canonical Reduction.
Let us now review the isolated system of $`N`$ charged scalar particles plus the electromagnetic field following Ref.. Just as on the hypersurface $`\mathrm{\Sigma }(\tau )`$ the positive energy particles are described by coordinates $`\stackrel{}{\eta }_i(\tau )`$ such that $`x_i^\mu (\tau )=z^\mu (\tau ,\stackrel{}{\eta }_i(\tau ))`$, so the electric charge of each particle is described in a semiclassical way by means of a pair of complex conjugate Grassmann variables $`\theta _i(\tau ),\theta _i^{}(\tau )`$ satisfying \[$`I_i=I_i^{}=\theta _i^{}\theta _i`$ is the generator of the $`U_{em}(1)`$ group of particle $`i`$\]
$`\theta _i^2=\theta _i^2=0,\theta _i\theta _i^{}+\theta _i^{}\theta _i=0,`$ (52)
$`\theta _i\theta _j=\theta _j\theta _i,\theta _i\theta _j^{}=\theta _j^{}\theta _i,\theta _i^{}\theta _j^{}=\theta _j^{}\theta _i^{},ij.`$ (53)
As said in the Introduction, the formal quantization procedure sends $`\theta _i^{}`$, $`\theta _i`$ into the Clifford algebra describing a two-level Fermi oscillator $`b_i^{}`$, $`b_i`$, and each Grassmann-valued electric charge $`Q_i=e\theta _i^{}\theta _i`$ goes into $`eb_i^{}b_i`$ with eigenvalues $`\pm e`$.
The standard electromagnetic potential $`A_\mu (z)`$ and the field strength $`F_{\mu \nu }(z)`$ are replaced on $`\mathrm{\Sigma }(\tau )`$ by Lorentz-scalar variables $`A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })`$ and $`F_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })`$ respectively, defined by
$`A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })=z_{\stackrel{ˇ}{A}}^\mu (\tau ,\stackrel{}{\sigma })A_\mu (z(\tau ,\stackrel{}{\sigma })),`$ (54)
$`F_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })=_{\stackrel{ˇ}{A}}A_{\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })_{\stackrel{ˇ}{B}}A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })=z_{\stackrel{ˇ}{A}}^\mu (\tau ,\stackrel{}{\sigma })z_{\stackrel{ˇ}{B}}^\nu (\tau ,\stackrel{}{\sigma })F_{\mu \nu }(z(\tau ,\stackrel{}{\sigma })).`$ (55)
The new potentials $`A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })`$ have built-in the covariant concept of “equal time” through their implicit dependence on the embeddings $`z^\mu (\tau ,\stackrel{}{\sigma })`$.
With $`d^3\mathrm{\Sigma }^\mu `$ the surface element of $`\mathrm{\Sigma }(\tau )`$ we have the following volume element of Minkowski space-time
$$d^4z=z_\tau ^\mu d\tau d^3\mathrm{\Sigma }_\mu =d\tau z_\tau ^\mu l_\mu \mathrm{\Gamma }d^3\sigma =\sqrt{g}d\tau d^3\sigma .$$
(56)
The action now depends on the configuration variables $`z^\mu (\tau ,\stackrel{}{\sigma })`$,$`A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{}{\eta }_i(\tau )`$, $`\theta _i(\tau )`$ and $`\theta _i^{}(\tau )`$,$`i=1,..,N`$, and consists of a “kinetic” piece for the complex Grassmann charges $`\frac{i}{2}[\theta _i^{}(\tau )\dot{\theta }_i(\tau )\dot{\theta }_i^{}(\tau )\theta _i(\tau )]𝑑\tau `$, the same particle kinetic piece as in the previous Section, the kinetic term $`d^4z(\frac{1}{4}F^{\mu \nu }F_{\mu \nu })`$ for the electromagnetic field, and the field-particle interaction term $`Q_iA_\mu (x_i(\tau ))\dot{x}_i^\mu \left(\tau \right)𝑑\tau `$:
$`S`$ $`=`$ $`{\displaystyle 𝑑\tau d^3\sigma (\tau ,\stackrel{}{\sigma })}={\displaystyle 𝑑\tau L(\tau )},`$ (57)
$`L(\tau )`$ $`=`$ $`{\displaystyle d^3\sigma (\tau ,\stackrel{}{\sigma })},`$ (58)
$`(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))[\theta _i^{}(\tau )\dot{\theta }_i(\tau )\dot{\theta }_i^{}(\tau )\theta _i(\tau )]`$ (62)
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))[m_i\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\sigma })+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}+`$
$`+Q_i(\tau )(A_\tau (\tau ,\stackrel{}{\sigma })+A_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau ))]`$
$`{\displaystyle \frac{1}{4}}\sqrt{g(\tau ,\stackrel{}{\sigma })}g^{\stackrel{ˇ}{A}\stackrel{ˇ}{C}}(\tau ,\stackrel{}{\sigma })g^{\stackrel{ˇ}{B}\stackrel{ˇ}{D}}(\tau ,\stackrel{}{\sigma })F_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })F_{\stackrel{ˇ}{C}\stackrel{ˇ}{D}}(\tau ,\stackrel{}{\sigma }),`$
$`Q_i(\tau )`$ $`=`$ $`e\theta _i^{}(\tau )\theta _i(\tau ).`$ (64)
Since $`A_\tau (\tau ,\stackrel{}{\sigma })`$ transforms as a $`\tau `$-derivative the action is still invariant under separate $`\tau `$\- and $`\stackrel{}{\sigma }`$-reparametrizations as in the free case. In addition it is invariant under the electromagnetic local gauge transformations and under the odd global phase transformations $`\delta \theta _ii\alpha \theta _i`$ , generated by the $`I_i`$’s. The $`Q_i=eI_i`$ are the constants of motion associated with this last symmetry \[$`\frac{d}{d\tau }Q_i(\tau )\stackrel{}{=}\mathrm{\hspace{0.17em}0}`$, where ’$`\stackrel{}{=}`$’ means evaluated on the solutions of the Euler-Lagrange equations; from now on we shall write $`Q_i`$ instead of $`Q_i(\tau )`$\].
Since the semiclassical approximation $`Q_i^2=0`$ regularizes the Coulomb self-energy, the criticism of Rohrlich to this action principle \[that the minimal coupling term and the electromagnetic field term are ill defined because they diverge on the worldlines of the particles\] does not apply.
The canonical momenta are \[$`E_{\stackrel{ˇ}{r}}=F_{\stackrel{ˇ}{r}\tau }`$ and $`B_{\stackrel{ˇ}{r}}=\frac{1}{2}ϵ_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}\stackrel{ˇ}{t}}F_{\stackrel{ˇ}{s}\stackrel{ˇ}{t}}`$ ($`ϵ_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}\stackrel{ˇ}{t}}=ϵ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}\stackrel{ˇ}{t}}`$) are the “electric” and “magnetic” fields respectively; for $`g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}\eta _{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}`$ one gets $`\pi ^{\stackrel{ˇ}{r}}=E_{\stackrel{ˇ}{r}}=E^{\stackrel{ˇ}{r}}`$\]
$`\rho _\mu (\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{(\tau ,\stackrel{}{\sigma })}{z_\tau ^\mu (\tau ,\stackrel{}{\sigma })}}={\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))m_i`$ (68)
$`{\displaystyle \frac{z_{\tau \mu }(\tau ,\stackrel{}{\sigma })+z_{\stackrel{ˇ}{r}\mu }(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )}{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\sigma })+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}}}+`$
$`+{\displaystyle \frac{\sqrt{g(\tau ,\stackrel{}{\sigma })}}{4}}[z_\mu ^\tau (\tau ,\stackrel{}{\sigma })g^{\stackrel{˘}{A}\stackrel{˘}{C}}(\tau ,\stackrel{}{\sigma })g^{\stackrel{˘}{B}\stackrel{˘}{D}}(\tau ,\stackrel{}{\sigma })2z_\mu ^{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })g^{\stackrel{˘}{B}\stackrel{˘}{D}}(\tau ,\stackrel{}{\sigma })g^{\overline{C}\tau }(\tau ,\stackrel{}{\sigma })`$
$`2z_\mu ^{\stackrel{ˇ}{D}}(\tau ,\stackrel{}{\sigma })g^{B\tau }(\tau ,\stackrel{}{\sigma })g^{\stackrel{˘}{A}\stackrel{˘}{C}}(\tau ,\stackrel{}{\sigma })]F_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })F_{\overline{C}\overline{D}}(\tau ,\stackrel{}{\sigma })`$
$`=`$ $`[(\rho _\nu l^\nu )l_\mu +(\rho _\nu z_{\stackrel{ˇ}{r}}^\nu )\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}z_{\stackrel{ˇ}{s}\mu }](\tau ,\stackrel{}{\sigma }),`$ (69)
$`\pi ^\tau (\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{L}{_\tau A_\tau (\tau ,\stackrel{}{\sigma })}}=0,`$ (71)
$`\pi ^{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{L}{_\tau A_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })}}={\displaystyle \frac{\gamma (\tau ,\stackrel{}{\sigma })}{\sqrt{g(\tau ,\stackrel{}{\sigma })}}}\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })(F_{\tau \stackrel{ˇ}{s}}g_{\tau \stackrel{ˇ}{v}}\gamma ^{\stackrel{ˇ}{v}\stackrel{ˇ}{u}}F_{\stackrel{ˇ}{u}\stackrel{ˇ}{s}})(\tau ,\stackrel{}{\sigma })=`$ (72)
$`=`$ $`{\displaystyle \frac{\gamma (\tau ,\stackrel{}{\sigma })}{\sqrt{g(\tau ,\stackrel{}{\sigma })}}}\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })(E_{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })+g_{\tau \stackrel{ˇ}{v}}(\tau ,\stackrel{}{\sigma })\gamma ^{\stackrel{ˇ}{v}\stackrel{ˇ}{u}}(\tau ,\stackrel{}{\sigma })ϵ_{\stackrel{ˇ}{u}\stackrel{ˇ}{s}\stackrel{ˇ}{t}}B_{\stackrel{ˇ}{t}}(\tau ,\stackrel{}{\sigma })),`$ (73)
$`\kappa _{i\stackrel{ˇ}{r}}(\tau )`$ $`=`$ $`{\displaystyle \frac{L(\tau )}{\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )}}=`$ (75)
$`=`$ $`m_i{\displaystyle \frac{g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}{\sqrt{g_{\tau \tau }(\tau ,\stackrel{}{\eta }_i(\tau ))+2g_{\tau \stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{s}}(\tau )}}}+`$ (76)
$`+`$ $`Q_iA_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau )),`$ (77)
$`\pi _{\theta i}(\tau )`$ $`=`$ $`{\displaystyle \frac{L(\tau )}{\dot{\theta }_i(\tau )}}={\displaystyle \frac{i}{2}}\theta _i^{}(\tau ),`$ (79)
$`\pi _{\theta ^{}i}(\tau )`$ $`=`$ $`{\displaystyle \frac{L(\tau )}{\dot{\theta }_i^{}(\tau )}}={\displaystyle \frac{i}{2}}\theta _i(\tau ).`$ (80)
The following Poisson brackets are assumed
$`\{z^\mu (\tau ,\stackrel{}{\sigma }),\rho _\nu (\tau ,\stackrel{}{\sigma }^{^{}}\}=\eta _\nu ^\mu \delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{^{}}),`$ (81)
$`\{A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma }),\pi ^{\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma }^{^{}})\}=\eta _{\stackrel{ˇ}{A}}^{\stackrel{ˇ}{B}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{^{}}),`$ (82)
$`\{\eta _i^{\stackrel{ˇ}{r}}(\tau ),\kappa _{j\stackrel{ˇ}{s}}(\tau )\}=\delta _{ij}\delta _{\stackrel{ˇ}{s}}^{\stackrel{ˇ}{r}},`$ (83)
$`\{\theta _i(\tau ),\pi _{\theta j}(\tau )\}=\delta _{ij},`$ (84)
$`\{\theta _i^{}(\tau ),\pi _{\theta ^{}j}(\tau )\}=\delta _{ij}.`$ (85)
The Grassmann momenta give rise to the second class constraints $`\pi _{\theta i}+\frac{i}{2}\theta _i^{}0`$, $`\pi _{\theta ^{}i}+\frac{i}{2}\theta _i0`$ \[$`\{\pi _{\theta i}+\frac{i}{2}\theta _i^{},\pi _{\theta ^{}j}+\frac{i}{2}\theta _j\}=i\delta _{ij}`$\]; $`\pi _{\theta i}`$ and $`\pi _{\theta ^{}i}`$ are then eliminated with the help of Dirac brackets
$$\{A,B\}{}_{}{}^{}=\{A,B\}i[\{A,\pi _{\theta i}+\frac{i}{2}\theta _i^{}\}\{\pi _{\theta ^{}i}+\frac{i}{2}\theta _i,B\}+\{A,\pi _{\theta ^{}i}+\frac{i}{2}\theta _i\}\{\pi _{\theta i}+\frac{i}{2}\theta _i^{},B\}],$$
(86)
so that the remaining Grassmann variables have the fundamental Dirac brackets \[which we will still denote $`\{.,.\}`$ for the sake of simplicity\]
$`\{\theta _i(\tau ),\theta _j(\tau )\}=\{\theta _i^{}(\tau ),\theta _j^{}(\tau )\}=0,`$ (87)
$`\{\theta _i(\tau ),\theta _j^{}(\tau )\}=i\delta _{ij}.`$ (88)
As in the free particle case of Section II, we obtain four primary constraints
$`_\mu `$ $`(\tau ,\stackrel{}{\sigma })=\rho _\mu (\tau ,\stackrel{}{\sigma })l_\mu (\tau ,\stackrel{}{\sigma })[T_{\tau \tau }(\tau ,\stackrel{}{\sigma })+{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\times `$ (90)
$`\sqrt{m_i^2\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })[\kappa _{i\stackrel{ˇ}{r}}(\tau )Q_iA_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })][\kappa _{i\stackrel{ˇ}{s}}(\tau )Q_iA_{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })]}]`$
$``$ $`z_{\stackrel{ˇ}{r}\mu }(\tau ,\stackrel{}{\sigma })\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\{T_{\tau \stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })+{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))[\kappa _{i\stackrel{ˇ}{s}}Q_iA_{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })]\}0,`$ (91)
where
$`T_{\tau \tau }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{1}{2}}({\displaystyle \frac{1}{\sqrt{\gamma }}}\pi ^{\stackrel{ˇ}{r}}g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}\pi ^{\stackrel{ˇ}{s}}{\displaystyle \frac{\sqrt{\gamma }}{2}}\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}\gamma ^{\stackrel{ˇ}{u}\stackrel{ˇ}{v}}F_{\stackrel{ˇ}{r}\stackrel{ˇ}{u}}F_{\stackrel{ˇ}{s}\stackrel{ˇ}{v}})(\tau ,\stackrel{}{\sigma }),`$ (92)
$`T_{\tau \stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`F_{\stackrel{ˇ}{s}\stackrel{ˇ}{t}}(\tau ,\stackrel{}{\sigma })\pi ^{\stackrel{ˇ}{t}}(\tau ,\stackrel{}{\sigma })=ϵ_{\stackrel{ˇ}{s}\stackrel{ˇ}{t}\stackrel{ˇ}{u}}\pi ^{\stackrel{ˇ}{t}}(\tau ,\stackrel{}{\sigma })B_{\stackrel{ˇ}{u}}(\tau ,\stackrel{}{\sigma })=`$ (93)
$`=`$ $`[\stackrel{}{\pi }(\tau ,\stackrel{}{\sigma })\times \stackrel{}{B}(\tau ,\stackrel{}{\sigma })]_{\stackrel{ˇ}{s}},`$ (94)
are the energy density and the Poynting vector respectively. We use the notation $`(\stackrel{}{\pi }\times \stackrel{}{B})_{\stackrel{ˇ}{s}}=(\stackrel{}{E}\times \stackrel{}{B})_{\stackrel{ˇ}{s}}`$ because it is consistent with $`ϵ_{\stackrel{ˇ}{s}\stackrel{ˇ}{t}\stackrel{ˇ}{u}}\pi ^{\stackrel{ˇ}{t}}B_{\stackrel{ˇ}{u}}`$ in the flat metric limit $`g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}\eta _{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}`$; in this limit $`T_{\tau \tau }\frac{1}{2}(\stackrel{}{E}^2+\stackrel{}{B}^2)`$.
This form of the constraint displays both the tangential $`z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })`$ and normal $`l^\mu (\tau ,\stackrel{}{\sigma })`$ components of the momentum $`\rho ^\mu (\tau ,\stackrel{}{\sigma })`$ conjugate to the embedding variables $`z^\mu (\tau ,\stackrel{}{\sigma })`$.
Again, being solved in terms of the momenta $`\rho _\mu (\tau ,\stackrel{}{\sigma })`$,these constraints are first class with exactly zero Poisson brackets ($`\{_\mu (\tau ,\stackrel{}{\sigma }),_\nu (\tau ,\stackrel{}{\sigma }^{^{}})\}=0`$) and their existence implies once again that the description of the system is independent of the choice of foliation.
Moreover, we have the (Lorentz scalar) primary constraints of the electromagnetic field connected with the gauge invariance of the action
$$\pi ^\tau (\tau ,\stackrel{}{\sigma })0.$$
(95)
From these constraints we construct the Dirac Hamiltonian. But first we must construct the canonical Hamiltonian $`H_C`$. The canonical Hamiltonian is
$`H_c`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\kappa _{i\stackrel{ˇ}{r}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau )+{\displaystyle d^3\sigma [\pi ^{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })_\tau A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma })\rho _\mu (\tau ,\stackrel{}{\sigma })z_\tau ^\mu (\tau ,\stackrel{}{\sigma })(\tau ,\stackrel{}{\sigma })]}=`$ (96)
$`=`$ $`{\displaystyle d^3\sigma [_{\stackrel{ˇ}{r}}(\pi ^{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })A_\tau (\tau ,\stackrel{}{\sigma }))A_\tau (\tau ,\stackrel{}{\sigma })\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })]}={\displaystyle d^3\sigma A_\tau (\tau ,\stackrel{}{\sigma })\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })},`$ (97)
after the elimination of a surface term.
Note that because of the $`\tau `$ and $`\stackrel{}{\sigma }`$ reparametrization invariance, $`H_C`$ nearly vanishes, except for the portion involving
$$\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })_{\stackrel{ˇ}{r}}\pi ^{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })\underset{i=1}{\overset{N}{}}Q_i\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )).$$
(98)
.
Thus the Dirac Hamiltonian is ($`\lambda ^\mu (\tau ,\stackrel{}{\sigma })`$ and $`\lambda _\tau (\tau ,\stackrel{}{\sigma })`$ are Dirac multipliers)
$$H_D=d^3\sigma [\lambda ^\mu (\tau ,\stackrel{}{\sigma })_\mu (\tau ,\stackrel{}{\sigma })+\lambda _\tau (\tau ,\stackrel{}{\sigma })\pi ^\tau (\tau ,\stackrel{}{\sigma })A_\tau (\tau ,\stackrel{}{\sigma })\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })].$$
(99)
The requirement that the primary constraints be $`\tau `$ independent ($`\{\pi ^\tau (\tau ,\stackrel{}{\sigma }),H_D\}0`$, $`\{^\mu (\tau ,\stackrel{}{\sigma }),H_D\}0`$) leads only to the Gauss’s law secondary constraint
$$\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })0.$$
(100)
Since the embedding variables $`z^\mu (\tau ,\stackrel{}{\sigma })`$ are the only configuration variables with Lorentz indices, the ten conserved generators of the Poincaré transformations are:
$`P^\mu `$ $`=`$ $`p_s^\mu ={\displaystyle d^3\sigma \rho ^\mu (\tau ,\stackrel{}{\sigma })},`$ (101)
$`J^{\mu \nu }`$ $`=`$ $`J_s^{\mu \nu }={\displaystyle d^3\sigma (z^\mu \rho ^\nu z^\nu \rho ^\mu )(\tau ,\stackrel{}{\sigma })},`$ (102)
(the subscript $`s`$ stands for hypersurface variable). From the first of these we obtain
$$\{z^\mu (\tau ,\stackrel{}{\sigma }),p_s^\nu \}=\eta ^{\mu \nu }.$$
(103)
We can restrict ourselves to foliations whose leaves are spacelike hyperplanes $`\mathrm{\Sigma }_H(\tau )`$ with constant timelike normal
$`b_\tau ^\mu (_\tau b_\tau ^\mu =0)`$, by imposing the following gauge fixings \[$`x_s^\mu `$ is an arbitrary origin\]
$$\zeta ^\mu (\tau ,\stackrel{}{\sigma })=z^\mu (\tau ,\stackrel{}{\sigma })x_s^\mu (\tau )b_{\stackrel{ˇ}{r}}^\mu (\tau )\sigma ^{\stackrel{ˇ}{r}}0.$$
(104)
In this expression $`b_{\stackrel{ˇ}{r}}^\mu (\tau )`$, $`\stackrel{ˇ}{r}=1,2,3`$ are three orthonormal vectors, such that the constant and future pointing normal to the hyperplane $`\mathrm{\Sigma }_H(\tau )`$ \[$`b_{\stackrel{ˇ}{A}}^\mu =(b_\tau ^\mu ,b_{\stackrel{ˇ}{r}}^\mu )`$ are orthonormal tetrads\] is
$$l^\mu (\tau ,\stackrel{}{\sigma })l^\mu =b_\tau ^\mu =\epsilon _{\alpha \beta \gamma }^\mu b_{\stackrel{ˇ}{1}}^\alpha (\tau )b_{\stackrel{ˇ}{2}}^\beta (\tau )b_{\stackrel{ˇ}{3}}^\gamma (\tau ).$$
(105)
Using the definitions of the vierbeins we obtain from the above gauge fixing the simplifications
$`z_{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })`$ $``$ $`b_{\stackrel{ˇ}{r}}^\mu (\tau ),`$ (106)
$`z_\tau ^\mu (\tau ,\stackrel{}{\sigma })`$ $``$ $`\dot{x}_s^\mu (\tau )+\dot{b}_{\stackrel{ˇ}{r}}^\mu (\tau )\sigma ^{\stackrel{ˇ}{r}},`$ (107)
$`g_{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })`$ $``$ $`\delta _{\stackrel{ˇ}{r}\stackrel{ˇ}{s}},\gamma ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })\delta ^{\stackrel{ˇ}{r}\stackrel{ˇ}{s}},\gamma (\tau ,\stackrel{}{\sigma })1,`$ (108)
as well as a natural decomposition of the Lorentz generators into orbital and spin portions
$`J_s^{\mu \nu }`$ $`=`$ $`x_s^\mu p_s^\nu x_s^\nu p_s^\mu +S_s^{\mu \nu },`$ (110)
$`S_s^{\mu \nu }=b_{\stackrel{ˇ}{r}}^\mu (\tau ){\displaystyle d^3\sigma \sigma ^{\stackrel{ˇ}{r}}\rho ^\nu (\tau ,\stackrel{}{\sigma })}b_{\stackrel{ˇ}{r}}^\nu (\tau ){\displaystyle d^3\sigma \sigma ^{\stackrel{ˇ}{r}}\rho ^\mu (\tau ,\stackrel{}{\sigma })}.`$
Here $`S_s^{\mu \nu }`$ is the spin part of the Lorentz generators.
These gauge fixings have the following Poisson brackets with the primary constraint $`_\mu (\tau ,\stackrel{}{\sigma })0`$
$$\{\zeta ^\mu (\tau ,\stackrel{}{\sigma }),_\nu (\tau ,\stackrel{}{\sigma })\}=\eta _\nu ^\mu \delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{}).$$
(111)
Therefore, we get a continuum set of second class constraints. They can be eliminated by forming the Dirac brackets
$$\{A,B\}^{}=\{A,B\}[\{A,\zeta ^\mu (\tau ,\stackrel{}{\sigma })\}\{_\mu (\tau ,\stackrel{}{\sigma }),B\}\{A,_\mu (\tau ,\stackrel{}{\sigma })\}\{\zeta ^\mu (\tau ,\stackrel{}{\sigma }),B\}].$$
(112)
For example, one finds that
$$\{x_s^\mu ,p_s^\nu \}^{}=\eta ^{\mu \nu }.$$
(113)
In this way the infinity of continuum hypersurface degrees of freedom $`z^\mu (\tau ,\stackrel{}{\sigma })`$, $`\rho _\mu (\tau ,\stackrel{}{\sigma })`$, are reduced to 20 : i) 8 are $`x_s^\mu (\tau )`$, $`p_s^\mu `$; ii) 12 are the 6 independent pairs of canonical variables hidden in $`b_{\stackrel{ˇ}{A}}^\mu `$ and $`S_s^{\mu \nu }=J_s^{\mu \nu }(x_s^\mu p_s^\nu x_s^\nu p_s^\mu )`$ \[they have the following brackets consistent with the orthonormality of the tetrads $`b_{\stackrel{ˇ}{A}}^\mu `$ : $`\{b_{\stackrel{ˇ}{A}}^\mu ,b_{\stackrel{ˇ}{B}}^\nu \}=0`$, $`\{S_s^{\mu \nu },b_{\stackrel{ˇ}{A}}^\rho \}=\eta ^{\rho \nu }b_{\stackrel{ˇ}{A}}^\mu \eta ^{\rho \mu }b_{\stackrel{ˇ}{A}}^\nu `$, $`\{S_s^{\mu \nu },S_s^{\alpha \beta }\}=C_{\gamma \delta }^{\mu \nu \alpha \beta }S_s^{\gamma \delta }`$, where $`C_{\gamma \delta }^{\mu \nu \alpha \delta }=\delta _\gamma ^\nu \delta _\delta ^\alpha \eta ^{\mu \beta }+\delta _\gamma ^\mu \delta _\delta ^\beta \eta ^{\nu \alpha }\delta _\gamma ^\nu \delta _\delta ^\beta \eta ^{\mu \alpha }\delta _\gamma ^\mu \delta _\delta ^\alpha \eta ^{\nu \beta }`$ are the structure constants of the Lorentz algebra\].
It can be shown that the following 10 first class constraints survive at the level of the Dirac brackets \[they are 10 combinations of the primary constraints whose gauge freedom is not fixed by the gauge fixings (104)\]
$`\stackrel{~}{^\mu }(\tau )`$ $`=`$ $`{\displaystyle d^3\sigma ^\mu (\tau ,\stackrel{}{\sigma })}=`$ (114)
$`=`$ $`p_s^\mu l^\mu [{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+`$ (115)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma [\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma })]+`$ (116)
$`+`$ $`b_{\stackrel{ˇ}{r}}^\mu (\tau )\left[{\displaystyle \underset{i=1}{\overset{N}{}}}[\kappa _{i\stackrel{ˇ}{r}}(\tau )Q_iA_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\eta }_i(\tau ))]+{\displaystyle d^3\sigma [\stackrel{}{\pi }\times \stackrel{}{B}]_{\stackrel{ˇ}{r}}(\tau ,\stackrel{}{\sigma })}\right]0,`$ (117)
and
$`\stackrel{~}{}^{\mu \nu }(\tau )`$ $`=`$ $`b_{\stackrel{ˇ}{r}}^\mu (\tau ){\displaystyle d^3\sigma \sigma ^{\stackrel{ˇ}{r}}^\nu (\tau ,\stackrel{}{\sigma })}b_{\stackrel{ˇ}{r}}^\nu (\tau ){\displaystyle d^3\sigma \sigma ^{\stackrel{ˇ}{r}}^\mu (\tau ,\stackrel{}{\sigma })}=`$ (118)
$`=`$ $`S_s^{\mu \nu }[b_{\stackrel{ˇ}{r}}^\mu (\tau )b_\tau ^\nu b_{\stackrel{ˇ}{r}}^\nu (\tau )b_\tau ^\mu ][{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^{\stackrel{ˇ}{r}}(\tau )\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+`$ (119)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma \sigma ^{\stackrel{ˇ}{r}}[\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma })]`$ (120)
$``$ $`[b_{\stackrel{ˇ}{r}}^\mu (\tau )b_{\stackrel{ˇ}{s}}^\nu (\tau )b_{\stackrel{ˇ}{r}}^\nu (\tau )b_{\stackrel{ˇ}{s}}^\mu (\tau )][{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^{\stackrel{ˇ}{r}}(\tau )[\kappa _i^{\stackrel{ˇ}{s}}(\tau )Q_iA^{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))]+`$ (121)
$`+`$ $`{\displaystyle }d^3\sigma \sigma ^{\stackrel{ˇ}{r}}[\stackrel{}{\pi }\times \stackrel{}{B}]^{_{\stackrel{ˇ}{s}}}(\tau ,\stackrel{}{\sigma })]0.`$ (122)
These constraints say that $`p_s^\mu `$ coincides with the total 4-momentum of the isolated system and that $`S_s^{\mu \nu }`$ is determined by its spin tensor.
The Dirac Hamiltonian becomes $`H_D=\stackrel{~}{\lambda }_\mu (\tau )\stackrel{~}{}^\mu (\tau )+\stackrel{~}{\lambda }_{\mu \nu }(\tau )\stackrel{~}{}^{\mu \nu }(\tau )+d^3\sigma [\lambda _\tau (\tau ,\stackrel{}{\sigma })\pi ^\tau (\tau ,\stackrel{}{\sigma })A_\tau (\tau ,\stackrel{}{\sigma })\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })]`$, where, due to the associated Hamilton equations, the new Dirac multipliers have the following interpretation: $`\stackrel{~}{\lambda }^\mu (\tau )\stackrel{}{=}\dot{x}_s^\mu (\tau )`$, $`\stackrel{~}{\lambda }^{\mu \nu }(\tau )=\stackrel{~}{\lambda }^{\nu \mu }(\tau )\stackrel{}{=}\frac{1}{2}[\dot{b}_{\stackrel{ˇ}{r}}^\mu (\tau )b_{\stackrel{ˇ}{r}}^\nu (\tau )b_{\stackrel{ˇ}{r}}^\mu (\tau )\dot{b}_{\stackrel{ˇ}{r}}^\nu (\tau )]`$.
Restricting our considerations to configurations with $`p_s^2>0`$, we make a further canonical reduction to the special foliation whose hyperplanes are orthogonal to $`p_s^\mu .`$ These hyperplanes are intrinsically determined by the system itself and are called the Wigner hyperplanes $`\mathrm{\Sigma }_W(\tau )`$. They can be identified by requiring the gauge fixings $`b_{\stackrel{ˇ}{A}}^\mu (\tau )L^\mu {}_{\nu =A}{}^{}(p_s,\stackrel{}{p}_s)`$ for the constraints $`\stackrel{~}{}^{\mu \nu }(\tau )0`$, where $`L^\mu {}_{\nu }{}^{}(p_s,\stackrel{}{p}_s)`$ is the standard Wigner boost for timelike Poincaré orbits. This implies $`l^\mu =b_\tau ^\mu p_s^\mu /\sqrt{p_s^2}`$.
The rest frame form of a timelike fourvector $`p^\mu `$ is $`\stackrel{}{p}{}_{}{}^{\mu }=\eta \sqrt{p^2}(1;\stackrel{}{0})=\eta ^{\mu o}\eta \sqrt{p^2}`$, $`\stackrel{}{p}{}_{}{}^{2}=p^2`$, where $`\eta =signp^0`$. Since we restricted ourselves to positive energy particles, $`\eta _i=+1`$, we shall put $`\eta =1`$. The standard Wigner boost transforming $`\stackrel{}{p}^\mu `$ into $`p^\mu `$ is
$`L^\mu {}_{\nu }{}^{}(p,\stackrel{}{p})`$ $`=`$ $`ϵ_\nu ^\mu (u(p))=`$ (123)
$`=`$ $`\eta _\nu ^\mu +2{\displaystyle \frac{p^\mu \stackrel{}{p}_\nu }{p^2}}{\displaystyle \frac{(p^\mu +\stackrel{}{p}^\mu )(p_\nu +\stackrel{}{p}_\nu )}{p\stackrel{}{p}+p^2}}=`$ (124)
$`=`$ $`\eta _\nu ^\mu +2u^\mu (p)u_\nu (\stackrel{}{p}){\displaystyle \frac{(u^\mu (p)+u^\mu (\stackrel{}{p}))(u_\nu (p)+u_\nu (\stackrel{}{p}))}{1+u^0(p)}},`$ (125)
$`\nu =0`$ $`ϵ_0^\mu (u(p))=u^\mu (p)=p^\mu /\sqrt{p^2},`$ (127)
$`\nu =r`$ $`ϵ_r^\mu (u(p))=(u_r(p);\delta _r^i{\displaystyle \frac{u^i(p)u_r(p)}{1+u^0(p)}}).`$ (128)
The inverse of $`L^\mu {}_{\nu }{}^{}(p,\stackrel{}{p})`$ is $`L^\mu {}_{\nu }{}^{}(\stackrel{}{p},p)`$, the standard boost to the system rest frame, defined by
$$L^\mu {}_{\nu }{}^{}(\stackrel{}{p},p)=L_\nu {}_{}{}^{\mu }(p,\stackrel{}{p})=L^\mu {}_{\nu }{}^{}(p,\stackrel{}{p})|_{\stackrel{}{p}\stackrel{}{p}}.$$
(129)
We also use these boosts to define the following vierbeins \[the $`ϵ_r^\mu (u(p))`$ ’s are also called polarization vectors; the indices $`r,s`$ will be used for $`A=`$1,2,3 and $`\overline{o}`$ for $`A`$=0\]
$`ϵ_A^\mu (u(p))=L^\mu {}_{A}{}^{}(p,\stackrel{}{p}),`$ (130)
$`ϵ_\mu ^A(u(p))=L^A{}_{\mu }{}^{}(\stackrel{}{p},p)=\eta ^{AB}\eta _{\mu \nu }ϵ_B^\nu (u(p)),`$ (131)
(132)
$`ϵ_\mu ^{\overline{o}}(u(p))=\eta _{\mu \nu }ϵ_o^\nu (u(p))=u_\mu (p),`$ (133)
$`ϵ_\mu ^r(u(p))=\delta ^{rs}\eta _{\mu \nu }ϵ_r^\nu (u(p))=(\delta ^{rs}u_s(p);\delta _j^r\delta ^{rs}\delta _{jh}{\displaystyle \frac{u^h(p)u_s(p)}{1+u^o(p)}}),`$ (134)
$`ϵ_o^A(u(p))=u_A(p),`$ (135)
which satisfy
$`ϵ_\mu ^A(u(p))ϵ_A^\nu (u(p))=\eta _\nu ^\mu ,`$ (136)
$`ϵ_\mu ^A(u(p))ϵ_B^\mu (u(p))=\eta _B^A,`$ (137)
$`\eta ^{\mu \nu }=ϵ_A^\mu (u(p))\eta ^{AB}ϵ_B^\nu (u(p))=u^\mu (p)u^\nu (p){\displaystyle \underset{r=1}{\overset{3}{}}}ϵ_r^\mu (u(p))ϵ_r^\nu (u(p)),`$ (138)
$`\eta _{AB}=ϵ_A^\mu (u(p))\eta _{\mu \nu }ϵ_B^\nu (u(p)),`$ (139)
$`p_\alpha {\displaystyle \frac{}{p_\alpha }}ϵ_A^\mu (u(p))=p_\alpha {\displaystyle \frac{}{p_\alpha }}ϵ_\mu ^A(u(p))=0.`$ (140)
With the Wigner rotation corresponding to the Lorentz transformation $`\mathrm{\Lambda }`$ being
$`R^\mu {}_{\nu }{}^{}(\mathrm{\Lambda },p)`$ $`=`$ $`[L(\stackrel{}{p},p)\mathrm{\Lambda }^1L(\mathrm{\Lambda }p,\stackrel{}{p})]^\mu {}_{\nu }{}^{}=\left(\begin{array}{cc}1& 0\\ 0& R^i{}_{j}{}^{}(\mathrm{\Lambda },p)\end{array}\right),`$ (143)
$`R^i{}_{j}{}^{}(\mathrm{\Lambda },p)`$ $`=`$ $`(\mathrm{\Lambda }^1)^i{}_{j}{}^{}{\displaystyle \frac{(\mathrm{\Lambda }^1)^i{}_{o}{}^{}p_{\beta }^{}(\mathrm{\Lambda }^1)^\beta _j}{p_\rho (\mathrm{\Lambda }^1)^\rho {}_{o}{}^{}+\eta \sqrt{p^2}}}`$ (145)
$``$ $`{\displaystyle \frac{p^i}{p^o+\eta \sqrt{p^2}}}[(\mathrm{\Lambda }^1)^o{}_{j}{}^{}{\displaystyle \frac{((\mathrm{\Lambda }^1)^o{}_{o}{}^{}1)p_\beta (\mathrm{\Lambda }^1)^\beta _j}{p_\rho (\mathrm{\Lambda }^1)^\rho {}_{o}{}^{}+\eta \sqrt{p^2}}}].`$ (146)
we have that the polarization vectors transform under the Poincaré transformations $`(a,\mathrm{\Lambda })`$ in the following way:
$$ϵ_r^\mu (u(\mathrm{\Lambda }p))=(R^1)_r{}_{}{}^{s}\mathrm{\Lambda }_{}^{\mu }{}_{\nu }{}^{}ϵ_{s}^{\nu }(u(p))\text{)}.$$
(147)
These boosts can be used to obtain the further canonical reduction referred to above. This takes place in two steps: i) firstly, one boosts to the rest frame the variables $`b_{\stackrel{ˇ}{A}}^\mu `$, $`S_s^{\mu \nu }`$, with the standard Wigner boost $`L^\mu {}_{\nu }{}^{}(p_s,\underset{s}{\overset{}{p}})`$ for timelike Poincaré orbits; ii) then, one adds the gauge-fixings $`b_{\stackrel{ˇ}{A}}^\mu L^\mu {}_{A}{}^{}(p_s,\underset{s}{\overset{}{p}})0`$ and goes to Dirac brackets. It can be shown that after this special gauge fixing the Lorentz scalar 3-indices $`\stackrel{ˇ}{r}`$ become Wigner spin 1 3-indices $`r`$. Therefore, we get a rest-frame instant form of dynamics with Wigner covariance.
The Lorentz generators become $`J_s^{\mu \nu }=\stackrel{~}{x}_s^\mu p_s^\nu \stackrel{~}{x}_s^\nu p_s^\mu +\stackrel{~}{S}_s^{\mu \nu }`$ with $`\stackrel{~}{S}_s^{\mu \nu }`$ given in Eq.(59) of Ref.. We now get $`\stackrel{~}{}^{\mu \nu }(\tau )0`$, i.e. $`S_s^{\mu \nu }`$ is forced to coincide with the spin tensor of the isolated system.
If we define the rest-frame spin tensor
$`\overline{S}_s^{AB}`$ $`=`$ $`ϵ_\mu ^A(u(p_s))ϵ_\nu ^B(u(p_s))S_s^{\mu \nu }[\eta _{\stackrel{ˇ}{r}}^A\eta _\tau ^B\eta _{\stackrel{ˇ}{r}}^B\eta _\tau ^A][{\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma \sigma ^{\stackrel{ˇ}{r}}[\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma })+`$ (148)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^{\stackrel{ˇ}{r}}(\tau )\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}]`$ (149)
$``$ $`[\eta _{\stackrel{ˇ}{r}}^A\eta _{\stackrel{ˇ}{s}}^B\eta _{\stackrel{ˇ}{r}}^B\eta _{\stackrel{ˇ}{s}}^A][{\displaystyle }d^3\sigma \sigma ^{\stackrel{ˇ}{r}}[\stackrel{}{\pi }\times \stackrel{}{B}]_{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\sigma })+`$ (150)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^{\stackrel{ˇ}{r}}(\tau )[\kappa _i^{\stackrel{ˇ}{s}}(\tau )Q_iA^{\stackrel{ˇ}{s}}(\tau ,\stackrel{}{\eta }_i(\tau ))]],`$ (151)
$`\overline{S}_s^{rs}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}(\eta _i^r(\tau )[\kappa _i^s(\tau )Q_iA^s(\tau ,\stackrel{}{\eta }_i(\tau ))]`$ (153)
$``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^s(\tau )[\kappa _i^r(\tau )Q_iA^s(\tau ,\stackrel{}{\eta }_i(\tau ))])+`$ (154)
$`+`$ $`{\displaystyle d^3\sigma \left(\sigma ^r[\stackrel{}{\pi }\times \stackrel{}{B}]^s(\tau ,\stackrel{}{\sigma })\sigma ^s[\stackrel{}{\pi }\times \stackrel{}{B}]^r(\tau ,\stackrel{}{\sigma })\right)},`$ (155)
$`\overline{S}_s^{\overline{0}r}`$ $``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\eta _i^r(\tau )\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}`$ (157)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^3\sigma \sigma ^r[\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma })}.`$ (158)
it can be shown that the form of $`\stackrel{~}{S}_s^{\mu \nu }`$ implies that the rest-frame “external” Poincaré generators are
$`J_s^{ij}`$ $`:`$ $`=\stackrel{~}{x}_s^ip_s^j\stackrel{~}{x}_s^jp_s^i+\delta ^{ir}\delta ^{js}\overline{S}_s^{rs},`$ (159)
$`J_s^{oi}`$ $`:`$ $`=\stackrel{~}{x}_s^op_s^i\stackrel{~}{x}_s^ip_s^o{\displaystyle \frac{\delta ^{ir}\overline{S}_s^{rs}p_s^s}{p_s^o+\eta _s\sqrt{p_s^2}}}.`$ (160)
Only in this special gauge do we get the separation of a decoupled “external” canonical non-covariant center of mass described by the 4 pairs $`\stackrel{~}{x}_s^\mu (\tau )`$, $`p_s^\mu `$, of canonical variables ($`\{\stackrel{~}{x}_s^\mu ,p_s^\nu \}^{}=\eta ^{\mu \nu }`$) identifying the Wigner hyperplane $`\mathrm{\Sigma }_W(\tau )`$ \[see Eq.(59) of Ref. for the expression of $`\stackrel{~}{x}_s^\mu (\tau )`$ in terms of $`x_s^\mu `$ and of the spin tensor\] and of the “internal” Wigner-covariant canonical variables $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{}{\kappa }_i(\tau )`$, $`A_A(\tau ,\stackrel{}{\sigma })`$, $`\pi ^A(\tau ,\stackrel{}{\sigma })`$ living inside the Wigner hyperplane and with the Dirac brackets coinciding with the original Poisson brackets ( $`\{\eta _i^{\stackrel{ˇ}{r}}(\tau ),\kappa _{j\stackrel{ˇ}{s}}(\tau )\}^{}=\delta _{ij}\delta _{\stackrel{ˇ}{s}}^{\stackrel{ˇ}{r}}`$, $`\{A_{\stackrel{ˇ}{A}}(\tau ,\stackrel{}{\sigma }),\pi ^{\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma }^{^{}})\}=\eta _{\stackrel{ˇ}{A}}^{\stackrel{ˇ}{B}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{^{}})`$).
As shown in Ref. one can replace $`\stackrel{~}{x}_s^\mu ,p_s^\mu `$ with a new canonical basis $`T_s=p_s\stackrel{~}{x}_s/\sqrt{p_s^2}=p_sx_s/\sqrt{p_s^2}`$ (it is the Lorentz-invariant rest frame time), $`\epsilon _s=\sqrt{p_s^2},\stackrel{}{z}_s=\sqrt{p_s^2}(\stackrel{~}{\stackrel{}{x}}_s\stackrel{~}{x}_s^o\stackrel{}{p}_s/p_s^o),\stackrel{}{k}_s=\stackrel{}{u}(p_s)`$ with $`\stackrel{}{z}_s`$ having the same covariance of the Newton-Wigner position operator under the little group $`O(3)`$ of the timelike Poincaré orbits. The 3-position canonical variable $`\stackrel{}{z}_s/ϵ_s`$ is the classical background of this operator and describes the decoupled “external” 3-center of mass, whose 4-position is $`\stackrel{~}{x}_s^\mu `$.
In this special gauge there is no restriction on $`p_s^\mu `$: the four velocity $`u^\mu (p_s)=p_s^\mu /\sqrt{p_s^2}=l^\mu `$ describes the orientation of the Wigner hyperplane with respect to an arbitrary Lorentz frame.
We obtain the following form for the constraints $`\stackrel{~}{^\mu }(\tau )0`$:
$`\stackrel{~}{^\mu }(\tau )`$ $`=`$ $`{\displaystyle d^3\sigma ^\mu (\tau ,\stackrel{}{\sigma })}=p_s^\mu `$ (162)
$`u^\mu (p)({\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma [\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma })+`$
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2})`$ (163)
$``$ $`ϵ_r^\mu (u(p))({\displaystyle }d^3\sigma [\stackrel{}{\pi }\times \stackrel{}{B}]^r(\tau ,\stackrel{}{\sigma })+`$ (164)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^r)0.`$ (165)
Their projections along the normal and the tangents to the Wigner hyperplane are
$`(\tau )`$ $`=`$ $`u^\mu (p_s)\stackrel{~}{}_\mu (\tau )=`$ (166)
$`=`$ $`\sqrt{p_s^2}({\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+`$ (167)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma [\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma }))0,`$ (168)
$$\stackrel{}{}_p(\tau )=\underset{i=1}{\overset{N}{}}[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]+d^3\sigma [\stackrel{}{\pi }\times \stackrel{}{B}](\tau ,\stackrel{}{\sigma })0.$$
(169)
The first one gives the mass spectrum of the isolated field plus particle system, while the other three say that the total 3-momentum of the $`N`$ charged particles plus fields vanishes inside the Wigner hyperplane $`\mathrm{\Sigma }_W(\tau )`$. This condition is the rest-frame condition identifying the Wigner hyperplane as the rest frame of the isolated system.
The Dirac Hamiltonian is now $`H_D=\stackrel{~}{\lambda }^\mu (\tau )\stackrel{~}{}_\mu (\tau )=\lambda (\tau )(\tau )\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau )`$, with $`\lambda (\tau )\dot{x}_{s\mu }(\tau )u^\mu (p_s)`$, $`\lambda _r(\tau )\dot{x}_{s\mu }(\tau )ϵ_r^\mu (u(p_s))`$.
The two additional electromagnetic constraints are
$`\pi ^\tau (\tau ,\stackrel{}{\sigma })`$ $``$ $`0,`$ (170)
$`\mathrm{\Gamma }(\tau ,\stackrel{}{\sigma })`$ $``$ $`0.`$ (171)
In the rest-frame instant form of dynamics on the Wigner hyperplanes they are Lorentz scalar constraints \[$`A_\tau (\tau ,\stackrel{}{\sigma })`$ and $`\pi ^\tau (\tau ,\stackrel{}{\sigma })`$ are Lorentz scalars, while $`\stackrel{}{A}(\tau ,\stackrel{}{\sigma })`$ and $`\stackrel{}{\pi }(\tau ,\stackrel{}{\sigma })`$ are spin-1 Wigner 3-vectors\].
We now eliminate the electromagnetic gauge degrees of freedom by decomposing the above spin-one Wigner 3-vector canonical field variables into their transverse and longitudinal components
$`\stackrel{}{A}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma }){\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\stackrel{}{}\stackrel{}{A}(\tau ,\stackrel{}{\sigma }),`$ (172)
$`\stackrel{}{\pi }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma }){\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\stackrel{}{}\stackrel{}{\pi }(\tau ,\stackrel{}{\sigma })`$ (173)
$``$ $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )),`$ (174)
with $`\mathrm{\Delta }=\stackrel{}{}^2.`$
We re-express everything in terms of the Dirac observables: i) $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$, $`\{\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma }),\stackrel{ˇ}{\pi }_{}^s(\tau ,\stackrel{}{\sigma }^{^{}})\}=P_{}^{rs}(\stackrel{}{\sigma })\delta ^3(\stackrel{}{\sigma }\stackrel{}{\sigma }^{^{}})`$ \[$`P_{}^{rs}(\stackrel{}{\sigma })=\delta ^{rs}+\frac{^r^s}{\mathrm{}}`$\] for the electromagnetic field; ii) $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )=\stackrel{}{\kappa }_i(\tau )+Q_i\frac{\stackrel{}{}}{\mathrm{}}\stackrel{}{}\stackrel{}{A}(\tau ,\stackrel{}{\sigma })`$ for the particles \[they now become dressed with a Coulomb cloud\]; iii) $`\stackrel{ˇ}{\theta }_i^{}(\tau )`$, $`\stackrel{ˇ}{\theta }_i(\tau )`$, such that $`Q_i=e\theta _i^{}\theta _i=e\stackrel{ˇ}{\theta }_i^{}\stackrel{ˇ}{\theta }_i`$.
This is known as the Wigner-covariant rest-frame radiation gauge. Note that $`\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))=\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{}{A}_{}(\tau ,\stackrel{}{\eta }_i(\tau )).`$ Using the Gauss law constraint $`\stackrel{}{}\stackrel{}{\pi }(\tau ,\stackrel{}{\sigma })\mathrm{\Sigma }_iQ_i\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i)`$ with $`\frac{1}{\mathrm{\Delta }}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i)`$ $`=1/(4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i|)`$ and integrating by parts we separate out the Coulomb portion of the rest frame energy from the field energy integral. A similar procedure on the field momentum integral simplifies the rest frame condition (and the expression for the internal angular momentum in the next section). Thus we find that the reduced form of the 4 constraints is
$`(\tau )`$ $`=`$ $`ϵ_s\{{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+(\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau )))^2}+`$ (175)
$`+`$ $`{\displaystyle \underset{ij}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+{\displaystyle }d^3\sigma {\displaystyle \frac{1}{2}}[\stackrel{ˇ}{\stackrel{}{\pi }}_{}^2+\stackrel{ˇ}{\stackrel{}{B}}^2](\tau ,\stackrel{}{\sigma })\}=`$ (176)
$`=`$ $`ϵ_sM0,`$ (177)
$`\stackrel{}{}_p(\tau )`$ $`=`$ $`\mathrm{\Sigma }_i\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )+{\displaystyle d^3\sigma [\stackrel{ˇ}{\stackrel{}{\pi }}_{}\times \stackrel{ˇ}{\stackrel{}{B}}](\tau ,\stackrel{}{\sigma })}0,`$ (178)
where $`ϵ_s=\sqrt{p_s^2}`$ and $`M`$ is the invariant mass of the isolated system. The rest-frame condition does not depend any more on the interaction as it must be in an instant form of dynamics. This procedure not only extracts the Coulomb potential from field theory but also regularizes the Coulomb interaction due to the semiclassical property $`Q_i^2=0`$.
If we add the gauge fixing $`T_s\tau 0`$ we get $`\lambda (\tau )=1`$ and the Dirac Hamiltonian for the evolution in the rest-frame time is $`H_D=M\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau )`$ \[see the more accurate discussion after Eq.(207)\].
The embedding corresponding to the Wigner hyperplanes in the gauge $`T_s\tau `$ is $`z^\mu (\tau ,\stackrel{}{\sigma })=x_s^\mu (\tau )+ϵ_r^\mu (u(p_s))\sigma ^r`$ with the origin of the 3-coordinates given by $`x_s^\mu (\tau )=x_s^\mu (0)+u^\mu (p_s)\tau +ϵ_r^\mu (u(p_s))_0^\tau 𝑑\tau ^{^{}}\lambda _r(\tau ^{^{}})`$ since $`\dot{x}_s^\mu (\tau )=u^\mu (p_s)+ϵ_r^\mu (u(p_s))\lambda _r(\tau )`$ \[instead for the “external” center of mass we have $`\dot{\stackrel{~}{x}}_s^\mu (\tau )=u^\mu (p_s)`$\]. The final canonical variables are: i) $`\stackrel{}{z}_s`$, $`\stackrel{}{k}_s`$ (the decoupled “external” 3-center of mass); ii) $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$ (the particle variables); iii) $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$ (the transverse radiation field). They are still restricted by the three rest-frame conditions $`\stackrel{}{}_p(\tau )0`$.
### B The Energy-Momentum Tensor.
The Euler-Lagrange equations from the action (64)) are
$`({\displaystyle \frac{}{z^\mu }}_{\stackrel{ˇ}{A}}{\displaystyle \frac{}{z_{\stackrel{ˇ}{A}}^\mu }})(\tau ,\stackrel{}{\sigma })=\eta _{\mu \nu }_{\stackrel{ˇ}{A}}[\sqrt{g}T^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}z_{\stackrel{ˇ}{B}}^\nu ](\tau ,\stackrel{}{\sigma })\stackrel{}{=}\mathrm{\hspace{0.17em}0},`$ (179)
$`({\displaystyle \frac{L}{\eta _i^r}}{\displaystyle \frac{d}{d\tau }}{\displaystyle \frac{L}{\dot{\eta }_i^r}})(\tau )\stackrel{}{=}\mathrm{\hspace{0.17em}0},`$ (180)
$`({\displaystyle \frac{}{A_{\stackrel{ˇ}{A}}}}_{\stackrel{ˇ}{B}}{\displaystyle \frac{}{_{\stackrel{ˇ}{B}}A_{\stackrel{ˇ}{A}}}})(\tau ,\stackrel{}{\sigma })\stackrel{}{=}\mathrm{\hspace{0.17em}0},`$ (181)
where we introduced the total energy-momentum tensor \[$`\dot{\eta }_i^{\stackrel{ˇ}{A}}(\tau )=(1;\dot{\eta }_i^{\stackrel{ˇ}{r}}(\tau ))`$\]
$`T^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`[{\displaystyle \frac{2}{\sqrt{g}}}{\displaystyle \frac{\delta S}{\delta g_{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}}}](\tau ,\stackrel{}{\sigma })=`$ (182)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )){\displaystyle \frac{m_i\dot{\eta }_i^{\stackrel{ˇ}{A}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{B}}(\tau )}{\sqrt{g}\sqrt{g_{\stackrel{ˇ}{C}\stackrel{ˇ}{D}}\dot{\eta }_i^{\stackrel{ˇ}{C}}(\tau )\dot{\eta }_i^{\stackrel{ˇ}{D}}(\tau )}}}+`$ (183)
$`+`$ $`[F^{\stackrel{ˇ}{A}\stackrel{ˇ}{C}}F_{\stackrel{ˇ}{C}}{}_{}{}^{\stackrel{ˇ}{B}}+{\displaystyle \frac{1}{4}}g^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}F^{\stackrel{ˇ}{C}\stackrel{ˇ}{D}}F_{\stackrel{ˇ}{C}\stackrel{ˇ}{D}}](\tau ,\stackrel{}{\sigma }).`$ (184)
When $`_{\stackrel{ˇ}{A}}[\sqrt{g}z_{\stackrel{ˇ}{B}}^\mu ](\tau ,\stackrel{}{\sigma })=0`$ as happens on the Wigner hyperplanes in the gauge $`T_s\tau 0`$, $`\stackrel{}{\lambda }(\tau )=0`$, we get the conservation of the energy-momentum tensor $`T^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}(\tau ,\stackrel{}{\sigma })`$, i.e. $`_{\stackrel{ˇ}{A}}T^{\stackrel{ˇ}{A}\stackrel{ˇ}{B}}\stackrel{}{=}\mathrm{\hspace{0.17em}0}`$. Otherwise there is compensation coming from the dynamics of the hypersurface.
On the Wigner hyperplanes the energy-momentum tensor becomes
$`T^{\tau \tau }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+{\displaystyle \frac{1}{2}}[\stackrel{}{\pi }^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma }),`$ (185)
$`T^{r\tau }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))[\kappa _i^r(\tau )Q_iA^r(\tau ,\stackrel{}{\eta }_i(\tau ))]+[\stackrel{}{\pi }\times \stackrel{}{B}](\tau ,\stackrel{}{\sigma }),`$ (186)
$`T^{rs}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )){\displaystyle \frac{[\kappa _i^r(\tau )Q_iA^r(\tau ,\stackrel{}{\eta }_i(\tau ))][\kappa _i^s(\tau )Q_iA^s(\tau ,\stackrel{}{\eta }_i(\tau ))]}{\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}}}`$ (187)
$``$ $`\left[{\displaystyle \frac{1}{2}}\delta ^{rs}[\stackrel{}{\pi }^2+\stackrel{}{B}^2][\pi ^r\pi ^s+B^rB^s]\right](\tau ,\stackrel{}{\sigma }).`$ (188)
Finally, after the canonical reduction, which eliminates the electromagnetic gauge degrees of freedom and the choice of the gauge $`T_s\tau `$, we get
$`T^{\tau \tau }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\sqrt{m_i^2+[\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+`$ (189)
$`+`$ $`{\displaystyle \frac{1}{2}}[\left(\stackrel{ˇ}{\stackrel{}{\pi }}_{}+{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\right)^2+\stackrel{}{B}^2](\tau ,\stackrel{}{\sigma }),`$ (190)
$`T^{r\tau }(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))[\stackrel{ˇ}{\kappa }_i^r(\tau )Q_i\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\eta }_i(\tau ))]+`$ (191)
$`+`$ $`[\left(\stackrel{ˇ}{\stackrel{}{\pi }}_{}+{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))\right)\times \stackrel{}{B}](\tau ,\stackrel{}{\sigma }),`$ (192)
$`T^{rs}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )){\displaystyle \frac{[\stackrel{ˇ}{\kappa }_i^r(\tau )Q_i\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\eta }_i(\tau ))][\stackrel{ˇ}{\kappa }_i^s(\tau )Q_i\stackrel{ˇ}{A}_{}^s(\tau ,\stackrel{}{\eta }_i(\tau ))]}{\sqrt{m_i^2+[\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}}}`$ (193)
$``$ $`[{\displaystyle \frac{1}{2}}\delta ^{rs}[(\stackrel{ˇ}{\stackrel{}{\pi }}_{}+{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )))^2+\stackrel{}{B}^2]`$ (194)
$``$ $`[(\stackrel{ˇ}{\stackrel{}{\pi }}_{}+{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )))^r(\stackrel{ˇ}{\stackrel{}{\pi }}_{}+{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )))^s+`$ (195)
$`+`$ $`B^rB^s]](\tau ,\stackrel{}{\sigma }).`$ (196)
## IV Internal Poincaré Algebra and Equations of Motion for the Electromagnetic Field and the N Charged Particles.
In this Section we first build a realization of the Poincaré algebra inside the Wigner hyperplane using results of the previous Section. Then, by identifying the Lorentz scalar rest-frame time $`T_s`$ with the invariant time $`\tau `$ labeling the hypersurfaces $`\mathrm{\Sigma }(\tau )`$, we arrive at the Dirac Hamiltonian $`H_D=M\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau )`$, in which the only gauge freedom left is the one associated with the rest-frame condition. We then obtain the Hamilton and Lagrange equations for fields and particles. Then we describe how to find the canonical “internal” center of mass $`\stackrel{}{q}_+`$ for fields and particles on the Wigner hyperplane. The natural gauge fixing to the rest-frame conditions $`\stackrel{}{}_p(\tau )0`$ are $`\stackrel{}{q}_+0`$: they imply $`\stackrel{}{\lambda }(\tau )=0`$ and the decoupling of the “internal” center of mass from the “internal” relative motions. In this way only the “external” decoupled 3-center of mass $`\stackrel{}{z}_s`$ remains (the Newton-Wigner-like 3-position which replaces the 4-center of mass $`\stackrel{~}{x}_s^\mu `$ in the gauge $`T_s\tau `$). A property of the particle accelerations of any order, which will be needed in the next Section, is derived.
### A Internal Poincaré Algebra
In the rest-frame instant form of the dynamics there is another realization of the Poincaré algebra besides the “external” one given in Eq.(160). This is the “internal” realization built in terms of the variables living inside each Wigner hyperplane. The associated generators of this internal Poincaré group are given by (for positive energies in the Wigner-covariant rest frame radiation gauge)
$`𝒫_{(int)}^\tau `$ $`=`$ $`M={\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+(\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau )))^2}+`$ (197)
$`+`$ $`{\displaystyle \underset{ij}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+{\displaystyle d^3\sigma \frac{1}{2}[\stackrel{ˇ}{\stackrel{}{\pi }}_{}^2+\stackrel{ˇ}{\stackrel{}{B}}^2](\tau ,\stackrel{}{\sigma })},`$ (198)
$`\stackrel{}{𝒫}_{(int)}`$ $`=`$ $`\stackrel{}{}_p=\stackrel{ˇ}{\stackrel{}{\kappa }}_+(\tau )+{\displaystyle d^3\sigma [\stackrel{ˇ}{\stackrel{}{\pi }}_{}\times \stackrel{ˇ}{\stackrel{}{B}}](\tau ,\stackrel{}{\sigma })}0,`$ (200)
$`𝒥_{(int)}^r`$ $`=`$ $`\epsilon ^{rst}\overline{S}_s^{st}={\displaystyle \underset{i=1}{\overset{N}{}}}(\stackrel{}{\eta }_i(\tau )\times \stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau ))^r+{\displaystyle }d^3\sigma (\stackrel{}{\sigma }\times [\stackrel{ˇ}{\stackrel{}{\pi }}_{}\times \stackrel{ˇ}{\stackrel{}{B}}]^r(\tau ,\stackrel{}{\sigma }),`$ (202)
$`𝒦_{(int)}^r`$ $`=`$ $`\overline{S}_s^{\overline{o}r}=\overline{S}_s^{r\overline{o}}={\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{}{\eta }_i(\tau )\sqrt{m_i^2+[\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))]^2}+`$ (203)
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}[{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_iQ_j[{\displaystyle \frac{1}{\mathrm{}_{\stackrel{}{\eta }_j}}}{\displaystyle \frac{}{\eta _j^r}}c(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau ))\eta _j^r(\tau )c(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau ))]+`$ (204)
$`+Q_i`$ $`{\displaystyle }d^3\sigma \stackrel{ˇ}{\pi }_{}^r(\tau ,\stackrel{}{\sigma })c(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))]{\displaystyle \frac{1}{2}}{\displaystyle }d^3\sigma \sigma ^r(\stackrel{ˇ}{\stackrel{}{\pi }}_{}^2+\stackrel{ˇ}{\stackrel{}{B}}^2)(\tau ,\stackrel{}{\sigma }),`$ (205)
in which $`\stackrel{ˇ}{\stackrel{}{\kappa }}_+(\tau )=\mathrm{\Sigma }_i\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$ and $`c(\stackrel{}{\eta }_i\stackrel{}{\eta }_j):=1/(4\pi |\stackrel{}{\eta }_j\stackrel{}{\eta }_i|)`$. The latter two Lorentz generators are determined as the components of the spin tensor $`\overline{S}_s^{AB}`$ defined in Eq.(158) inside each Wigner hyperplane.
The Dirac Hamiltonian is
$$H_D=\lambda (\tau )(\tau )\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau ).$$
(206)
As already said, if we add the gauge-fixing
$$\chi =T_s\tau 0,T_s\frac{p_s\stackrel{~}{x}_s}{\sqrt{p_s^2}}=\frac{p_sx_s}{\sqrt{p_s^2}},$$
(207)
implying that the Lorentz scalar parameter $`\tau `$ labelling the leaves of the foliation of Minkowski spacetime with Wigner hyperplanes coincides with the rest-frame time $`T_s`$ of the decoupled point particle clock (the “external” center of mass) $`\stackrel{~}{x}_s^\mu `$, its conservation in $`\tau `$ will imply $`\lambda (\tau )=1`$ so that, after taking the Dirac brackets associated with the second class constraints $`ϵ_sM0`$ and $`T_s\tau 0`$ (this eliminates $`T_s`$ and $`ϵ_s`$), the final Dirac Hamiltonian in this gauge would be $`H_D=\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau )`$. However, if we wish to reintroduce the evolution in $`\tau T_s`$ in this frozen phase space \[containing the canonical variables $`\stackrel{}{z}_s`$, $`\stackrel{}{k}_s`$, $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$\] we must use the Hamiltonian
$$H_D=M\stackrel{}{\lambda }(\tau )\stackrel{}{}_p(\tau ),$$
(208)
because $`M=𝒫_{(int)}^\tau `$ is the invariant mass and the “internal” energy generator of the isolated system \[it is like with the frozen Hamilton-Jacobi theory, in which the time evolution can be reintroduced by using the energy generator of the Poincaré group as Hamiltonian\].
The only remaining first class constraints are the rest-frame conditions $`\stackrel{}{}_p(\tau )0`$. The Dirac multipliers $`\stackrel{}{\lambda }(\tau )`$ describe the remaining gauge freedom on the location of the “internal” center of mass on the Wigner hyperplanes. In the next Subsection we will study the natural gauge fixings for these first class constraints. After this final canonical reduction the isolated system will be described by the decoupled “external” 3-center-of-mass canonical variables $`\stackrel{}{z}_s`$, $`\stackrel{}{k}_s`$ and by relative Wigner-covariant degrees of freedom on the Wigner hyperplane, with the invariant mass $`M`$ as the Hamiltonian for the evolution in $`\tau T_s`$.
### B The equations of motion for particles and fields.
The Hamilton-Dirac equations associated to the previous Hamiltonian are
$`\dot{\stackrel{}{\eta }}_i(\tau )`$ $`\stackrel{}{=}`$ $`{\displaystyle \frac{\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))}{\sqrt{m_i^2+(\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau )))^2}}}\stackrel{}{\lambda }(\tau ),`$ (209)
$`\dot{\stackrel{ˇ}{\stackrel{}{\kappa }}}_i(\tau )\stackrel{}{=}`$ $`{\displaystyle \underset{ki}{}}{\displaystyle \frac{Q_iQ_k(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau ))}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau )^3}}+`$ (210)
$`+`$ $`Q_i(\dot{\eta }_i^u(\tau )+\lambda ^u(\tau )){\displaystyle \frac{}{\stackrel{}{\eta }_i}}\stackrel{ˇ}{A}_{}^u(\tau ,\stackrel{}{\eta }_i(\tau ))],`$ (211)
$`\stackrel{ˇ}{\stackrel{}{\kappa }}_+(\tau )`$ $`+`$ $`{\displaystyle d^3\sigma [\stackrel{ˇ}{\stackrel{}{\pi }}_{}\times \stackrel{ˇ}{\stackrel{}{B}}](\tau ,\stackrel{}{\sigma })}0.`$ (212)
in which $`\stackrel{}{=}`$ means evaluated on the equations of motion.
The Hamilton-Dirac equations for the fields are
$`\dot{\stackrel{ˇ}{A}}_r(\tau ,\stackrel{}{\sigma })\stackrel{}{=}\stackrel{ˇ}{\pi }_r(\tau ,\stackrel{}{\sigma })[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{A}_r(\tau ,\stackrel{}{\sigma }),`$ (213)
$`\dot{\stackrel{ˇ}{\pi }}_{}^r(\tau ,\stackrel{}{\sigma })`$ $`\stackrel{}{=}`$ $`\mathrm{\Delta }\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{\pi }_{}^r(\tau ,\stackrel{}{\sigma })+`$ (214)
$``$ $`{\displaystyle \underset{i}{}}Q_iP_{}^{rs}(\stackrel{}{\sigma })\dot{\eta }_i^s(\tau )\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )).`$ (215)
The associated Lagrangian, obtained by means of an inverse Legendre transformation , is
$`L_R(\tau )`$ $`=`$ $`\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau ){\displaystyle d^3\sigma \dot{\stackrel{ˇ}{\stackrel{}{A}}}_{}(\tau ,\stackrel{}{\sigma })\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })}H_R(\tau )=`$ (216)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\left[m_i\sqrt{1(\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau ))^2}+Q_i[\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau )]\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))\right]+`$ (217)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+`$ (218)
$`+`$ $`{\displaystyle d^3\sigma [\frac{(\dot{\stackrel{ˇ}{\stackrel{}{A}}}_{}+[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{\stackrel{}{A}}_{})^2}{2}\frac{\stackrel{ˇ}{\stackrel{}{B}}^2}{2}](\tau ,\stackrel{}{\sigma })}.`$ (219)
Here $`\stackrel{}{\lambda }(\tau )`$ is now interpreted as a non-linear Lagrange multiplier needed to get the rest-frame conditions $`\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}0`$. Its Euler-Lagrange equations$`\frac{d}{dt}\frac{L_R}{\dot{\eta }_{ir}}\stackrel{}{=}\frac{L_R}{\eta _{ir}};\frac{L_R}{\stackrel{}{\lambda }}\stackrel{}{=}\mathrm{\hspace{0.17em}0}`$ yield
$`{\displaystyle \frac{d}{d\tau }}[m_i{\displaystyle \frac{\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau )}{\sqrt{1(\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau ))^2}}}+Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))]\stackrel{}{=}`$ (220)
$`\stackrel{}{=}`$ $`{\displaystyle \underset{ki}{}}{\displaystyle \frac{Q_iQ_k(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau )}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau )^3}}+Q_i(\dot{\eta }_i^u(\tau )+\lambda ^u(\tau )){\displaystyle \frac{}{\stackrel{}{\eta }_i}}\stackrel{ˇ}{A}_{}^u(\tau ,\stackrel{}{\eta }_i(\tau )),`$ (221)
$``$ $`\ddot{\stackrel{ˇ}{A}}_{}^r(\tau ,\stackrel{}{\sigma }){\displaystyle \frac{d}{d\tau }}\{[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })\}\stackrel{}{=}`$ (224)
$`\stackrel{}{=}\mathrm{\Delta }\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })+[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\{\dot{\stackrel{ˇ}{A}}_{}^r(\tau ,\stackrel{}{\sigma })+[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })\}`$
$``$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_iP_{}^{rs}(\stackrel{}{\sigma })[\dot{\eta }_i^s(\tau )+\lambda ^s(\tau )]\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )),`$ (225)
and \[these are the Lagrangian rest-frame conditions\]
$`{\displaystyle \underset{i=1}{\overset{N}{}}}[m_i{\displaystyle \frac{\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau )}{\sqrt{1(\dot{\stackrel{}{\eta }}(\tau )+\stackrel{}{\lambda }(\tau ))^2}}}+Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))]+`$ (226)
$`+{\displaystyle d^3\sigma \underset{r}{}[(\stackrel{}{}\stackrel{ˇ}{A}_{}^r)(\dot{\stackrel{ˇ}{A}}_{}^r+[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{A}_{}^r)](\tau ,\stackrel{}{\sigma })}\stackrel{}{=}\mathrm{\hspace{0.17em}0}.`$ (227)
The Lagrangian expression for the conserved invariant mass $`M=𝒫_{(int)}^\tau `$ is
$`E_{rel}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{m_i}{\sqrt{1(\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau ))^2}}}+{\displaystyle \underset{i>j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+`$ (228)
$`+`$ $`{\displaystyle d^3\sigma \frac{1}{2}[\stackrel{ˇ}{\stackrel{}{E}}_{}^2+\stackrel{ˇ}{\stackrel{}{B}}^2](\tau ,\stackrel{}{\sigma })}=const.`$ (229)
Eq.(225) may be rewritten as
$`{\displaystyle \frac{d}{d\tau }}(m_i{\displaystyle \frac{\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau )}{\sqrt{1+(\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau ))^2}}})`$ (230)
$`\stackrel{}{=}`$ $`{\displaystyle \underset{ki}{}}{\displaystyle \frac{Q_iQ_k(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau ))}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau )^3}}+`$ (231)
$`+`$ $`Q_i[\stackrel{ˇ}{\stackrel{}{E}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))+(\dot{\stackrel{}{\eta }}_i(\tau )+\stackrel{}{\lambda }(\tau ))\times \stackrel{ˇ}{\stackrel{}{B}}(\tau ,\stackrel{}{\eta }_i(\tau ))],`$ (232)
where the notation $`\stackrel{ˇ}{E}_{}^r=\stackrel{ˇ}{\dot{A}}_{}^r[\stackrel{}{\lambda }(\tau )\stackrel{}{}]\stackrel{ˇ}{A}_{}^r=\stackrel{ˇ}{\pi }_{}^r`$ has been introduced.
Eqs.(232) and (225) are the rest-frame analogues of the usual equations for charged particles in an external electromagnetic field and of the electromagnetic field with external particle sources in which both particles and electromagnetic field are dynamical. Eq.(227) defines the rest frame by using the total (Wigner spin 1) 3-momentum of the isolated system formed by the particles plus the electromagnetic field. Eq.(229 ) gives the constant invariant mass of the isolated system: the electromagnetic self-energy of the particles has been regularized by the Grassmann-valued electric charges \[$`Q_i^2=0`$\] so that the invariant mass is finite.
### C The “internal” center of mass and the last gauge fixing.
The rest-frame conditions $`\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}0`$ show that there are still 3 gauge degrees of freedom among the reduced canonical variables $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$ on each Wigner hyperplane. They correspond to our freedom in the choice of the point of the Wigner hyperplane that locates the “internal” 3-center of mass $`\stackrel{}{q}_+`$ of the isolated system. After the gauge fixing $`\stackrel{}{q}_+0`$ only Wigner-covariant relative variables are left on the Wigner hyperplane and there is no double counting of the center of mass \[only the decoupled canonical non-covariant “external” one $`\stackrel{}{z}_s`$, $`\stackrel{}{k}_s`$ is left\].
In Refs. there was a naive choice $`\stackrel{}{\eta }_+=\frac{1}{N}_{i=1}^N\stackrel{}{\eta }_i`$ of the Wigner spin 1 3-vector conjugate to $`\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}`$. Then, after realizing that $`\stackrel{}{\eta }_+0`$ does not imply $`\stackrel{}{\lambda }(\tau )=0`$, in Ref. a different choice $`\stackrel{}{q}_+`$ was made by utilizing the group-theoretical results of Ref.: now the time constancy of the gauge fixings $`\stackrel{}{q}_+0`$ implies $`\stackrel{}{\lambda }(\tau )=0`$. Moreover, the nonrelativistic limit of $`\stackrel{}{q}_+`$ is now the unique nonrelativistic center of mass.
In this gauge we get the simplest description of the dynamics on the Wigner hyperplanes: $`\stackrel{}{\sigma }=\stackrel{}{q}_+0`$ implies that the “internal” center of mass is put at the origin $`x_s^\mu (\tau )`$ of the coordinates \[$`z^\mu (\tau ,\stackrel{}{\sigma })=x_s^\mu (\tau )+ϵ_r^\mu (u(p_s))\sigma ^r`$ with $`x_s^\mu (\tau )=x_s^\mu (0)+u^\mu (p_s)\tau `$\] of the Wigner hyperplane. In this gauge the origin acquires the property $`\dot{x}_s^\mu (\tau )=u^\mu (p_s)`$ and becomes also the “external” Fokker-Pryce center of inertia of the isolated system \[in a future paper there will be a more detailed analysis of these problems\].
In order to find $`\stackrel{}{q}_+`$ one must take advantage of the “internal” realization of the Poincaré algebra inside the Wigner hyperplane. Ref. implies the following definition of the canonical “internal” 3-center of mass
$`\stackrel{}{q}_+`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{𝒦}_{(int)}}{\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}}}+`$ (233)
$`+`$ $`{\displaystyle \frac{\stackrel{}{𝒥}_{(int)}\times \stackrel{}{𝒫}_{(int)}}{\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}[𝒫_{(int)}^\tau +\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}]}}`$ (235)
$`+{\displaystyle \frac{\stackrel{}{𝒦}_{(int)}\stackrel{}{𝒫}_{(int)}\stackrel{}{𝒫}_{(int)}}{𝒫_{(int)}^\tau \sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}[𝒫_{(int)}^\tau +\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}]}}`$
$`{\displaystyle \genfrac{}{}{0pt}{}{}{\stackrel{}{𝒫}_{(int)}0}}`$ $`{\displaystyle \frac{\stackrel{}{𝒦}_{(int)}}{𝒫_{(int)}^\tau }}=\stackrel{}{R}_+.`$ (236)
Imposing the rest-frame condition, it is seen that $`\stackrel{}{q}_+`$ weakly coincides with the noncanonical “internal” Møller center of energy $`\stackrel{}{R}_+`$. In that same limit it is also equal to the “internal” Fokker-Pryce center of inertia defined by
$$\stackrel{}{Y}_+=\stackrel{}{q}_++\frac{\stackrel{}{S}_{(int)}\times \stackrel{}{𝒫}_{(int)}}{\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}[𝒫_{(int)}^\tau +\sqrt{(𝒫_{(int)}^\tau )^2(\stackrel{}{𝒫}_{(int)})^2}]},$$
(237)
where
$$\stackrel{}{S}_{(int)}\stackrel{}{𝒥}_{(int)}\stackrel{}{q}_+\times \stackrel{}{𝒫}_{(int)}\stackrel{}{S}_s.$$
(238)
With the gauge fixing condition $`\stackrel{}{q}_+0`$ and with $`T_S=\tau `$ one finds the following expression for the origin $`x_s^\mu (\tau )`$ of the coordinates on the Wigner hyperplane \[$`x_s^\mu (0)`$ is arbitrary\]
$$x_s^{(\stackrel{}{q}_+)\mu }(T_s)=x_s^\mu (0)+u^\mu (p_s)T_s.$$
(239)
It can be shown that this coincides with the covariant noncanonical “external” Fokker-Pryce center of inertia $`Y^\mu (\tau )`$. However, it is different from both the “external” center of mass $`\stackrel{~}{x}_s^\mu (\tau )`$ and the “external” center of energy of Møller $`R^\mu (\tau )`$.
Since $`\frac{d}{d\tau }\stackrel{}{q}_+\stackrel{}{=}\{\stackrel{}{q}_+,M\stackrel{}{\lambda }(\tau )\stackrel{}{}_p\}=\stackrel{}{\lambda }(\tau )0`$, there is no gauge freedom left and we could eliminate the variables $`\stackrel{}{q}_+`$, $`\stackrel{}{𝒫}_{(int)}=\stackrel{}{}_p`$ and look for a canonical basis of (Dirac observable) relative variables on the Wigner hyperplane.
Instead of doing that \[see Ref.\], in this paper we will go on to work with all the variables $`\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$, but we shall restrict their equations of motion to the gauge $`\stackrel{}{\lambda }(\tau )=0`$ without explicitly introducing $`\stackrel{}{q}_+0`$.
The equations of motion for the particles and for the electromagnetic field then become
$`{\displaystyle \frac{d}{d\tau }}(m_i{\displaystyle \frac{\dot{\stackrel{}{\eta }}_i(\tau )}{\sqrt{1\dot{\stackrel{}{\eta }}_i^2(\tau )}}})`$ $`\stackrel{}{=}`$ $`{\displaystyle \underset{ki}{}}{\displaystyle \frac{Q_iQ_k(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau ))}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_k(\tau )^3}}+`$ (240)
$`+`$ $`Q_i[\stackrel{ˇ}{\stackrel{}{E}}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))+\dot{\stackrel{}{\eta }}_i(\tau )\times \stackrel{ˇ}{\stackrel{}{B}}(\tau ,\stackrel{}{\eta }_i(\tau ))],`$ (241)
$`\mathrm{}\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\ddot{\stackrel{ˇ}{A}}_{}^r(\tau ,\stackrel{}{\sigma })+\mathrm{\Delta }\stackrel{ˇ}{A}_{}^r(\tau ,\stackrel{}{\sigma })\stackrel{}{=}J_{}^r(\tau ,\stackrel{}{\sigma })=`$ (242)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_iP_{}^{rs}(\stackrel{}{\sigma })\dot{\eta }^s(\tau )\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))=`$ (243)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\dot{\eta }^s(\tau )(\delta ^{rs}+{\displaystyle \frac{^r^s}{\mathrm{\Delta }}})\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))=`$ (244)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\dot{\eta }^s(\tau )[\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))+`$ (245)
$`+`$ $`{\displaystyle d^3\sigma ^{}\frac{\pi ^{rs}(\stackrel{}{\sigma }\stackrel{}{\sigma }^{})}{\stackrel{}{\sigma }\stackrel{}{\sigma }^{}^3}\delta ^3(\stackrel{}{\sigma }^{}\stackrel{}{\eta }_i(\tau ))},`$ (246)
with
$$\pi ^{rs}(\stackrel{}{\sigma }\stackrel{}{\sigma }^{})=\delta ^{rs}3(\sigma ^r\sigma _{}^{}{}_{}{}^{r})(\sigma ^s\sigma _{}^{}{}_{}{}^{s})/(\stackrel{}{\sigma }\stackrel{}{\sigma }^{})^2.$$
(247)
We point out that defining $`\stackrel{}{\beta }_i(\tau )=\dot{\stackrel{}{\eta }}_i(\tau )=\frac{d\stackrel{}{\eta }_i(\tau )}{d\tau }=\frac{1}{c}\frac{d\stackrel{}{\eta }_i^{^{}}(t)}{dt}`$ \[$`\tau =ct`$, $`\stackrel{}{\eta }_i^{^{}}(t)=\stackrel{}{\eta }_i(\tau )`$; even if we use everywhere $`c=1`$, we have momentarily reintroduced it\] and $`\stackrel{}{\beta }_i^{(h)}=d^h\stackrel{}{\beta }_i/d\tau ^h`$ and writing the particle equations of motion as (no sum over $`i`$)
$$\frac{d}{d\tau }(m_i\frac{\stackrel{}{\beta }_i(\tau )}{\sqrt{1\stackrel{}{\beta }_i^2(\tau )}})=\frac{m_i}{\sqrt{1\stackrel{}{\beta }_i^2(\tau )}}[\stackrel{}{\beta }_i^{(1)}+\stackrel{}{\beta }_i\frac{\stackrel{}{\beta }_i^{(1)}\stackrel{}{\beta }_i}{1\stackrel{}{\beta }_i^2(\tau )}]\stackrel{}{=}Q_i\stackrel{}{F}_i,$$
(248)
we obtain
$$m_i\frac{\stackrel{}{\beta }_i^{(1)}\stackrel{}{\beta }_i}{(1\stackrel{}{\beta }_i^2(\tau ))^{3/2}}\stackrel{}{=}Q_i\stackrel{}{\beta }_i\stackrel{}{F}_i,$$
(249)
so that
$$\stackrel{}{\beta }_i^{(1)}\stackrel{}{=}\frac{\sqrt{1\stackrel{}{\beta }_i^2(\tau )}}{m_i}Q_i(\stackrel{}{F}_i\stackrel{}{\beta }_i\stackrel{}{\beta }_i\stackrel{}{F}_i).$$
(250)
Thus in general we will have for every $`h1`$
$$\stackrel{}{\beta }_i^{(h)}\stackrel{}{=}Q_i\stackrel{}{G}_i,$$
(251)
so that using the Grassmann property of the charges
$$Q_i\stackrel{}{\beta }_i^{(h)}\stackrel{}{=}0,h1.$$
(252)
This will lead to important simplifications later allowing us to drop acceleration dependent terms in the force.
Due to the projector $`P_{}^{rs}(\stackrel{}{\sigma })`$ required by the rest-frame radiation gauge, the sources of the transverse (Wigner spin 1) vector potential becomes non - local and one has a system of integrodifferential equations (like with the equations generated by Fokker-Tetrode actions) with the open problem of how to define an initial value problem.
The Lagrangian equations identifying the rest frame become
$`{\displaystyle \underset{i=1}{\overset{N}{}}}(\eta _im_i{\displaystyle \frac{\dot{\stackrel{}{\eta }}_i(\tau )}{\sqrt{1\dot{\stackrel{}{\eta }}^2(\tau )}}}+Q_i\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\eta }_i(\tau )))+`$ (253)
$`+{\displaystyle d^3\sigma \underset{r}{}[(\stackrel{}{}\stackrel{ˇ}{A}_{}^r)\dot{\stackrel{ˇ}{A}}_{}^r](\tau ,\stackrel{}{\sigma })}\stackrel{}{=}\mathrm{\hspace{0.17em}0}.`$ (254)
## V Electromagnetic Lienard-Wiechert Potentials.
In this Section we will study the Lienard-Wiechert solutions of the previous radiation gauge field equations in the gauge $`\stackrel{}{\lambda }(\tau )=0`$ in absence of incoming radiation by using the results of Ref.. We shall study the $`retarded`$, $`advanced`$ and $`\frac{1}{2}(retarded+advanced)`$ Lienard-Wiechert potentials. Using the Smart-Winter expansion for retarded and advanced time dependence, we obtain an infinite series form of the retarded and advanced Lienard-Wiechert potentials depending on instantaneous accelerations of every order. It will be shown that the results of the previous Section imply that on the solutions of the particle equations of motion the higher accelerations decouple due to the semiclassical regularization $`Q_i^2=0`$. We show that this implies that the $`\frac{1}{2}(retardedadvanced)`$ Lienard-Wiechert potential vanishes at the semiclassical level and that there is only one semiclassical Lienard-Wiechert potential: $`retarded=advanced=\frac{1}{2}(retarded+advanced)`$. This allows us to re-express the semiclassical Lienard-Wiechert potential in terms of particle canonical coordinates and momenta \[the same can be done for the Lienard-Wiechert electric field, as it will be shown in the next Section\]. Therefore, we get a Hamiltonian description of the Lienard-Wiechert semiclassical solution: this sector of solutions can be identified as the symplectic submanifold of the space of solutions of the electromagnetic field equations determined by two pairs of second class constraints, which force the electromagnetic field to coincide with the semiclassical Lienard-Wiechert one.. After having gone to Dirac brackets with respect to them, we get a reduced phase space with only particles and we find a canonical basis $`\stackrel{~}{\stackrel{}{\eta }}_i`$, $`\stackrel{~}{\stackrel{}{\kappa }}_i`$ for these brackets for arbitrary $`N`$. In the new variables the rest-frame condition becomes $`_{i=1}^N\stackrel{~}{\stackrel{}{\kappa }}_i0`$, as expected in an instant form of dynamics.
### A Grassmann truncated form of the advanced and retarded Lienard Wiechert Solutions.
Here we develop the Grassmann truncated forms for the Lienard-Wiechert vector potential \[see the next Section for the transverse electric and magnetic fields\]. The $`\frac{1}{2}(retarded+advanced)`$ solutions are given (for $`\stackrel{}{\lambda }(\tau )=0`$) by \[for the sake of notational simplicity we will use the notation $`\stackrel{}{\kappa }_i(\tau )`$, $`\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{}{\pi }_{}(\tau ,\stackrel{}{\sigma })`$, instead of $`\stackrel{ˇ}{\stackrel{}{\kappa }}_i(\tau )`$, $`\stackrel{ˇ}{\stackrel{}{A}}_{}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{ˇ}{\stackrel{}{\pi }}_{}(\tau ,\stackrel{}{\sigma })`$, from now on\]
$`A_S^r(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{1}{2}}[A_+^r+A_{}^r](\tau ,\stackrel{}{\sigma })=`$ (255)
$`=`$ $`{\displaystyle \frac{1}{2}}𝒫_{}^{rs}(\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{2\pi }}{\displaystyle 𝑑\tau _1d^3\sigma _1[\theta (\tau \tau _1)+\theta (\tau \tau _1)]}`$ (257)
$`\delta [(\tau \tau _1)^2(\stackrel{}{\sigma }\stackrel{}{\sigma }_1)^2]\dot{\eta }_i^s(\tau _1)\delta ^3(\stackrel{}{\sigma }_1\stackrel{}{\eta }_i(\tau _1))=`$
$`=`$ $`𝒫_{}^{rs}(\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{2\pi c}}{\displaystyle 𝑑t_1\delta [(tt_1)^2\frac{1}{c^2}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(ct_1))^2]\beta _i^s(ct_1)}:=`$ (258)
$`:`$ $`={\displaystyle \underset{i=1}{\overset{N}{}}}Q_iA_{Si}^r(\tau ,\stackrel{}{\sigma }),`$ (259)
in which we have put $`\tau =ct`$, $`\stackrel{}{\beta }_i(\tau )=\dot{\stackrel{}{\eta }}_i(\tau )=\frac{1}{c}\frac{d\stackrel{}{\eta }_i^{^{}}(t)}{dt}`$ and $`\stackrel{}{A}_+=\stackrel{}{A}_{RET}`$ ($`\stackrel{}{A}_{}=\stackrel{}{A}_{ADV}`$) for the retarded (advanced) solution. The equation for $`t_1`$ is $`c^2(tt_1)^2=(\stackrel{}{\sigma }\stackrel{}{\eta }_i(ct_1))^2`$ with the two solutions being
$`t_{i+}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{1}{c}}\tau _{i+}(\tau ,\stackrel{}{\sigma })=t{\displaystyle \frac{1}{c}}r_{i+}(\tau _{i+}(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })={\displaystyle \frac{\tau }{c}}T_{i+}(\tau ,\stackrel{}{\sigma }),`$ (260)
$`t_i(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{1}{c}}\tau _i(\tau ,\stackrel{}{\sigma })=t+{\displaystyle \frac{1}{c}}r_i(\tau _i(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })={\displaystyle \frac{\tau }{c}}+T_i(\tau ,\stackrel{}{\sigma }),`$ (261)
for the retarded and for the advanced case respectively. The light cone delta function is
$`\delta [(\tau \tau _1)^2(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau _1))^2]={\displaystyle \frac{1}{c^2}}\delta [(tt_1)^2{\displaystyle \frac{1}{c^2}}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(ct_1))^2]=`$ (262)
$`={\displaystyle \frac{\delta [\tau _1\tau _{i+}(\tau ,\stackrel{}{\sigma })]}{2|\tau \tau _1\stackrel{}{\beta }_i(\tau _1)(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau _1))|}}+{\displaystyle \frac{\delta [\tau _1\tau _i(\tau ,\stackrel{}{\sigma })]}{2|\tau \tau _1\stackrel{}{\beta }_i(\tau _1)(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau _1))|}}.`$ (263)
The relative space location between the field point and the retarded or advanced particle position is
$$\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }))=\stackrel{}{r}_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })=r_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })\widehat{r}_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma }),$$
(264)
and its length is related to the time interval by
$`r_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })=|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }))|=cT_{i\pm }(\tau ,\stackrel{}{\sigma })=|\tau \tau _{i\pm }(\tau ,\stackrel{}{\sigma })|,`$ (265)
(266)
$`\tau \tau _{i\pm }(\tau ,\stackrel{}{\sigma })=\pm cT_{i\pm }(\tau ,\stackrel{}{\sigma })=\pm r_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma }).`$ (267)
The effective spatial interval is defined by
$$\rho _{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })=r_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })[1\stackrel{}{\beta }_i(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }))\widehat{r}_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })].$$
(268)
In terms of these variables, the retarded, advanced and time symmetric solutions are
$`A_\pm ^r(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}𝒫_{}^{rs}(\stackrel{}{\sigma }){\displaystyle \frac{\beta _i^s(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }))}{\rho _{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })}},`$ (269)
$`A_S^r(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_iA_{Si}^r(\tau ,\stackrel{}{\sigma })={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{8\pi }}𝒫_{}^{rs}(\stackrel{}{\sigma })\left[{\displaystyle \frac{\beta _i^s(\tau _{i+}(\tau ,\stackrel{}{\sigma }))}{\rho _{i+}(\tau _{i+}(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })}}+{\displaystyle \frac{\beta _i^s(\tau _i(\tau ,\stackrel{}{\sigma }))}{\rho _i(\tau _i(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })}}\right].`$ (270)
We use the Smart-Wintner expansion
$`f(\tau _{i\pm })`$ $`=`$ $`f(\tau cT_{i\pm }(\tau _{i\pm }((\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })))=f(\tau [\pm r_{i\pm }(\tau _{i\pm }(\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })])=`$ (271)
$`=`$ $`f(\tau )+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^k}{k!}}{\displaystyle \frac{d^{k1}}{d\tau ^{k1}}}\left[\left(\pm r_i(\tau ,\stackrel{}{\sigma })\right)^k{\displaystyle \frac{df(\tau )}{d\tau }}\right]=`$ (272)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^k}{k!}}{\displaystyle \frac{d^k}{d\tau ^k}}\left[\left(\pm r_i(\tau ,\stackrel{}{\sigma })\right)^k[1\stackrel{}{\beta }_i(\tau )\widehat{r}_i(\tau ,\stackrel{}{\sigma })]f(\tau )\right],`$ (273)
where
$`\stackrel{}{r}_i(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`r_i(\tau ,\stackrel{}{\sigma })\widehat{r}_i(\tau ,\stackrel{}{\sigma })=\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )=\stackrel{}{r}_{i\pm }(\tau _\pm (\tau ,\stackrel{}{\sigma }),\stackrel{}{\sigma })|_{\tau _{i\pm }(\tau ,\stackrel{}{\sigma })=\tau },`$ (274)
$`f(\tau _{i\pm })`$ $`=`$ $`{\displaystyle \frac{\beta _i^s(\tau _{i\pm })}{\rho _{i\pm }(\tau _{i\pm })}}={\displaystyle \frac{\beta _i^s(\tau _{i\pm })}{r_{i\pm }(\tau _{i\pm })[1\stackrel{}{\beta }_i(\tau _{i\pm })\widehat{r}_{i\pm }(\tau ,\stackrel{}{\sigma })]}}.`$ (275)
and where the last line in Eq.(273) is identical to the previous one since $`\frac{dr_i(\tau ,\stackrel{}{\sigma })}{d\tau }=\stackrel{}{\beta }_i(\tau )\widehat{r}_i(\tau ,\stackrel{}{\sigma })`$.
Hence we get
$$A_\pm ^r(\tau ,\stackrel{}{\sigma })=\underset{i=1}{\overset{N}{}}\frac{Q_i}{4\pi }𝒫_{}^{rs}(\stackrel{}{\sigma })\underset{k=0}{\overset{\mathrm{}}{}}\frac{()^k}{k!}\frac{d^k}{d\tau ^k}[r_i^{k1}(\tau ,\stackrel{}{\sigma })\beta _i^s(\tau )],A_S^r=\frac{1}{2}(A_+^r+A_{}^r).$$
(276)
In order to evaluate the above derivatives we need the Leibnitz formula for the $`k`$th derivative of the product $`f(\tau )g(\tau )`$
$$\frac{d^k}{d\tau ^k}(fg)=\underset{m=0}{\overset{k}{}}\frac{k!}{m!(km)!}\frac{d^mf}{d\tau ^m}\frac{d^{km}g}{d\tau ^{km}},$$
(277)
Thus we get
$`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^k}{k!}}{\displaystyle \frac{d^k}{d\tau ^k}}[r_i^{k1}\beta _i^s]`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^k}{k!}}{\displaystyle \underset{m=0}{\overset{k}{}}}{\displaystyle \frac{k!}{m!(km)!}}{\displaystyle \frac{d^mr^{k1}}{d\tau ^m}}{\displaystyle \frac{d^{km}\beta ^s}{d\tau ^{km}}}=`$ (278)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^k}{m!(km)!}}{\displaystyle \frac{d^mr^{k1}}{d\tau ^m}}{\displaystyle \frac{d^{km}\beta ^s}{d\tau ^{km}}}`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{h=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^{h+m}}{m!h!}}{\displaystyle \frac{d^h\beta ^s}{d\tau ^h}}{\displaystyle \frac{d^mr^{h+m1}}{d\tau ^m}}.`$ (279)
Using the notation $`\beta ^{(h)s}=\frac{d^h\beta ^s}{d\tau ^h}`$ we obtain the following expression for the vector potential
$$A_\pm ^r(\tau ,\stackrel{}{\sigma })=\underset{i=1}{\overset{N}{}}\frac{Q_i}{4\pi }𝒫_{}^{rs}(\stackrel{}{\sigma })\underset{h=0}{\overset{\mathrm{}}{}}\frac{()^h}{h!}\beta _i^{(h)s}(\tau )\varphi _{i\pm ,h}(\tau ,\stackrel{}{\sigma }),$$
(280)
in which
$$\varphi _{i\pm ,h}(\tau ,\stackrel{}{\sigma })=\underset{m=0}{\overset{\mathrm{}}{}}\frac{()^m}{m!}\frac{d^mr_i^{h+m1}(\tau ,\stackrel{}{\sigma })}{d\tau ^m}=\underset{m=0}{\overset{\mathrm{}}{}}\frac{()^m}{m!}\frac{d^m}{d\tau ^m}[\sqrt{(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))^2}]^{m+h1}.$$
(281)
In order to display the result of the evaluation of the derivative we use the formula
$$\frac{d^m}{d\tau ^m}R(f(\tau ))=\underset{n=0}{\overset{m}{}}\underset{n_1n_{2..}}{}\frac{m!}{n_1!n_2!..}\frac{d^nR(f(\tau ))}{df^n}|_{f=f(\tau )}(\frac{1}{1!}\frac{df(\tau )}{d\tau })^{n_1}(\frac{1}{2!}\frac{d^2f(\tau )}{d\tau ^2})^{n_2}\mathrm{}$$
(282)
(with the summations restricted so that $`_rn_r=n,_rrn_r=m`$ ) to obtain
$`\varphi _{i\pm ,h}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^m}{m!}}{\displaystyle \underset{n=0}{\overset{m}{}}}{\displaystyle \underset{n_1n_{2..}}{}}{\displaystyle \frac{m!}{n_1!n_2!..}}`$ (284)
$`{\displaystyle \frac{^nr_i^{m+h1}(\tau ,\stackrel{}{\sigma })}{\stackrel{}{r}_i^n}}\left({\displaystyle \frac{\stackrel{}{\beta }_i(\tau )}{1!}}\right)^{n_1}\left({\displaystyle \frac{\stackrel{}{\beta }_i^{(1)}(\tau )}{2!}}\right)^{n_2}\mathrm{}`$
In this expression the symbol $``$ represents a scalar product between the tensors to the left and to the right with the summation$`_rn_r=n`$ indicating how the indices would be matched. Changing the $`m`$ summation index to $`k=mn`$ we obtain
$`\varphi _{i\pm ,h}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n_1n_{2..}}{}}{\displaystyle \frac{()^{k+n}()^{\mathrm{\Sigma }_rn_r=n}}{n_1!n_2!..}}`$ (286)
$`{\displaystyle \frac{^nr_i^{k+n+h1}(\tau ,\stackrel{}{\sigma })}{\stackrel{}{r}_i^n}}\left({\displaystyle \frac{\stackrel{}{\beta }_i(\tau )}{1!}}\right)^{n_1}\left({\displaystyle \frac{\stackrel{}{\beta }_i^{(1)}(\tau )}{2!}}\right)^{n_2}\mathrm{}`$
(In this latter summation $`_rn_r=n,_rrn_r=n+k.`$)
Now we can take advantage of the Grassmann charges to significantly simplify the above multi-summations. As we have seen above, with a semiclassical $`Q_i`$ there are no accelerations on shell ($`Q_i\stackrel{}{\beta }_i^{(h)}\stackrel{}{=}0`$) in the equations of motion of the particle ‘$`i`$’, since both the Coulomb potential and the Lienard-Wiechert Lorentz force on particle ‘$`i`$’ produced by the other particles, i.e. $`Q_i[\stackrel{}{E}_{}(\tau ,\stackrel{}{\eta }_i(\tau ))+\stackrel{}{\beta }_i(\tau )\times \stackrel{}{B}(\tau ,\stackrel{}{\eta }_i(\tau ))]`$, are proportional to $`Q_i`$. Therefore, the full set of Hamilton equations (212), (215) for both fields and particles imply that at the semiclassical level we have a natural “order reduction” of the final particle equation of motion in the Lienard-Wiechert sector \[only second order differential equations\].
One effect of this truncation is the elimination of multi-particle forces; all the interactions will be pairwise, in both the Lagrangian and Hamiltonian formalisms. This was to be expected since the rest-frame instant form is an equal-time description of the $`N`$ particle system: (acceleration-independent) 3-body,.. N-body forces appear as soon as we go to a description with no concept of equal time, like in the standard approach with $`N`$ first class constraints .
Thus the only contributing indices are $`n_2=n_3=..=0,n_1=n`$ and our expression for the transverse vector potentials simplify to
$`A_\pm ^r(\tau ,\stackrel{}{\sigma })`$ $`\stackrel{}{=}`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}𝒫_{}^{rs}(\stackrel{}{\sigma })\beta _i^s(\tau )\varphi _{i\pm ,0}(\tau ,\stackrel{}{\sigma })\stackrel{}{=}`$ (287)
$`\stackrel{}{=}`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}𝒫_{}^{rs}(\stackrel{}{\sigma })\beta _i^s(\tau ){\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\pm )^n}{n!}}{\displaystyle \frac{^nr_i^{n1}(\tau ,\stackrel{}{\sigma })}{\stackrel{}{r}_i^n}}\left({\displaystyle \frac{\stackrel{}{\beta }_i(\tau )}{1!}}\right)^n,`$ (288)
$`A_S^r(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{1}{2}}(A_+^r+A_{}^r)(\tau ,\stackrel{}{\sigma }).`$ (289)
Since $`r_i=\sqrt{\stackrel{}{r}_i^2}`$, we see that for odd n=2m+1 we get
$$\frac{^{2m+1}}{\stackrel{}{r}_i^{2m+1}}(\sqrt{\stackrel{}{r}_i^2})^{2m}=\frac{^{2m+1}}{\stackrel{}{r}_i^{2m+1}}(\stackrel{}{r}_i^2)^m=0,$$
(290)
and this implies the equality of the retarded, advanced and symmetric Lienard-Wiechert potentials on-shell
$`A_S^r(\tau ,\stackrel{}{\sigma })`$ $`\stackrel{}{=}`$ $`A_\pm ^r(\tau ,\stackrel{}{\sigma })\stackrel{}{=}{\displaystyle \underset{i=1,iu}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}𝒫_{}^{rs}(\stackrel{}{\sigma })\beta _i^s(\tau )`$ (292)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}\left(\stackrel{}{\beta }_i(\tau ){\displaystyle \frac{^{2m}}{\stackrel{}{r}_i^{2m}}}\right)r_i^{2m1}(\tau ,\stackrel{}{\sigma }).`$
Therefore, at the semiclassical level there is only one Lienard-Wiechert sector with a uniquely determined standard action-at-a-distance interaction.
We use a tensor notation to write the transverse symmetric vector potential above as
$$\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })\stackrel{}{=}\underset{i=1}{\overset{N}{}}\frac{Q_i}{4\pi }𝒫_{}\dot{\stackrel{}{\eta }}_i\underset{m=0}{\overset{\mathrm{}}{}}\frac{\dot{\stackrel{}{\eta }}_{ij_1}(\tau )..\dot{\stackrel{}{\eta }}_{ij_{2m}}(\tau )}{(2m)!}\frac{^{2m}|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}}{\sigma _{j_1}..\sigma _{j_{2m}}}.$$
(293)
Using the definition of the Coulomb projection operator
$$𝒫(\stackrel{}{\sigma })_{hk}F(\stackrel{}{\sigma })=\delta _{hk}F(\stackrel{}{\sigma })\frac{1}{4\pi }d^3\sigma ^{}\frac{^2}{\sigma _h\sigma _k}\frac{1}{|\stackrel{}{\sigma }^{}\stackrel{}{\sigma }|}F(\stackrel{}{\sigma }^{}),$$
(294)
and compactifying the notation still further we obtain \[$`\stackrel{}{}_\sigma =/\stackrel{}{\sigma }`$\]
$`\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })\stackrel{}{=}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}[\dot{\stackrel{}{\eta }}_i(\tau )(\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}`$ (295)
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle }d^3\sigma ^{}[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma ){\displaystyle \frac{1}{|\stackrel{}{\sigma }^{}\stackrel{}{\sigma }|}}](\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma ^{})^{2m}|\stackrel{}{\sigma }^{}\stackrel{}{\eta }_i(\tau )|^{2m1}].`$ (296)
Integration by parts and changing from $`\frac{}{\stackrel{}{\sigma }^{}}`$ to $`\frac{}{\stackrel{}{\sigma }}`$ and translation gives
$`\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })`$ $`\stackrel{}{=}`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}[\dot{\stackrel{}{\eta }}_i(\tau )(\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}`$ (298)
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle }d^3\sigma ^{}\left(\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m+1}{\displaystyle \frac{1}{|\stackrel{}{\sigma }^{}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))|}}\right)\sigma _{}^{}{}_{}{}^{2m1}].`$
The integral above is finite, and thus we can view it as the $`\mathrm{\Lambda }\mathrm{}`$ limit of an integral with a cutoff $`\mathrm{\Lambda }`$ and take the derivatives out. The integral is thus of the form
$$\frac{1}{4\pi }\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m+1}d^3\sigma ^{}\frac{\sigma _{}^{}{}_{}{}^{2m1}}{|\stackrel{}{\sigma }^{}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ))|},$$
(299)
and
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle _\mathrm{\Lambda }}d^3\sigma ^{}{\displaystyle \frac{\sigma _{}^{}{}_{}{}^{2m1}}{|\stackrel{}{\sigma }^{}(\stackrel{}{\sigma }\stackrel{}{\eta }_i)|}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^\mathrm{\Lambda }}𝑑\sigma ^{}\sigma _{}^{}{}_{}{}^{2m+1}{\displaystyle _1^1}{\displaystyle \frac{dz}{\sqrt{\stackrel{}{\sigma ^{}}^2+(\stackrel{}{\sigma }\stackrel{}{\eta }_i)^22\sigma ^{}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|z}}}`$ (300)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^\mathrm{\Lambda }}𝑑\sigma ^{}\sigma _{}^{}{}_{}{}^{2m+1}{\displaystyle \frac{1}{\sigma ^{}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}\sqrt{\stackrel{}{\sigma ^{}}^2+(\stackrel{}{\sigma }\stackrel{}{\eta }_i)^22\sigma |\stackrel{}{\sigma }\stackrel{}{\eta }_i|}`$ (301)
$`=`$ $`{\displaystyle \frac{1}{2|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}{\displaystyle _0^\mathrm{\Lambda }}𝑑\sigma ^{}\sigma _{}^{}{}_{}{}^{2m}(|\stackrel{}{\sigma }^{}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|||\stackrel{}{\sigma }^{}+|\stackrel{}{\sigma }\stackrel{}{\eta }_i||)`$ (302)
$`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }^{2m+1}}{2m+1}}{\displaystyle \frac{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m+1}}{(2m+1)(2m+2)}}.`$ (303)
Note that the $`\mathrm{\Lambda }`$ cutoff will get killed by the $`\sigma `$ derivatives. Thus, we obtain
$`\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })`$ $`\stackrel{}{=}`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{(2m)!}}\dot{\stackrel{}{\eta }}_i(\tau )(\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m}|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}`$ (305)
$`{\displaystyle \frac{1}{(2m+2)!}}\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m+1}]:=`$
$`:`$ $`={\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\dot{\stackrel{}{\eta }}_i(\tau )).`$ (306)
Using the first half of particle Hamilton equations (212) \[with $`\stackrel{}{\lambda }(\tau )=0`$\] in the form $`\dot{\stackrel{}{\eta }_i}=\stackrel{}{\kappa }_i/\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}+O(Q_i)`$, we can, as shown in Appendix A, arrive at the following closed form of the vector potential \[$`\stackrel{}{\eta }_i=\stackrel{}{\eta }_i(\tau )`$, $`\stackrel{}{\kappa }_i=\stackrel{}{\kappa }_i(\tau )`$\]
$`\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })\stackrel{}{=}{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau )),`$ (307)
(308)
$`\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)={\displaystyle \frac{1}{4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}}}`$ (309)
$`\stackrel{}{\kappa }_i(𝐈{\displaystyle \frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^2}})({\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}{\sqrt{m_i^2+(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}}}1)\times `$ (310)
$`{\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}{\stackrel{}{\kappa }_i^2(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}}].`$ (311)
### B Lienard-Wiechert Second-Class Constraints, their Dirac Brackets and the New Canonical Variables.
Thus far we have the reduced phase space of $`N`$ charged particles plus the transverse electromagnetic field. This is a well defined isolated system with a global Darboux basis \[ $`\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i,\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma }),\stackrel{}{\pi }_{}(\tau ,\stackrel{}{\sigma })`$\] and a well defined physical Hamiltonian, the invariant mass $`M=𝒫_{(int)}^\tau `$. All possible configurations of motion take place in this reduced phase space. The space of solutions of Hamilton’s equations is a symplectic space in that there is a definition of Poisson brackets on the space of solutions. The question arises whether one can select a subset of solutions of the equations of motion which is still a symplectic manifold: an arbitrarily chosen set of solutions will not form a symplectic manifold. The method we propose here is to add by hand a set of second class constraints “compatible with the equations of motion” which amounts to the selection of a symplectic submanifold of the symplectic manifold of solutions.
The above Grassmann truncated semiclassical Lienard Wiechert solution $`\stackrel{}{A}_S`$ for the vector potential with $`\stackrel{}{\pi }_S=\stackrel{}{E}_S=\frac{}{\tau }\stackrel{}{A}_S`$ for the canonical conjugate field momentum \[see Eq.(359) in the next Section\] and provide us such a set of second class constraints by way of
$`\stackrel{}{\chi }_1(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))0,`$ (312)
$`\stackrel{}{\chi }_2(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\stackrel{}{\pi }_{}(\tau ,\stackrel{}{\sigma }){\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\stackrel{}{\pi }_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))0.`$ (313)
These constraints allow us to eliminate the canonical degrees of freedom of the radiation field and to get the symmetric Lienard-Wiechert reduced phase space, in which there are only particle degrees of freedom. This has an immediate and important consequence: the independent variables $`\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i`$ will no longer be canonical when one imposes these constraints by way of modified Dirac brackets.
Now in order to compute the effects of these constraints we must use them in the construction of Dirac brackets. This requires that we compute the 6x6 matrix of brackets
$$\left(\genfrac{}{}{0pt}{}{\{\stackrel{}{\chi }_1,\stackrel{}{\chi }_1\}}{\{\stackrel{}{\chi }_2,\stackrel{}{\chi }_1\}}\genfrac{}{}{0pt}{}{\{\stackrel{}{\chi }_1,\stackrel{}{\chi }_2\}}{\{\stackrel{}{\chi }_2,\stackrel{}{\chi }_2\}}\right).$$
(314)
It turns out that this matrix bracket is relatively simple, due to the Grassmann charges. Consider, for example the case of two particles. The particle or Lienard-Wiechert parts of the matrix bracket vanish since $`Q_1^2=0=Q_2^2`$ and cross terms vanish because they involve Poisson brackets of particle one variables with particle two variables. Thus the only part of the 6x6 matrix bracket that contributes is from the field variables. It has the form
$$\{\stackrel{}{\chi }_1(\tau ,\stackrel{}{\sigma }_1),\stackrel{}{\chi }_2(\tau ,\stackrel{}{\sigma }_2)\}=(𝐈\frac{\stackrel{}{}\stackrel{}{}}{\stackrel{}{}^2})\delta ^3(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2),$$
(315)
and since
$$\{\stackrel{}{\chi }_1,\stackrel{}{\chi }_1\}=0=\{\stackrel{}{\chi }_2,\stackrel{}{\chi }_2\},$$
(316)
only the 3x3 off diagonal portion contributes.
In order to have a well defined Dirac bracket we need to use a modified form of the Dirac bracket in which the inverse of the matrix of constraint Poisson brackets is used. Calling this matrix $`C`$, we define $`\stackrel{~}{C}^1`$ so that $`C\stackrel{~}{C}^1=(𝐈\frac{\stackrel{}{}\stackrel{}{}}{\stackrel{}{}^2})\delta ^3(\stackrel{}{\sigma }_1\stackrel{}{\sigma }_2)`$. But the transverse form of the delta function allows us to use the idempotent property of the projector to show that the inverse of $`C`$ in this sense is just $`C`$ itself. In that case for two functions $`f(\stackrel{}{\kappa }_i,\stackrel{}{\eta }_i),g(\stackrel{}{\kappa }_i,\stackrel{}{\eta }_i)`$ of the particle variables the Dirac bracket becomes
$`\{f,g\}^{}`$ $`=`$ $`\{f,g\}`$ (319)
$`[{\displaystyle }d^3\sigma \{f,{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))\}\{{\displaystyle \underset{j}{}}Q_j\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j(\tau ),\stackrel{}{\kappa }_j(\tau )),g\}`$
$`\{f,{\displaystyle \underset{j}{}}Q_j\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j(\tau ),\stackrel{}{\kappa }_j(\tau ))\}\{{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau )),g\}].`$
This bracket will lead to a new symplectic manifold by altering the basic commutation relations and providing us with new canonical variables. Toward this end we define the following scalar function.
$$𝒦=\underset{i=1}{\overset{N1}{}}\underset{j=i+1}{\overset{N}{}}Q_iQ_j𝒦_{ij}(\stackrel{}{\kappa }_i,\stackrel{}{\kappa }_j;\stackrel{}{\eta }_i\stackrel{}{\eta }_j),$$
(320)
in which
$`𝒦_{ij}`$ $`=`$ $`{\displaystyle }d^3\stackrel{}{\sigma }[\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j)`$ (321)
$``$ $`\stackrel{}{A}_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j)\stackrel{}{\pi }_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)]=`$ (322)
$`=`$ $`𝒦_{ij}(\stackrel{}{\kappa }_i,\stackrel{}{\kappa }_j;\stackrel{}{\eta }_i\stackrel{}{\eta }_j)=𝒦_{ji}.`$ (323)
Let $`\stackrel{~}{\stackrel{}{\eta }_i}=\stackrel{}{\eta }_i+\stackrel{}{\alpha }_i,`$ $`\stackrel{~}{\stackrel{}{\kappa }_i}=\stackrel{}{\kappa }_i+\stackrel{}{\beta }_i,`$ $`i=1,2,..N`$ where
$`\stackrel{}{\alpha }_i`$ $`=`$ $`a_i{\displaystyle \underset{j=i+1}{\overset{N}{}}}Q_iQ_j_{\kappa _i}𝒦_{ij}+\overline{a}_i{\displaystyle \underset{j=1}{\overset{i1}{}}}Q_iQ_j_{\kappa _i}𝒦_{ji},`$ (324)
$`\stackrel{}{\beta }_i`$ $`=`$ $`b_i{\displaystyle \underset{j=i+1}{\overset{N}{}}}Q_iQ_j_{\eta _i}𝒦_{ij}+\overline{b}_i{\displaystyle \underset{j=1}{\overset{i1}{}}}Q_iQ_j_{\eta _i}𝒦_{ji}.`$ (325)
Since they do not appear in these equations, we may choose $`\overline{a}_1=\overline{b}_1=a_N=b_N=0`$.
We determine relations between the unknown coefficients by requiring that $`\stackrel{~}{\stackrel{}{\eta }_i},\stackrel{~}{\stackrel{}{\kappa }_j}`$ be independent canonical variables. So, for example, (for $`k<l`$)
$`\{\stackrel{~}{\stackrel{}{\eta }_k},\stackrel{~}{\stackrel{}{\eta }_l}\}^{}`$ $`=`$ $`\{\stackrel{~}{\stackrel{}{\eta }_k},\stackrel{~}{\stackrel{}{\eta }_l}\}`$ (326)
$``$ $`[{\displaystyle }d^3\sigma \{\stackrel{~}{\stackrel{}{\eta }_k},{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)\}\{{\displaystyle \underset{j}{}}Q_j\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j),\stackrel{~}{\stackrel{}{\eta }_l}\}`$ (328)
$`\{\stackrel{~}{\stackrel{}{\eta }_k},{\displaystyle \underset{j}{}}Q_j\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j)\}\{{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i),\stackrel{~}{\stackrel{}{\eta }_l}\}]=`$
$`=`$ $`\{\stackrel{}{\eta }_k,\stackrel{}{\alpha }_l\}+\{\stackrel{}{\alpha }_k,\stackrel{}{\eta }_l\}+Q_kQ_l_{\kappa _k}_{\kappa _l}𝒦_{kl}=0.`$ (329)
Then using the expressions for $`\stackrel{}{\alpha }_i`$ leads to
$$a_l\overline{a}_k=1,k>l;\overline{a}_la_k=1,l>k.$$
(330)
Solving this gives
$`\overline{a}_2`$ $`=`$ $`\overline{a}_3=..=\overline{a}_N:=\overline{a},`$ (331)
$`a_1`$ $`=`$ $`a_2=..a_{N1}:=a.`$ (332)
Similarly, requiring that $`\stackrel{~}{\{\stackrel{}{\kappa }_k},\stackrel{~}{\stackrel{}{\kappa }_l}\}^{}=0`$ leads to
$`\overline{b}_2`$ $`=`$ $`\overline{b}_3=..=b_N:=\overline{b},`$ (333)
$`b_1`$ $`=`$ $`b_2=..=b_{N1}:=b.`$ (334)
Requiring that $`\{\stackrel{~}{\stackrel{}{\eta }_i},`$ $`\stackrel{~}{\stackrel{}{\kappa }_i}\}^{}=\stackrel{}{\stackrel{}{1}}`$ leads to
$$a_i+b_i=0,i=1,..,N1;\overline{a}_i+\overline{b}_i=0,i=2,..,N.$$
(335)
This condition implies that the first two conditions are equivalent to one another. The requirement that $`\{\stackrel{~}{\stackrel{}{\eta }_k},`$ $`\stackrel{~}{\stackrel{}{\kappa }_l}\}^{}=0`$ for $`k<l`$ leads to
$$a_k+\overline{b}_l=1;k<l,$$
(336)
or
$$a_k\overline{a}_l=1;k<l,$$
(337)
which is also the same as the first condition. While for $`l<k`$ it leads to
$$\overline{a}_k+b_l=1;l<k,$$
(338)
or
$$\overline{a}_ka_l=1;l<k,$$
(339)
which again is the same as the first condition. This leaves us with just two unknowns $`a`$ and $`\overline{a}`$.
So in summary we have
$`\stackrel{~}{\stackrel{}{\eta }_i}`$ $`=`$ $`\stackrel{}{\eta }_i+a{\displaystyle \underset{j=i+1}{\overset{N}{}}}Q_iQ_j\stackrel{}{}_{\kappa _i}𝒦_{ij}+\overline{a}{\displaystyle \underset{j=1}{\overset{i1}{}}}Q_iQ_j\stackrel{}{}_{\kappa _i}𝒦_{ji},`$ (340)
$`\stackrel{~}{\stackrel{}{\kappa }_i}`$ $`=`$ $`\stackrel{}{\kappa }_ia{\displaystyle \underset{j=i+1}{\overset{N}{}}}Q_iQ_j\stackrel{}{}_{\eta _i}𝒦_{ij}\overline{a}{\displaystyle \underset{j=1}{\overset{i1}{}}}Q_iQ_j\stackrel{}{}_{\eta _i}𝒦_{ji}.`$ (341)
Let us rewrite the rest frame condition Eq.(178)
$`\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}={\displaystyle \underset{i1}{\overset{N}{}}}\stackrel{}{\kappa }_i+{\displaystyle d^3\sigma [\stackrel{}{\pi }_{}\times \stackrel{}{B}](\tau ,\stackrel{}{\sigma })}=`$ (342)
$`=`$ $`{\displaystyle \underset{i1}{\overset{N}{}}}\stackrel{}{\kappa }_i+{\displaystyle \underset{i<j}{}}Q_iQ_j{\displaystyle }d^3\sigma [\stackrel{}{\pi }_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)\times (\stackrel{}{}_\sigma \times \stackrel{}{A}_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j))+`$ (344)
$`+\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j)\times (\stackrel{}{}_\sigma \times \stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i))]0.`$
in these new canonical variables. Let us expand the cross products, integrate by parts and use the transverse gauge condition to get
$$\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}=\underset{i1}{\overset{N}{}}\stackrel{}{\kappa }_i+\underset{i<j}{}Q_iQ_j\stackrel{}{}_{\eta _j}𝒦_{ij}=0.$$
(345)
If we choose $`a=\overline{a}=\frac{1}{2}`$ this becomes
$$\stackrel{}{}_p=\stackrel{}{𝒫}_{(int)}=\stackrel{~}{\stackrel{}{\kappa }}_+=\underset{i=1}{\overset{N}{}}\stackrel{~}{\stackrel{}{\kappa }_i}=0,$$
(346)
like in the case of $`N`$ either free or interacting particles on the Wigner hyperplane \[in an instant form of dynamics the Poincaré generators $`\stackrel{}{𝒫}_{(int)}`$ do not depend on the interaction\].
In other words the rest frame condition is simply that the sum of the $`N`$ (new) canonical momentum is zero. Note that this same choice gives
$$\underset{i=1}{\overset{N}{}}\stackrel{~}{\stackrel{}{\eta }_i}=\underset{i1}{\overset{N}{}}\stackrel{}{\eta }_i\underset{i<j}{}Q_iQ_j\stackrel{}{}_{\kappa _j}𝒦_{ij}.$$
(347)
Therefore, with the choice $`a=\overline{a}=\frac{1}{2},`$ Eq.(341) defining the final canonical variables becomes
$`\stackrel{~}{\stackrel{}{\eta }_i}`$ $`=`$ $`\stackrel{}{\eta }_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{}}Q_iQ_j\stackrel{}{}_{\kappa _j}𝒦_{ij},`$ (348)
$`\stackrel{~}{\stackrel{}{\kappa }_i}`$ $`=`$ $`\stackrel{}{\kappa }_i{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{}}Q_iQ_j\stackrel{}{}_{\eta _i}𝒦_{ij},`$ (349)
with $`𝒦_{ij}`$ given by Eq.(323).
In the next section we will re-express the other internal Poincaré generators $`M=𝒫_{(int)}^\tau `$, $`\stackrel{}{𝒥}_{(int)}`$, $`\stackrel{}{𝒦}_{(int)}`$ of Eqs.(205) and the internal center-of-mass coordinate $`\stackrel{}{q}_+\stackrel{}{𝒦}_{(int)}/M`$ of Eq.(236) in these final canonical variables.
## VI The Exact Darwin Hamiltonian from the Invariant Mass.
In this Section our aim is to use the explicit semiclassical Lienard Wiechert solution of Eqs.(306), (311) for the transverse vector potential to obtain an explicit form of the instantaneous action-at-a-distance potentials present in the invariant mass $`M=𝒫_{(int)}^\tau `$ of Eq.(205) after the elimination of the electromagnetic degrees of freedom. In its phase space form this Hamiltonian for the $`\tau T_s`$-evolution will contain:
i) a vector potential $`\stackrel{}{V}_i(\tau )=Q_i_{ij}^{1..N}Q_j\stackrel{}{A}_{Sj}(\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau ),\stackrel{}{\kappa }_j(\tau ))`$, minimally coupled to $`\stackrel{}{\kappa }_i(\tau )`$, under the square root kinetic energy term of each particle ‘i’;
ii) a scalar potential $`U(\tau )=\frac{1}{2}d^3\sigma [\stackrel{}{\pi }_S^2+\stackrel{}{B}_S^2](\tau ,\stackrel{}{\sigma })`$, coming from the field energy, which adds to the Coulomb potentials.
Due to $`Q_i^2=0`$ we can extract the vector potentials from the square roots: the semiclassical contribution from the vector potentials is a new effective scalar potential $`U_1=_{i=1}^NQ_i\frac{\stackrel{}{\kappa }_i\stackrel{}{V}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}}`$ and we get the complete Darwin potential $`V_{DAR}=U+U_1`$ added to the Coulomb one. In this form the invariant mass becomes the exact semiclassical Darwin Hamiltonian and $`V_{DAR}`$ is the Darwin potential to all orders of $`1/c^2`$ for every $`N`$. If we call $`V_{LOD}`$ the lowest $`1/c^2`$ order historical Darwin potential \[see Eq.(395)\], we have $`U_1=2V_{LOD}+U_{1HOD}`$ \[see Eq.(392)\], $`U=V_{LOD}U_{1HOD}+U_{HOD}`$ \[see Eq.(416)\], $`V_{DAR}=V_{LOD}+U_{HOD}`$. \[“$`LOD`$” and “$`HOD`$” mean lowest and higher order in $`1/c^2`$ respectively\]. When we re-express the invariant mass in terms of the final canonical variables, there is an extra contribution $`U_{HOD}^{}`$ coming from the square roots, so that at the end the final Darwin potential is
$$\stackrel{~}{V}_{DAR}=V_{DAR}+U_{HOD}^{}=V_{LOD}+V_{HOD};V_{HOD}=U_{HOD}+U_{HOD}^{}.$$
(350)
In this Section we shall evaluate the complete (to all orders in $`1/c^2`$) Darwin potential in the old (no longer canonical) variables and then we shall re-express it in terms of the new final canonical variables. We begin by obtaining the contribution of the field energy integrals, expressed in terms of the canonical particle variables. Using the truncation properties of the Grassmann charges we extract from the kinetic piece the vector potential portion and combine it with the field energy integral. In addition to the naive kinetic part (expressed in terms of the old canonical momentum) we obtain the rest frame Coulomb part, a part that generalizes the standard Darwin interaction and a double infinite series containing all higher order corrections. As an extra check on these generalized Darwin interactions in $`M`$ we obtain an independent derivation of this series using the Lagrangian expression for the invariant mass. Then we express the kinetic portion in terms of the new canonical momentum to obtain the final form of the complete Darwin Hamiltonian. This interaction Hamiltonian contains no $`N`$-body forces and is a sum of two-body portions. In the center-of-mass rest frame the double infinite series can be summed exactly to obtain a closed form expression for the special case of two particles. Then all the generators of the “internal” Poincaire‘ algebra and the energy-momentum tensor are expressed in the new variables in the $`N`$-particle case. As with the three-momentum we find that the internal angular momentum does not depend on the interaction. Also we obtain the “internal” center of mass $`\stackrel{}{q}_+`$ and there are some comments on how to find a collective variable (replacing the center of mass) for a cluster of particles interacting with the remaining ones of the isolated system.
### A Field Energy and Momentum Integrals.
Although we have summed exactly to a closed form (311) the semiclassical Lienard-Wiechert solution, we use the series form (306) for finding the expression for the invariant mass $`M`$, since the closed form provides no simplification in obtaining that expression. From the above we can find expressions for the semiclassical Lienard-Wiechert electric and magnetic fields. For the electric field we find \[$`\ddot{\stackrel{}{\eta }}_i(\tau )`$ does not contribute due to Eq.(252), and the same is true of $`\dot{\stackrel{}{\kappa }}_i(\tau )`$; $`\frac{}{\tau }|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^n=\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma |\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^n`$\]
$`\stackrel{}{E}_S(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`\stackrel{}{\pi }_S(\tau ,\stackrel{}{\sigma })={\displaystyle \frac{\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })}{\tau }}=`$ (351)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{(2m)!}}\dot{\stackrel{}{\eta }}_i(\tau )(\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}`$ (352)
$``$ $`{\displaystyle \frac{1}{(2m+2)!}}\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m+1}]=`$ (353)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))=`$ (354)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{\stackrel{}{\kappa }_i(\tau )\stackrel{}{}_\sigma }{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}}}\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))=`$ (355)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i\times `$ (359)
$`{\displaystyle \frac{1}{4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^2}}[\stackrel{}{\kappa }_i(\tau )\stackrel{}{\kappa }_i(\tau ){\displaystyle \frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|}}{\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}}{[m_i^2+(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2]^{3/2}}}+`$
$`+{\displaystyle \frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|}}({\displaystyle \frac{\stackrel{}{\kappa }_i^2(\tau )+(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2}{\stackrel{}{\kappa }_i^2(\tau )(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2}}({\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}}{\sqrt{m_i^2+(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2}}}1)+`$
$`+{\displaystyle \frac{(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}}{[m_i^2+(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2]^{3/2}}})].`$
The magnetic field is \[$`B_S^r=ϵ^{rsu}\frac{}{\sigma ^s}A_S^u`$\]
$`\stackrel{}{B}_S(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i}{\overset{1..N}{}}}Q_i\stackrel{}{B}_{Si}(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau ),\stackrel{}{\kappa }_i(\tau ))=`$ (360)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}\dot{\stackrel{}{\eta }}_i(\tau )\times \stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_i(\tau )\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^{2m1}=`$ (361)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}Q_i{\displaystyle \frac{1}{4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|^2}}{\displaystyle \frac{m_i^2\stackrel{}{\kappa }_i(\tau )\times \frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|}}{[m_i^2+(\stackrel{}{\kappa }_i(\tau )\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )|})^2]^{3/2}}}.`$ (362)
In these expressions $`\dot{\stackrel{}{\eta }}_i(\tau )`$ may be replaced with $`\stackrel{}{\kappa }_i(\tau )/\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}`$. We have also given the closed form of the fields.
Let us remark that the Feynman-Wheeler complete absorber assumption is violated in the Maxwell theory, since it would imply (see for instance Ref. , where there is the definition of radiation in the Feynman-Wheeler theory) $`\stackrel{}{E}_S(\tau ,\stackrel{}{\sigma })=\stackrel{}{E}_\pm (\tau ,\stackrel{}{\sigma })=0`$ and $`\stackrel{}{B}_S(\tau ,\stackrel{}{\sigma })=\stackrel{}{B}_\pm (\tau ,\stackrel{}{\sigma })=0`$ everywhere (inside and outside the absorbers).
For the energy we need $`\stackrel{}{E}_S^2+\stackrel{}{B}_S^2`$ and for the momentum we need $`\stackrel{}{E}_S\times \stackrel{}{B}_S`$. In Appendix B we evaluate these and show that \[$`\stackrel{}{\eta }_{ij}(\tau )=\eta _{ij}(\tau )\widehat{\eta }_{ij}(\tau )=\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )`$; $`\stackrel{}{}_{ij}=/\stackrel{}{\eta }_{ij}`$\]
$`U(\tau )={\displaystyle \frac{1}{2}}{\displaystyle d^3\sigma (\stackrel{}{E}_S^2+\stackrel{}{B}_S^2)(\tau ,\stackrel{}{\sigma })}:={\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}h_1(\dot{\stackrel{}{\eta }_i},\dot{\stackrel{}{\eta }_j},\stackrel{}{\eta }_{ij})=`$ (363)
$`=`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[\dot{\stackrel{}{\eta }_i}\dot{\stackrel{}{\eta }_j}{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m+1}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$ (365)
$``$ $`{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m+2}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}+`$ (366)
$`+`$ $`\dot{\stackrel{}{\eta }_i}\dot{\stackrel{}{\eta }_j}{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n}\stackrel{}{}_{ij}^2\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$ (367)
$``$ $`{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m+1}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}],`$ (368)
and
$`{\displaystyle d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })}:={\displaystyle \underset{i<j}{\overset{1..N}{}}}Q_iQ_j\stackrel{}{h}_1(\dot{\stackrel{}{\eta }_i},\dot{\stackrel{}{\eta }_j},\stackrel{}{\eta }_{ij})=`$ (369)
$`=`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}[{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(\stackrel{}{}_{ij}[\dot{\stackrel{}{\eta }_i}\dot{\stackrel{}{\eta }_j}{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m+1}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$ (371)
$``$ $`{\displaystyle \frac{(\dot{\stackrel{}{\eta }_i}\stackrel{}{}_{ij})^{2m+2}(\dot{\stackrel{}{\eta }_j}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+3}}{(2n+2m+4)!}}])+(ij)].`$ (372)
### B The complete Darwin Hamiltonian in terms of the old canonical variables
Our first aim is to express the invariant mass $`M`$, i.e. the Hamiltonian for the $`\tau T_s`$-evolution, in terms of the original canonical variables. Later we will obtain this Hamiltonian in terms of the new canonical variables. The Hamiltonian $`M`$ and the internal 3-momentum are
$`M`$ $`=`$ $`𝒫_{(int)}^\tau ={\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+(\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau )))^2}+`$ (373)
$`+`$ $`{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+{\displaystyle d^3\sigma \frac{1}{2}[\stackrel{}{\pi }_S^2(\tau ,\stackrel{}{\sigma })+\stackrel{}{B}_S^2(\tau ,\stackrel{}{\sigma })]}=`$ (374)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )\stackrel{}{V}_i(\tau )]^2}+{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi |\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )|}}+U(\tau )=`$ (375)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}+{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi |\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )|}}+V_{DAR}(\tau ),`$ (376)
$`V_{DAR}(\tau )`$ $`:=`$ $`V_{LOD}(\tau )+U_{HOD}(\tau ),`$ (378)
$`\stackrel{}{𝒫}_{(int)}`$ $`=`$ $`\stackrel{}{}_p(\tau )=\stackrel{}{\kappa }_+(\tau )+{\displaystyle d^3\sigma [\stackrel{}{\pi }_S\times \stackrel{}{B}_S](\tau ,\stackrel{}{\sigma })}0,`$ (380)
The first line of $`M`$ is
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+(\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau )))^2}={\displaystyle \underset{i=1}{\overset{N}{}}}\left(\sqrt{m_i^2+\stackrel{}{\kappa }_i(\tau )^2}{\displaystyle \frac{\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau ))}{\sqrt{m_i^2+\stackrel{}{\kappa }_i(\tau )^2}}}\right)=`$ (381)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}+U_1,`$ (384)
$`U_1={\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\stackrel{}{\kappa }_i\stackrel{}{V}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}}},V_{DAR}(\tau )=U_1(\tau )+U(\tau )=V_{LOD}(\tau )+U_{HOD}(\tau ),`$
in which the vector potential is given by the semiclassical Lienard-Wiechert transverse potential Eq.(306). By re-expressing this transverse vector potential in terms of the momenta \[one may use in that expression either the new or old canonical variables because of the Grassmann truncation\] we obtain
$`\stackrel{}{A}_S(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{(2m)!}}{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_\sigma )^{2m}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m1}`$ (386)
$`{\displaystyle \frac{1}{(2m+2)!}}\stackrel{}{}_\sigma ({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m+1}],`$
so that for the scalar potential $`U_1`$ we get
$`U_1`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\stackrel{}{\kappa }_i\stackrel{}{V}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}}}={\displaystyle \underset{i=1}{\overset{N}{}}}\left({\displaystyle \frac{\stackrel{}{\kappa }_iQ_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau ))}{\sqrt{m_i^2+\stackrel{}{\kappa }_i(\tau )^2}}}\right)=`$ (387)
$`=`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}[({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}})`$ (391)
$`\left(({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m}+({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m}\right){\displaystyle \frac{\eta _{ij}^{2m1}}{(2m)!}}`$
$`(({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m+1}+`$
$`+({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }j^2}}}\stackrel{}{}_{ij})^{2m+1}){\displaystyle \frac{\eta _{ij}^{2m+1}}{(2m+2)!}}].`$
The lowest order part of this kinetic contribution (the $`m=0`$ term) is twice the familiar Darwin interaction (but with $`m_i\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}`$; strictly speaking this is a higher order correction). Thus we have
$$U_1=\underset{i=1}{\overset{N}{}}\frac{\stackrel{}{\kappa }_i\stackrel{}{V}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}}=2V_{LOD}(\tau )+U_{1HOD},$$
(392)
with
$`V_{LOD}={\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta _{ij}}}`$ (394)
$`({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}){\displaystyle \frac{\eta _{ij}}{2}})=`$
$`=`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{8\pi \eta _{ij}}}\left({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}+({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\widehat{\eta }_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\widehat{\eta }_{ij})\right).`$ (395)
The standard form for the historical Darwin term is the above but with $`\sqrt{m_i^2+\kappa _i^2}m_i`$.
The remaining compensating part of the familiar Darwin interaction plus all higher order parts come from the field energy $`U(\tau ).`$ In terms of momentum variables, using Grassmann truncations, the expression Eq.(368 ) for the field energy integral simply ) becomes
$`U(\tau )={\displaystyle \frac{1}{2}}{\displaystyle d^3\sigma (\stackrel{}{E}_S^2+\stackrel{}{B}_S^2)(\tau ,\stackrel{}{\sigma })}={\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}h_1({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}},\stackrel{}{\eta }_{ij})=`$ (396)
$`={\displaystyle \underset{i<j}{\overset{N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\times `$ (397)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+1}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$ (398)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\eta _{ij}^{2n+2m+3}}{(2n+2m+4)!}}+`$ (399)
$`+{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n}\eta _{ij}^{2n+2m1}}{(2n+2m)!}}`$ (400)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+1}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}].`$ (401)
Single infinite sum pieces can be split off from the double infinite sum in the last two lines of the above expression for the field energy integral. This naturally separates out the compensating portion of the familiar lowest order Darwin parts plus all remaining higher order Darwin parts, including a piece that cancels exactly $`U_{1HOD}.`$
$`U(\tau )={\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}h_1({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}},\stackrel{}{\eta }_{ij})=`$ (402)
$`=`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}\left({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta _{ij}}}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}){\displaystyle \frac{\eta _{ij}}{2}}\right)+`$ (409)
$`+{\displaystyle \underset{i<j}{\overset{N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}`$
$`[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+1}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\eta _{ij}^{2n+2m+3}}{(2n+2m+4)!}}]+`$
$`+{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}[{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}`$
$`(({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m}+({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m}))\eta _{ij}^{2m1}`$
$`{\displaystyle \frac{1}{(2m+2)!}}({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m+1}+`$
$`+`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m+1})\eta _{ij}^{2m+1}]+`$ (413)
$`+{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\times `$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\eta _{ij}^{2n+2m+3}}{(2n+2m+4)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+3}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+3}\eta _{ij}^{2n+2m+5}}{(2n+2m+6)!}}]=`$
$`=`$ $`V_{LOD}U_{1HOD}+U_{HOD},`$ (416)
$`V_{DAR}(\tau )=U_1(\tau )+U(\tau )=V_{LOD}(\tau )+U_{HOD}(\tau ).`$
In this form we see that the first summation is $`V_{LOD}`$ and combines with the gauge part of the kinetic piece \[$`2V_{LOD}(\tau )`$\] to give the familiar lowest order Darwin piece $`V_{LOD}(\tau )`$. The third set of summations is $`U_{1HOD}`$ and exactly cancels with the corresponding term in the kinetic piece. The second and fourth set of summations ( the double sums which begin at higher order in $`(1/c^2)`$) we define as $`U_{HOD}`$, coming only from $`U(\tau )`$. Altogether we obtain \[$`\eta _{ij}=|\stackrel{}{\eta }_{ij}|=|\stackrel{}{\eta }_i\stackrel{}{\eta }_j|`$\]
$`M`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+(\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau )))^2}+`$ (417)
$`+`$ $`{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}+{\displaystyle d^3\sigma \frac{1}{2}[\stackrel{}{\pi }_S^2+\stackrel{}{B}_S^2](\tau ,\stackrel{}{\sigma })}=`$ (418)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i(\tau )\stackrel{}{V}_i(\tau )]^2}+{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi |\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j)(\tau )|}}+U(\tau )=`$ (419)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}+{\displaystyle \underset{ij}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{|\stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )|}}+V_{LOD}(\tau )+V_{HOD}(\tau )=`$ (420)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i(\tau )^2}+{\displaystyle \underset{ij}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{}{\eta }_i(\tau )\stackrel{}{\eta }_j(\tau )}}`$ (422)
$``$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}\left({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta _{ij}}}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}){\displaystyle \frac{\eta _{ij}}{2}}\right)+`$ (423)
$`+`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\times `$ (428)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+1}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+1}\eta _{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}+`$
$`+{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\eta _{ij}^{2n+2m+3}}{(2n+2m+4)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+3}(\frac{\stackrel{}{\kappa }_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+3}\eta _{ij}^{2n+2m+5}}{(2n+2m+6)!}}].`$
In Appendix C we show how the multiple directional derivatives in the generalized higher order Darwin interactions can be evaluated in general, thereby obtaining a more readily usable form. (The expression for $`M`$ and for $`V_{DAR}`$ for arbitrary $`N`$ is an immediate generalization of the case for $`N=2`$ and for simplicity of notation the details in the appendices will be limited to the two-body case.) However, because of the significant complexity of the above form it will be of value to first obtain an alternative derivation of this series. For $`N`$=2 we have the following Lagrangian expression for the invariant mass \[see Eq.(229) with $`\stackrel{}{\lambda }(\tau )=0`$; $`h_1`$ is defined in Eq.(368)\]
$`E_{rel}`$ $`=`$ $`h(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })={\displaystyle \frac{m_1}{\sqrt{1\dot{\stackrel{}{\eta }_1}^2}}}+{\displaystyle \frac{m_2}{\sqrt{1\dot{\stackrel{}{\eta }_2}^2}}}+{\displaystyle \frac{Q_1Q_2}{4\pi |\stackrel{}{\eta }|}}+{\displaystyle \frac{Q_1Q_2}{4\pi }}h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta }):=`$ (430)
$`:`$ $`=h_0(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })+{\displaystyle \frac{Q_1Q_2}{4\pi }}h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta }).`$ (431)
In order to find the Hamiltonian $`H(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2;\stackrel{}{\eta })`$ from $`h_1`$ we must demand that Hamilton’s equation be satisfied . We use the Dirac bracket since we have used the constraint as a strong condition on the dynamical variables. This will lead to a set of differential equations for $`H(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2;\stackrel{}{\eta })`$ and to a result that agrees exactly with the above expression for $`M`$. The details of this analysis are given in Appendix D.
Here we mention a comparison of these cross checked results (valid to all order of $`1/c^2`$) with approximate results obtained elsewhere. In one obtains a single time Lagrangian by expanding the symmetric Green function in the Fokker action, used by Wheeler and Feynman, to all orders in $`1/c^2`$. From this Lagrangian one obtains the Legendre Hamiltonian $`\stackrel{~}{h}(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })+`$ terms involving the acceleration and all higher order derivatives. Ignoring those higher order accelerations one finds the same Legendre Hamiltonian as above. From that Hamiltonian the authors obtain a final Hamiltonian that although agreeing through order $`1/c^2`$ with the results above (including the standard Darwin interaction to that order) they differ from our common results (above and in Appendix D) at order $`1/c^4`$ (no terms of higher order are computed in ). The failure to obtain results that agree with our result here is the neglect there of using the proper Dirac brackets in the Kerner reduction. The reason that those Dirac brackets were not used is that the authors took as a starting point the Fokker action, not taking into account that this action itself is a result of imposing constraints on the solutions of the electromagnetic field equations.
By using Eqs.(349) we get
$`\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}`$ $`=`$ $`\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2+\stackrel{~}{\stackrel{}{\kappa }_i}{\displaystyle \underset{ji}{}}Q_iQ_j\stackrel{}{}_{\eta _i}𝒦_{ij}}=`$ (432)
$`=`$ $`\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}+{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }_i}_{ji}Q_iQ_j\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij}}{2\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}}},`$ (434)
with
$`Q_iQ_j𝒦_{ij}(\stackrel{}{\kappa }_{i,}\stackrel{}{\kappa }_{j,}\stackrel{}{\eta }_i\stackrel{}{\eta }_j)`$ $`=`$ $`Q_iQ_j\stackrel{~}{𝒦}_{ij}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}),`$ (435)
$`\stackrel{~}{\stackrel{}{\kappa }_i}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij}`$ $`=`$ $`{\displaystyle }d^3\sigma [(\stackrel{~}{\stackrel{}{\kappa }_i}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{}{A}_{Si}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}))\stackrel{}{\pi }_{Si}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})`$ (437)
$`\stackrel{}{A}_{Si}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})(\stackrel{~}{\stackrel{}{\kappa }_i}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{}{\pi }_{Si}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}))],`$
so that $`M=𝒫_{(int)}^\tau `$ becomes (due to Grassmann truncation we can replace the old variables by the new variables in the interaction terms)
$`M=𝒫_{(int)}^\tau ={\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}+{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }_i}_{ji}Q_iQ_j\stackrel{}{}_{\stackrel{~}{\eta }_i}𝒦_{ij}}{2\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}}}+{\displaystyle \underset{ij}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}}}`$
$``$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}\left({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_j^2}}}{\displaystyle \frac{1}{\stackrel{~}{\eta }_{ij}}}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_j^2}}}\stackrel{}{}_{ij}){\displaystyle \frac{\stackrel{~}{\eta }_{ij}}{2}}\right)+`$ (438)
$`+`$ $`{\displaystyle \underset{i<j}{\overset{1..N}{}}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{(}\stackrel{}{\kappa }_i}^2}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\times `$ (443)
$`{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+1}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+1}\stackrel{~}{\eta }_{ij}^{2n+2m+1}}{(2n+2m+2)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\stackrel{~}{\eta }_{ij}^{2n+2m+3}}{(2n+2m+4)!}}+`$
$`+{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\stackrel{~}{\eta }_{ij}^{2n+2m+3}}{(2n+2m+4)!}}`$
$`{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+3}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+3}\stackrel{~}{\eta }_{ij}^{2n+2m+5}}{(2n+2m+6)!}}].`$
where $`\stackrel{}{}_{ij}=/\stackrel{~}{\eta }_{ij}`$. We emphasize at this point that the interaction terms are all two-body forces; there are no $`N`$ body forces in our semiclassical treatment.
### C <br>Complete Darwin Hamiltonian in terms of the Final Canonical Variables for the N-Body Problem. <br>
In this subsection we will present the final expression for $`M=𝒫_{(int)}^\tau `$ in the final canonical variables for arbitrary $`N`$ . Now from Eq.(345)above we can replace $`_{i<j}_{\eta _i}Q_iQ_j𝒦_{ij}`$ by $`d^3\sigma [\stackrel{}{\pi }_S\times \stackrel{}{B}_S](\tau ,\stackrel{}{\sigma })`$. We have already performed that integral (see Eq.(LABEL:d14)) for the two-body problem and the result for the $`N`$ body problem is an immediate generalization. Using that complete expression we obtain \[$`\eta _{ij}=|\stackrel{}{\eta }_{ij}|=|\stackrel{}{\eta }_i\stackrel{}{\eta }_j|`$\]
$`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i(\tau )^2}=\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}(\tau )^2}+U_{HOD}^{}(\tau )`$ (451)
$`U_{HOD}^{}(\tau )={\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{8\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{(2n+2m+2)!}}`$
$`\times (({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+2}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n}+`$
$`+2({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+1}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+1}+`$
$`({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+2})\stackrel{~}{\eta }_{ij}^{2n+2m+1}`$
$`{\displaystyle \frac{1}{(2n+2m+4)!}}(({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+3}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+1}+`$
$`+2({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+2}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+2}+`$
$`+`$ $`({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+1}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+3})\stackrel{~}{\eta }_{ij}^{2n+2m+3}].`$ (452)
We now combine this portion of the complete Darwin Hamiltonian coming from the vector potentials $`\stackrel{}{V}_i(\tau )`$ in the kinetic terms with the potential $`U(\tau )`$ arising from the field energy integral \[which contains $`V_{LOH}(\tau )`$\] to obtain the final form of the invariant mass in the rest frame
$`M`$ $`=`$ $`𝒫_{(int)}^\tau ={\displaystyle \underset{i=1}{\overset{M}{}}}\sqrt{m_i^2+(\stackrel{}{\kappa }_i(\tau )Q_i\stackrel{}{A}_S(\tau ,\stackrel{}{\eta }_i(\tau )))^2}+`$ (453)
$`+`$ $`{\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \eta _{ij}}}+{\displaystyle d^3\sigma \frac{1}{2}[\stackrel{}{\pi }_S^2+\stackrel{}{B}_S^2](\tau ,\stackrel{}{\sigma })}=`$ (454)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+[\stackrel{}{\kappa }_i^2(\tau )\stackrel{}{V}_i(\tau )]^2}+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \eta _{ij}}}+U(\tau )=`$ (455)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{}{\kappa }_i^2(\tau )}+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \eta _{ij}}}+V_{DAR}(\tau )=[V_{DAR}=V_{LOD}+V_{HOD}]`$ (456)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}+{\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi \eta _{ij}}}+\stackrel{~}{V}_{DAR}(\tau ),[\stackrel{~}{V}_{DAR}=V_{DAR}+U_{HOD}^{}]`$ (458)
$`\stackrel{~}{V}_{DAR}=`$ $`{\displaystyle \underset{i<j}{}}{\displaystyle \frac{Q_iQ_j}{4\pi }}\left({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\stackrel{~}{\eta }}}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij}){\displaystyle \frac{\stackrel{~}{\eta }_{ij}}{2}}\right)`$ (460)
$``$ $`{\displaystyle \frac{Q_iQ_j}{8\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{(2n+2m+2)!}}`$ (462)
$`\times (({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m+2}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2n}+`$
$`+`$ $`({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2n+2})\stackrel{~}{\eta }_{ij}^{2n+2m+1}`$ (463)
$``$ $`{\displaystyle \frac{1}{(2n+2m+4)!}}(({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2m+3}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{}_{ij})^{2n+1}+`$ (464)
$`+`$ $`({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2m+1}({\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}\stackrel{}{_{ij}})^{2n+3})\stackrel{~}{\eta }_{ij}^{2n+2m+3}]`$ (465)
$``$ $`2{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+2}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{~}{\stackrel{}{\kappa }}_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+2}\stackrel{~}{\eta }_{ij}^{2n+2m+3}}{(2n+2m+4)!}}+`$ (466)
$`+`$ $`2{\displaystyle \frac{(\frac{\stackrel{~}{\stackrel{}{\kappa }}_i}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_{i}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2m+3}(\frac{\stackrel{~}{\stackrel{}{\kappa }}_j}{\sqrt{m_j^2+\stackrel{}{\kappa }_{j}^{}{}_{}{}^{2}}}\stackrel{}{}_{ij})^{2n+3}\stackrel{~}{\eta }_{ij}^{2n+2m+5}}{(2n+2m+6)!}}].`$ (467)
Notice that the $`N`$ body interaction Hamiltonian is a sum of individual 2 body Hamiltonians.(Grassmann truncation eliminates $`N`$-body forces). In this sense these interactions correspond to a sum of disconnected and spectator Feynman diagrams.
### D Further Reductions of the N-Body Darwin Hamiltonian and Closed Form Solutions for the Two-Body Problem.
It is of interest to see if we can make further simplifications by deriving an expression for the multiple derivatives. We show in Appendix C that the problematic expression $`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^a(\stackrel{}{\kappa }_j\stackrel{}{}_{ij})^b\stackrel{~}{\eta }^{a+b1}`$ is obtained in terms of powers of \[$`\stackrel{~}{\eta }_{ij}=|\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}|;\widehat{\eta }_{ij}=(\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})/\stackrel{~}{\eta }_{ij}`$\]
$$\mathrm{cos}^2\varphi _{ij}:=\frac{(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{\kappa }_i\widehat{\eta }_{ij}\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2}{(\stackrel{}{\kappa }_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)}.$$
(468)
For $`N=2`$ this expression reduces to 1 if $`\stackrel{~}{\stackrel{}{\kappa }}_1=\stackrel{~}{\stackrel{}{\kappa }}_2`$ (the center of mass rest frame condition for the two body problem) . The result for $`a=2m+1,b=2n+1`$ is
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m+1}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n+1}\stackrel{~}{\eta }_{ij}^{2(m+n)+1}=`$ (469)
$`=`$ $`2{\displaystyle \frac{[(2m+2n+1)!!]^2[(n+m+1)!]^2}{(2(m+n+1))!}}`$ (473)
$`{\displaystyle \frac{(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij}\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^m(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^n}{\stackrel{~}{\eta }_{ij}}}\times `$
$`{\displaystyle \underset{l=0}{\overset{n}{}}}{\displaystyle \underset{k=0}{\overset{nl}{}}}{\displaystyle \underset{h=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{l}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{l+mn}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2(nl)+1}{2k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{k}{h}}\right)(1)^{k+h}(\mathrm{cos}^2\varphi _{ij})^{nl+hk}.`$
For the two body case in the center-of-mass rest frame $`\stackrel{~}{\stackrel{}{\kappa }}_1=\stackrel{~}{\stackrel{}{\kappa }}_2`$ $`:=\stackrel{~}{\stackrel{}{\kappa }}`$ this expression reduces to
$$\frac{[(2m+2n+1)!!]^2}{\stackrel{~}{\eta }}(\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2)^{m+n+1}.$$
(474)
In the general case this expression can be more simply written in terms of $`\mathrm{cos}(2k+1)\varphi _{ij}`$ \[see Eq.(C96)\]. Using the identity
$$\frac{[(2m+2n+1)!!][(n+m+1)!]}{(2(m+n+1))!}=\frac{1}{2^{m+n+1}},$$
(475)
one obtains for $`mn`$ the expression
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m+1}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n+1}\stackrel{~}{\eta }_{ij}^{2(m+n)+1}=`$ (476)
$`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{\eta }_{ij}}}{\displaystyle \frac{(2(m+n+1))!(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^{m+1/2}(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^{n+1/2}}{2^{m+n+2}}}`$ (479)
$`\times {\displaystyle \underset{k=0}{\overset{n}{}}}\mathrm{cos}(2k+1)\varphi _{ij}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{nk}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{mk}}\right)`$
Hence
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m+3}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n+1}\stackrel{~}{\eta }_{ij}^{2(m+n)+1}=`$ (480)
$`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{\eta }_{ij}}}{\displaystyle \frac{(2(m+n+2))!(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^{m+3/2}(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^{n+1/2}}{2^{m+n+3}}}`$ (483)
$`\times {\displaystyle \underset{k=0}{\overset{n}{}}}\mathrm{cos}(2k+1)\varphi \left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{nk}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+3}{m+1k}}\right)`$
By similar methods one finds that
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m+2}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n+2}\stackrel{~}{\eta }_{ij}^{2(m+n)+3}=`$ (484)
$`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{\eta }_{ij}}}{\displaystyle \frac{(2(m+n+2))!(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^{m+1}(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^{n+1}}{2^{m+n+4}}}\times `$ (487)
$`({\displaystyle \underset{k=0}{\overset{n}{}}}\mathrm{cos}(2k+1)\varphi _{ij}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{nk}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{mk}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{2m+2}{m+1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+2}{n+1}}\right)`$
To determine the other combination note
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n}\stackrel{~}{\eta }_{ij}^{2(m+n)1}=`$ (488)
$`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{\eta }_{ij}}}{\displaystyle \frac{(2(m+n))!(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^m(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^n}{2^{m+n+2}}}\times `$ (491)
$`({\displaystyle \underset{k=0}{\overset{n1}{}}}\mathrm{cos}(2k+1)\varphi _{ij}\left({\displaystyle \genfrac{}{}{0pt}{}{2n1}{n1k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m1}{m1k}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{2m}{m}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)`$
Thus
$`(\stackrel{~}{\stackrel{}{\kappa }}_i\stackrel{}{}_{ij})^{2m+2}(\stackrel{~}{\stackrel{}{\kappa }}_j\stackrel{}{}_{ij})^{2n}\stackrel{~}{\eta }_{ij}^{2(m+n)+1}=`$ (492)
$`=`$ $`{\displaystyle \frac{2}{\stackrel{~}{\eta }_{ij}}}{\displaystyle \frac{(2(m+n+1))!(\stackrel{~}{\stackrel{}{\kappa }}_i^2(\stackrel{~}{\stackrel{}{\kappa }}_i\widehat{\eta }_{ij})^2)^{m+1}(\stackrel{~}{\stackrel{}{\kappa }}_j^2(\stackrel{~}{\stackrel{}{\kappa }}_j\widehat{\eta }_{ij})^2)^n}{2^{m+n+3}}}\times `$ (495)
$`{\displaystyle \underset{k=0}{\overset{n1}{}}}\mathrm{cos}(2k+1)\varphi _{ij}\left({\displaystyle \genfrac{}{}{0pt}{}{2n1}{n1k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{m+1k}}\right)+\left({\displaystyle \genfrac{}{}{0pt}{}{2m+2}{m+1}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n}{n}}\right)`$
Unfortunately, the $`k`$ sums cannot be performed analytically to known closed forms for the $`N`$-body problem.
In the special case of the two body system we can obtain a closed form if we use the rest frame condition $`\stackrel{~}{\stackrel{}{\kappa }}_1+\stackrel{~}{\stackrel{}{\kappa }}_2=0`$ . The expression we get in this way may be used with the Dirac brackets associated with $`\stackrel{~}{\stackrel{}{\kappa }}_1+\stackrel{~}{\stackrel{}{\kappa }}_20,\stackrel{}{q}_+0`$, so that the final reduced phase contains only $`\stackrel{~}{\eta }=|\stackrel{~}{\stackrel{}{\eta }_1}\stackrel{~}{\stackrel{}{\eta }_2}|`$ and $`\stackrel{~}{\stackrel{}{\kappa }}:=\stackrel{~}{\stackrel{}{\kappa }}_1=\stackrel{~}{\stackrel{}{\kappa }}_2.`$ Using the identity in Eq.(474) the higher order Darwin part $`V_{HOD}`$ $`=U_{HOD}+U_{HOD\text{ }}^{}`$ becomes
$`V_{HOD}`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi \stackrel{~}{\eta }}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[\stackrel{~}{\stackrel{}{\kappa }}^2{\displaystyle \frac{[(2m+2n+1)!!]^2}{(2n+2m+2)!}}[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{n+m+1}`$ (501)
$`\times ({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2m+1}({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2n+1}({\displaystyle \frac{1}{m_1^2+\stackrel{}{\kappa }^2}}+{\displaystyle \frac{1}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})+`$
$`+{\displaystyle \frac{[(2m+2n+3)!!]^2}{(2n+2m+4)!}}`$
$`[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{n+m+2}(({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2m+1}({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2n+1}({\displaystyle \frac{1}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}+{\displaystyle \frac{1}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})+`$
$`+2\stackrel{~}{\stackrel{}{\kappa }}^2{\displaystyle \frac{[(2m+2n+3)!!]^2}{(2n+2m+4)!}}[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta }]^2)^{n+m+2}({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{}{\kappa }^2}}})^{2m+3}({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2n+3}`$
$`2{\displaystyle \frac{[(2m+2n+5)!!]^2}{(2n+2m+6)!}}[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{n+m+3}({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2m+3}({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})^{2n+3}].`$
We use
$$\frac{[(2m+2n+1)!!]^2}{(2n+2m+2)!}=\frac{()^{n+m}}{2(n+m+1)}\left(\genfrac{}{}{0pt}{}{3/2}{n+m}\right),$$
(503)
and let $`m+n=l`$ so that $`0ml`$ and $`0l<\mathrm{}`$. Then we perform the $`m`$ sum using
$$\underset{m=0}{\overset{l}{}}\left(\frac{x}{y}\right)^m=\frac{y^{l+1}x^{l+1}}{y^l(yx)},$$
(504)
and obtain
$`V_{HOD}`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi \stackrel{~}{\eta }}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}[\stackrel{~}{\stackrel{}{\kappa }}^2{\displaystyle \frac{()^l}{2(l+1)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l}}\right)[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{l+1}[{\displaystyle \frac{(\frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}(\frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}}{m_1^2m_2^2}}]`$ (512)
$`\times \sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}({\displaystyle \frac{1}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}+{\displaystyle \frac{1}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})+`$
$`+{\displaystyle \frac{()^{l+1}}{2(l+2)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l+1}}\right)[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta }]^2)^{l+2}[{\displaystyle \frac{(\frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}(\frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}}{m_1^2m_2^2}}]`$
$`\times \sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}({\displaystyle \frac{1}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}+{\displaystyle \frac{1}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})+`$
$`+\stackrel{~}{\stackrel{}{\kappa }}^2{\displaystyle \frac{()^{l+1}}{(l+2)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l+1}}\right)[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{l+2}\times `$
$`[{\displaystyle \frac{(\frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}(\frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}}{m_1^2m_2^2}}]({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})`$
$`{\displaystyle \frac{()^l}{(l+3)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l+2}}\right)[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]^{l+3}\times `$
$`[{\displaystyle \frac{(\frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}(\frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}})^{2l+2}}{m_1^2m_2^2}}]({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})].`$
We use
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{2(l+1)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l}}\right)x^l`$ $`=`$ $`{\displaystyle \frac{1}{x}}((1x)^{1/2}1),`$ (513)
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^{l+1}}{2(l+2)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l+1}}\right)x^l`$ $`=`$ $`{\displaystyle \frac{1}{x^2}}(2(1x)^{1/2}2x)={\displaystyle \frac{1}{x^2}}(2(1x)^{1/2}+(1x)3),`$ (514)
$`{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{()^l}{2(l+3)}}\left({\displaystyle \genfrac{}{}{0pt}{}{3/2}{l+2}}\right)x^l`$ $`=`$ $`{\displaystyle \frac{1}{x^3}}((1x)^{1/2}{\displaystyle \frac{3}{8}}(1x)^2+{\displaystyle \frac{5}{4}}(1x){\displaystyle \frac{15}{8}}),`$ (515)
and finally determine
$`V_{HOD}`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi \stackrel{~}{\eta }(m_1^2m_2^2)}}\times `$ (524)
$`(\stackrel{~}{\stackrel{}{\kappa }}^2[\sqrt{{\displaystyle \frac{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}\sqrt{{\displaystyle \frac{m_1^2+\stackrel{}{\kappa }^2}{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}](\sqrt{{\displaystyle \frac{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}}+\sqrt{{\displaystyle \frac{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})+`$
$`+[(m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2)[2\sqrt{{\displaystyle \frac{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}+{\displaystyle \frac{m_2^2+(\stackrel{}{\kappa }\widehat{\eta })^2}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}3]`$
$`(m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2)[2\sqrt{{\displaystyle \frac{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}+{\displaystyle \frac{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}3]]\times `$
$`(\sqrt{{\displaystyle \frac{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}}+\sqrt{{\displaystyle \frac{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})+`$
$`+2\kappa ^2[(m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2)[2\sqrt{{\displaystyle \frac{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}+{\displaystyle \frac{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}3]`$
$`(m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2)[2\sqrt{{\displaystyle \frac{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}+{\displaystyle \frac{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}3]]\times `$
$`({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})({\displaystyle \frac{1}{m_2^2+\sqrt{\stackrel{~}{\stackrel{}{\kappa }}^2}}})`$
$`2[(m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2)^2(\sqrt{{\displaystyle \frac{m_2^2+\stackrel{}{\kappa }^2}{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}{\displaystyle \frac{3}{8}}\left({\displaystyle \frac{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_2^2+\stackrel{}{\kappa }^2}}\right)^2+`$
$`+`$ $`{\displaystyle \frac{5}{4}}\left({\displaystyle \frac{m_2^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}\right){\displaystyle \frac{15}{8}})`$ (526)
$`(m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2)^2(\sqrt{{\displaystyle \frac{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}}{\displaystyle \frac{3}{8}}\left({\displaystyle \frac{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}\right)^2+`$
$`+`$ $`{\displaystyle \frac{5}{4}}\left({\displaystyle \frac{m_1^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}\right){\displaystyle \frac{15}{8}})]\times `$ (528)
$`({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})).`$
So our final two-body expression is
$`M`$ $`=`$ $`\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}+\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}+{\displaystyle \frac{Q_1Q_2}{4\pi \stackrel{~}{\eta }}}+\stackrel{~}{V}_{DAR},`$ (530)
$`\stackrel{~}{V}_{DAR}=V_{LOD}+V_{HOD},V_{HOD}=U_{HOD}+U_{HOD}^{}`$
$`V_{LOD}`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi \eta }}[\stackrel{~}{\stackrel{}{\kappa }}^2(\stackrel{}{\kappa }\widehat{\eta })^2]({\displaystyle \frac{1}{\sqrt{m_1^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}})({\displaystyle \frac{1}{\sqrt{m_2^2+\stackrel{~}{\stackrel{}{\kappa }}^2}}}).`$ (532)
in which we have used the rest-frame condition $`\stackrel{~}{\stackrel{}{\kappa }}_1=\stackrel{~}{\stackrel{}{\kappa }}_2:=\stackrel{~}{\stackrel{}{\kappa }}`$. Note that $`Q_1Q_2\stackrel{}{\kappa }=Q_1Q_2\stackrel{~}{\stackrel{}{\kappa }}.`$
In the equal mass limit $`m_1m_2:=m`$
$`V_{HOD}=+{\displaystyle \frac{Q_1Q_2}{8\pi \stackrel{~}{\eta }}}\times `$ (533)
$`{\displaystyle \frac{m^2[3\stackrel{~}{\stackrel{}{\kappa }}^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]2\stackrel{}{\kappa }^2[\stackrel{~}{\stackrel{}{\kappa }}^23(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]\sqrt{\frac{m^2+\kappa ^2}{m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}[3\stackrel{~}{\stackrel{}{\kappa }}^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2][m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]}{(m^2+\kappa ^2)[m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]}},`$ (534)
so that
$`M=𝒫_{(int)}^\tau =2\sqrt{m^2+\stackrel{~}{\stackrel{}{\kappa }}^2}+{\displaystyle \frac{Q_1Q_2}{4\pi \stackrel{~}{\eta }}}+{\displaystyle \frac{Q_1Q_2}{8\pi \stackrel{~}{\eta }}}\times `$ (536)
$`({\displaystyle \frac{m^2[3\stackrel{~}{\stackrel{}{\kappa }}^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]2\stackrel{~}{\stackrel{}{\kappa }}^2[\stackrel{~}{\stackrel{}{\kappa }}^23(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]\sqrt{\frac{m^2+\stackrel{~}{\stackrel{}{\kappa }}^2}{m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2}}2[\stackrel{~}{\stackrel{}{\kappa }}^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2][m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]}{(m^2+\stackrel{~}{\stackrel{}{\kappa }}^2)[m^2+(\stackrel{~}{\stackrel{}{\kappa }}\widehat{\eta })^2]}}).`$ (537)
Our result differs from the appropriate terms at order $`1/c^4`$ with those in , and although it does agree with their Darwin interaction at order $`1/c^2`$. Since each of these latter three sources took as a starting point the Fokker particle Lagrangian they have not used, as we have done here, the canonical variables which would come from using the pair of secondary constraints arising from the reduction of the field plus particle Lagrangian to the Fokker action. We point out that Molina et al in , like earlier work by Golubenkov and Smorodinski obtain order $`1/c^4`$ corrections to the Hamiltonian. However, unlike them these authors not only use the Coulomb force law to replace acceleration dependent terms on the Lagrangian but also include the effects of that substitution on the choice of canonical variables by viewing it as a constraint (thus including Dirac brackets), in the reduction to Hamiltonian forms. (Note that since these three approaches, unlike ours, do not use Grassmann charges, they contain acceleration driven terms not contained in our approach. The comparisons we are talking about here with our approach refer to the terms not driven by acceleration dependent Lagrangian potentials).
It is of interest that Hamilton’s equations has a solution for circular orbits just as does the Schild solution for Feynman-Wheeler electrodynamics . In Appendix E we show how this comes about for equal masses. The case of unequal masses is similar.
### E $`𝒥_{(int)}^r`$, $`𝒦_{(int)}^r`$, $`\stackrel{}{q}_+`$ and the Energy-Momentum Tensor in the Final Canonical Variables.
Eqs.(346) and (Eq.(467) \[(Eq.(532)) for $`N=2`$\] give $`\stackrel{}{𝒫}_{(int)}`$ and $`M=𝒫_{(int)}^\tau `$ in terms of the final canonical variables. For the internal angular momentum we get
$`𝒥_{(int)}^r`$ $`=`$ $`\epsilon ^{rst}\overline{S}_s^{st}={\displaystyle \underset{i=1}{\overset{N}{}}}(\stackrel{}{\eta }_i(\tau )\times \stackrel{}{\kappa }_i(\tau ))^r+{\displaystyle d^3\sigma (\stackrel{}{\sigma }\times [\stackrel{}{\pi }_{}\times \stackrel{}{B}](\tau ,\stackrel{}{\sigma }))^r}=`$ (539)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}[(\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{}{\alpha }_i)\times (\stackrel{~}{\stackrel{}{\kappa }_i}\stackrel{}{\beta }_i)]^r+`$ (540)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{\overset{N}{}}}Q_iQ_j[{\displaystyle }d^3\sigma (\stackrel{}{\sigma }\times [\stackrel{}{\pi }_{Si}\times (\stackrel{}{}_\sigma \times \stackrel{}{A}_{Sj})]+ij).`$ (541)
Using the forms for $`\stackrel{}{\alpha }_i`$ and $`\stackrel{}{\beta }_i`$ given in Eqs.(349) together with the expression for $`𝒦_{ij}`$ and expanding the cross products in the integral we obtain
$`𝒥_{(int)}^r`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}(\stackrel{~}{\stackrel{}{\eta }_i}\times \stackrel{~}{\stackrel{}{\kappa }_i})^r+`$ (545)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{\overset{N}{}}}Q_iQ_j\{(\stackrel{}{\eta }_i\times _{\eta _i}+\stackrel{}{\kappa }_i\times _{\kappa _i})^r{\displaystyle d^3\sigma (\stackrel{}{A}_{Si}\stackrel{}{\pi }_{Sj}\stackrel{}{A}_{Sj}\stackrel{}{\pi }_{Si})}\}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{\overset{N}{}}}Q_iQ_j[{\displaystyle }d^3\sigma \epsilon ^{ruv}\sigma ^u[\stackrel{}{\pi }_{Si}_\sigma ^v\stackrel{}{A}_{Sj}+\stackrel{}{\pi }_{Sj}_\sigma ^v\stackrel{}{A}_{Si}`$
$`\stackrel{}{\pi }_{Si}\stackrel{}{}_\sigma A^v{}_{Sj}{}^{}\stackrel{}{\pi }_{Sj}\stackrel{}{}_\sigma A^v{}_{Si}{}^{}].`$
Using the transverse nature of the field together with vanishing surface terms we find that
$$\stackrel{}{𝒥}_{(int)}=\underset{i=1}{\overset{N}{}}\stackrel{~}{\stackrel{}{\eta }_i}\times \stackrel{~}{\stackrel{}{\kappa }_i},$$
(546)
if
$`(\stackrel{}{\eta }_i\times \stackrel{}{}_{\eta _i}+\stackrel{}{\kappa }_i\times \stackrel{}{}_{\kappa _i}){\displaystyle d^3\sigma (\stackrel{}{A}_{Si}\stackrel{}{\pi }_{Sj}\stackrel{}{A}_{Sj}\stackrel{}{\pi }_{Si})}`$ (547)
$`=`$ $`{\displaystyle }d^3\sigma [\stackrel{}{A}_{Si}\times \stackrel{}{\pi }_{Sj}\stackrel{}{\sigma }\times (\stackrel{}{}_\sigma A^k{}_{Si}{}^{})\pi _{Sj}^k+\stackrel{}{A}_{Sj}\times \stackrel{}{\pi }_{Si}+\stackrel{}{\sigma }\times (A^k{}_{Si}{}^{}\stackrel{}{}_{\sigma }^{}\pi _{Sj}^k)].`$ (548)
We have given explicit forms forms for $`\stackrel{}{A}_{Si}`$ and $`\stackrel{}{\pi }_{Si}`$ in Eqs.(311) and (359) respectively. Their general forms are
$`\stackrel{}{A}_{Si}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}[\stackrel{}{\kappa }_if_i(\kappa _i^2,\stackrel{}{\kappa }_i\widehat{\rho }_i)+\widehat{\rho }_ig_i(\kappa _i^2,\stackrel{}{\kappa }_i\widehat{\rho }_i)],`$ (549)
$`\stackrel{}{\pi }_{Si}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi |\stackrel{}{\sigma }\stackrel{}{\eta }_i|^2}}[\stackrel{}{\kappa }_ih_i(\kappa _i^2,\stackrel{}{\kappa }_i\widehat{\rho }_i)+\widehat{\rho }_ic_i(\kappa _i^2,\stackrel{}{\kappa }_i\widehat{\rho }_i)].`$ (550)
in which $`\widehat{\rho }_i:=\frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}`$ . Using
$$\stackrel{}{}_\sigma \frac{1}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^n}=\frac{n\widehat{\rho }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{n+1}},\stackrel{}{}_\sigma \widehat{\rho }_i=\frac{1}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}(𝐈\widehat{\rho }_i\widehat{\rho }_i),$$
(551)
we find that
$$[(\stackrel{}{\eta }_i\times \stackrel{}{}_{\eta _i}+\stackrel{}{\kappa }_i\times \stackrel{}{}_{\kappa _i})\stackrel{}{A}_{Si}]\stackrel{}{\pi }_{Sj}=\stackrel{}{\sigma }\times (\stackrel{}{}_\sigma A^k{}_{Si}{}^{})\pi _{Sj}^k+\stackrel{}{A}_{Si}\times \stackrel{}{\pi }_{Sj},$$
(552)
while
$$A^k{}_{Sj}{}^{}[(\stackrel{}{\eta }_i\times \stackrel{}{}_{\eta _i}+\stackrel{}{\kappa }_i\times \stackrel{}{}_{\kappa _i})\pi _{Si}^k=A^k{}_{Sj}{}^{}\stackrel{}{\sigma }\times \stackrel{}{}_\sigma \pi ^k{}_{Si}{}^{}+\stackrel{}{A}_{Sj}\times \stackrel{}{\pi }_{Si},$$
(553)
thus verifying Eq.(546).
Thus as expected, in an instant form of dynamics, the internal angular momentum does not depend upon the interaction. On the other hand, the interaction-dependent internal boosts have the form of Eq.(205). Using the above forms for the final dynamical variables we obtain
$`\stackrel{}{𝒦}_{(int)}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\stackrel{~}{\stackrel{}{\eta }_i}[\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}+`$ (556)
$`+{\displaystyle \frac{\stackrel{~}{\stackrel{}{\kappa }_i}_{ji}Q_iQ_j[\stackrel{}{}_{\stackrel{~}{\eta }_i}\frac{1}{2}\stackrel{~}{𝒦}_{ij}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})2\stackrel{}{A}{}_{Sj}{}^{}(\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})]}{2\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}}}]`$
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{}}Q_iQ_j\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }_i}^2}\stackrel{}{}_{\stackrel{~}{\kappa }_i}\stackrel{~}{𝒦}_{ij}(\stackrel{~}{\stackrel{}{\kappa }_i},\stackrel{~}{\stackrel{}{\kappa }_j},\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j})+`$
$`+`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{}}{\displaystyle \frac{Q_iQ_j}{8\pi }}{\displaystyle \frac{\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}}{|\stackrel{~}{\stackrel{}{\eta }_i}\stackrel{~}{\stackrel{}{\eta }_j}|}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{}}{\displaystyle \frac{Q_iQ_j}{4\pi }}{\displaystyle d^3\sigma \frac{\stackrel{ˇ}{\pi }_{Sj}^r(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_j},\stackrel{~}{\stackrel{}{\kappa }_j})}{|\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_i}|}}`$ (559)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{ji}{}}Q_iQ_j{\displaystyle }d^3\sigma \stackrel{}{\sigma }[\stackrel{}{\pi }_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_i},\stackrel{~}{\stackrel{}{\kappa }_i})\stackrel{}{\pi }_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_j},\stackrel{~}{\stackrel{}{\kappa }_j})+`$
$`+\stackrel{}{B}{}_{Si}{}^{}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_i},\stackrel{~}{\stackrel{}{\kappa }_i})\stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }_j},\stackrel{~}{\stackrel{}{\kappa }_j})]=`$
$`=`$ $`𝒫_{(int)}^\tau \stackrel{}{R}_+.`$ (560)
In the last line the internal Møller center of energy (236) is shown explicitly. This equation allows us to express the internal canonical center of mass $`\stackrel{}{q}_+\stackrel{}{R}_+`$ defined in Eq.(236) in terms of the final canonical variables. The natural gauge fixing to the rest frame constraints $`\stackrel{}{𝒫}_{(int)}=_{i=1}^N\stackrel{~}{\stackrel{}{\kappa }_i}0`$ is $`\stackrel{}{q}_+0.`$
From Eq.(196) we get the following expression for the conserved energy-momentum tensor:
$`T^{\tau \tau }(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s))\sqrt{m_i^2+[\stackrel{}{\kappa }_i(T_s)Q_i{\displaystyle \underset{ji}{}}Q_j\stackrel{}{A}_{Sj}(T_s,\stackrel{}{\eta }_i(T_s))]^2}+`$ (561)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{ij}{\overset{1..N}{}}}Q_iQ_j[(\stackrel{}{\pi }_{Si}+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(\tau )))(\stackrel{}{\pi }_{Sj}+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_j(\tau )))+`$ (562)
$`+`$ $`\stackrel{}{B}_{Si}\stackrel{}{B}_{Sj}](T_s,\stackrel{}{\sigma }),`$ (563)
$`T^{r\tau }(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s))[\kappa _i^r(T_s)Q_i{\displaystyle \underset{ji}{}}Q_jA_{Sj}^r(T_s,\stackrel{}{\eta }_i(T_s))]+`$ (564)
$`+`$ $`Q_i{\displaystyle \underset{ij}{\overset{1..N}{}}}Q_j[(\stackrel{}{\pi }_{Si}+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s)))\times \stackrel{}{B}_{Sj}+(ij)](T_s,\stackrel{}{\sigma }),`$ (565)
$`T^{rs}(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s))`$ (567)
$`{\displaystyle \frac{[\kappa _i^r(T_s)Q_i_{ji}Q_jA_{Sj}^r(T_s,\stackrel{}{\eta }_i(T_s))][\kappa _i^s(T_s)Q_i_{ji}Q_jA_{Sj}^s(T_s,\stackrel{}{\eta }_i(T_s))]}{\sqrt{m_i^2+[\stackrel{}{\kappa }_i(T_s)Q_i_{ji}Q_j\stackrel{}{A}_{Sj}(T_s,\stackrel{}{\eta }_i(T_s))]^2}}}`$
$``$ $`{\displaystyle \underset{ij}{\overset{1..N}{}}}Q_iQ_j[{\displaystyle \frac{1}{2}}\delta ^{rs}[(\stackrel{}{\pi }_{Si}+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s)))(\stackrel{}{\pi }_{Sj}+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_j(T_s)))+`$ (568)
$`+`$ $`\stackrel{}{B}_{Si}\stackrel{}{B}_{Sj}]`$ (569)
$``$ $`[(\pi _{Si}^r+{\displaystyle \frac{^r}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_i(T_s)))(\pi _{Sj}^s+{\displaystyle \frac{^s}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{}{\eta }_j(T_s)))+`$ (570)
$`+`$ $`B_{Si}^rB_{Sj}^s]](T_s,\stackrel{}{\sigma }).`$ (571)
By using Eqs.(349) we get \[$`\stackrel{~}{\stackrel{}{\eta }}_i=\stackrel{~}{\stackrel{}{\eta }}_i(T_s)`$, $`\stackrel{~}{\stackrel{}{\kappa }}_i=\stackrel{~}{\stackrel{}{\kappa }}_i(T_s)`$; $`\stackrel{}{A}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)`$ is given in Eq.(311); $`\stackrel{~}{𝒦}_{ij}=\stackrel{~}{𝒦}_{ij}(\stackrel{~}{\stackrel{}{\kappa }}_i,\stackrel{~}{\stackrel{}{\kappa }}_j,\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j)`$ is given in Eqs.(320), (323)\]
$`T^{\tau \tau }(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_iQ_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij})`$ (573)
$`\sqrt{m_i^2+\left[\stackrel{~}{\stackrel{}{\kappa }}_iQ_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j(\stackrel{}{A}_{Sj}(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij})\right]^2}+`$
$`+`$ $`{\displaystyle \frac{1}{2}}Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}[(\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)`$ (575)
$`(\stackrel{}{E}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)+`$
$`+`$ $`\stackrel{}{B}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)\stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)]=`$ (576)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}[\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)`$ (578)
$`\left(1{\displaystyle \frac{Q_i_{ji}Q_j\stackrel{~}{\stackrel{}{\kappa }}_i[\stackrel{}{A}_{Sj}(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+\frac{1}{2}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij}]}{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}\right)+`$
$`+`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{}_\sigma \delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)Q_i{\displaystyle \underset{ji}{}}Q_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij}]+`$ (579)
$`+`$ $`{\displaystyle \frac{1}{2}}Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}[(\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)`$ (581)
$`(\stackrel{}{E}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)+`$
$`+`$ $`\stackrel{}{B}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)\stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)],`$ (582)
$`T^{r\tau }(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_iQ_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij})`$ (585)
$`\left[\stackrel{~}{\kappa }_i^rQ_i{\displaystyle \underset{ji}{}}Q_j\left(A_{Sj}^r(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^r\stackrel{~}{𝒦}_{ij}\right)\right]+`$
$`+`$ $`Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j[(\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i))\times `$ (587)
$`\times \stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+(ij)]=`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}[\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)`$ (589)
$`[\stackrel{~}{\kappa }_i^rQ_i{\displaystyle \underset{ji}{}}Q_j\left(A_{Sj}^r(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^r\stackrel{~}{𝒦}_{ij}\right)]+`$
$`+`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{}_\sigma \delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)\stackrel{~}{\kappa }_i^rQ_i{\displaystyle \underset{ji}{}}Q_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij}]+`$ (590)
$`+`$ $`Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j[(\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i))\times `$ (592)
$`\times \stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+(ij)],`$
$`T^{rs}(T_s,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i+{\displaystyle \frac{1}{2}}{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_iQ_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij})`$ (597)
$`{\displaystyle \frac{1}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}}\left(1+{\displaystyle \frac{Q_i_{ji}^{1..N}Q_j\stackrel{~}{\stackrel{}{\kappa }}_i[\stackrel{}{A}_{Sj}(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+\frac{1}{2}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij}]}{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}\right)`$
$`\left[\stackrel{~}{\kappa }_i^rQ_i{\displaystyle \underset{ji}{}}Q_j\left(A_{Sj}^r(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^r\stackrel{~}{𝒦}_{ij}\right)\right]`$
$`\left[\stackrel{~}{\kappa }_i^sQ_i{\displaystyle \underset{ji}{}}Q_j\left(A_{Sj}^s(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^s\stackrel{~}{𝒦}_{ij}\right)\right]`$
$``$ $`Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j[{\displaystyle \frac{1}{2}}\delta ^{rs}([\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)]`$ (599)
$`[\stackrel{}{E}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)]+`$
$`+`$ $`\stackrel{}{B}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)\stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j))`$ (600)
$``$ $`[E_{Si}^r(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{^r}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)][E_{Sj}^s(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{^s}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)]`$ (601)
$``$ $`B_{Si}^r(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)B_{Sj}^s(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)]=`$ (602)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{1}{\sqrt{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}}[\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\kappa }}_i)`$ (604)
$`(\stackrel{~}{\kappa }_i^r\stackrel{~}{\kappa }_i^s[1+{\displaystyle \frac{Q_i_{ji}^{1..N}Q_j\stackrel{~}{\stackrel{}{\kappa }}_i[\stackrel{}{A}_{Sj}(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+\frac{1}{2}\stackrel{}{}_{\stackrel{~}{\eta }_i}\stackrel{~}{𝒦}_{ij}]}{m_i^2+\stackrel{~}{\stackrel{}{\kappa }}_i^2}}]`$
$``$ $`Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j[\stackrel{~}{\kappa }_i^r(A_{Sj}^s(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^s\stackrel{~}{𝒦}_{ij})+`$ (605)
$`+`$ $`\stackrel{~}{\kappa }_i^s(A_{Sj}^r(\stackrel{~}{\stackrel{}{\eta }}_i\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{1}{2}}_{\stackrel{~}{\eta }_i}^r\stackrel{~}{𝒦}_{ij})])+`$ (606)
$`+`$ $`{\displaystyle \frac{1}{2}}\stackrel{}{}_\sigma \delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)\stackrel{~}{\kappa }_i^r\stackrel{~}{\kappa }_i^sQ_i{\displaystyle \underset{ji}{}}Q_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij}]`$ (607)
$``$ $`Q_i{\displaystyle \underset{ji}{\overset{1..N}{}}}Q_j[{\displaystyle \frac{1}{2}}\delta ^{rs}([\stackrel{}{E}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)]`$ (609)
$`[\stackrel{}{E}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{\stackrel{}{}}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)]+`$
$`+`$ $`\stackrel{}{B}_{Si}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)\stackrel{}{B}_{Sj}(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j))`$ (610)
$``$ $`[E_{Si}^r(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)+{\displaystyle \frac{^r}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i)][E_{Sj}^s(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)+{\displaystyle \frac{^s}{\mathrm{}}}\delta ^3(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j)]`$ (611)
$``$ $`B_{Si}^r(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_i,\stackrel{~}{\stackrel{}{\kappa }}_i)B_{Sj}^s(\stackrel{}{\sigma }\stackrel{~}{\stackrel{}{\eta }}_j,\stackrel{~}{\stackrel{}{\kappa }}_j)].`$ (612)
In the final canonical variables there is a dipole term \[gradient of delta function\] for each particle like it happens with spinning particles: the role of the spin is taken by $`\frac{1}{2}Q_i_{ji}^{1..N}Q_j\stackrel{}{}_{\stackrel{~}{\kappa }_j}\stackrel{~}{𝒦}_{ij}`$.
Let us remark that following we can define Dixon’s multipolar expansion of the energy momentum tensor about an arbitrary point on the Wigner hyperplane. There, the requirement of having the mass dipole vanishing could be shown to be equivalent to the identification of the point with $`\stackrel{}{R}_+\stackrel{}{q}_+`$, namely the internal center of mass.
For a cluster of $`n`$ particles inside the isolated $`N`$body system, we can now define a non-conserved energy-momentum tensor $`T_{(n)}^{AB}(T_s,\stackrel{}{\sigma })`$ by collecting in the previous equation all the terms depending on the canonical coordinates $`\stackrel{~}{\stackrel{}{\eta }_i},`$ $`\stackrel{~}{\stackrel{}{\kappa }_i}`$ of the $`n`$ particles of the cluster. This cluster energy-momentum tensor depends also on the canonical coordinates of the other $`Nn`$ particles: for the cluster these are “external fields”. If we make a Dixon multipole expansion of this cluster energy-momentum tensor with respect to an arbitrary point on the Wigner hyperplane and we require that the cluster mass dipole vanish, then we can identify a Møller center of energy $`\stackrel{}{R}_+^{(n)}`$ for the cluster. This is the only collective configuration variable which can be defined for a non-isolated cluster, which has no internal conserved Poincaré algebra associated with it, besides the nonconserved cluster 3-momentum $`\stackrel{}{𝒫}_{(int)}^{(n)}=_{i\epsilon \{n\}}\stackrel{~}{\stackrel{}{\kappa }_i}`$. Moreover, while $`T_{(n)}^{\tau \tau }(T_s,\stackrel{}{\sigma })`$ is the (non-conserved) energy density of the cluster, by analogy with the theory of dissipative fluids we can say that: i) $`q^\mu (T_s,\stackrel{}{\sigma })=ϵ_r^\mu (u(p_s))T_{(n)}^{r\tau }(T_s,\stackrel{}{\sigma })`$ is the heat flow; ii) $`𝒫(T_s,\stackrel{}{\sigma })=\frac{1}{3}_rT_{(n)}^{rr}(T_s,\stackrel{}{\sigma })`$ is the pressure; iii) $`T_{(n)(an)}^{rs}(T_s,\stackrel{}{\sigma })=T_{(n)}^{rs}(T_s,\stackrel{}{\sigma })\frac{1}{3}_uT_{(n)}^{uu}(T_s,\stackrel{}{\sigma })`$ is the shear (or anisotropic) stress tensor.
## VII Conclusions.
In this paper we analyzed how to extract the action-at-a-distance interparticle potential hidden in the semiclassical Lienard-Wiechert solution of the electromagnetic field equations, a subset of the solutions of the equations of motion for the isolated system formed by $`N`$ scalar charged particles plus the electromagnetic field. The problem is formulated in the Wigner-covariant rest-frame instant form of dynamics, which is defined on the Wigner hyperplanes orthognal to the total time-like four-momentum of the isolated system and which requires the choice of the sign of the energy of the particles (in this paper we considered only positive energies).
This was possible due to the semiclassical approximation of using Grassmann-valued electric charges ($`Q_i^2=0`$, $`Q_iQ_j0`$ for $`ij`$) as an alternative to the extended electron models used for the regularization of the Coulomb self energies. How this happens was shown in Ref., where the Coulomb potential was extracted from the electromagnetic potential by making the canonical reduction of the electromagnetic gauge freedom via the Shanmugadhasan canonical transformation. This is equivalent to the use of a Wigner-covariant radiation (or Coulomb) gauge in the rest-frame instant form.
Ref. presented the retarded Lienard-Wiechert solution for the transverse electromagnetic field in the rest frame instant radiation gauge: in this gauge, due to the transversality, the retarded Lienard-Wiechert potential associated with each charged particle depends on the whole past history of the other particles. At the semiclassical level a single charged particle with Grassmann-valued electric charge does not radiate even if it has a non-trivial Lienard-Wiechert potential, avoiding therefore the acausal features of the Abraham-Lorentz-Dirac equations, and has no mass renormalization. However, a system of $`N`$ charged particles produces, by virtue of the interference terms from the various retarded Lienard-Wiechert potentials of the particles, a radiation which reproduces the standard Larmor expression for radiation in the wave zone, when the particles are considered as external sources of the electromagnetic field and their equations of motion are not used.
If instead the particles are considered dynamical, the use of their equations of motion and of the semiclassical approximation lead to a drastic simplification of the Lienard-Wiechert potentials and fileds. Indeed, if we make an equal time expansion of the delay by expressing these potentials and fields in terms of particle coordinates, velocities and accelerations of every order, it turns out that all the accelerations decouple at the semiclassical level due to the particle equations of motion.
Therefore, at the semiclassical level the retarded, advanced and symmetric Lienard-Wiechert potentials and the electric and magnetic fields coincide and depend only on the positions and velocities of the particles, so that we can find their phase space expression in terms of particle positions and momenta.
In this way the semiclassical Lienard-Wiechert potential and fields can be reinterpreted as scalar and vector interparticle instantaneous action-at-a-distance potentials. It is then possible to identify a semiclassical reduced phase space containing only particles by eliminating the electromagnetic field by adding by hand second class contraints which force the transverse potential and electric field canonical variables to coincide with the semiclassical Lienard-Wiechert ones in the absence of incoming radiation: $`\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma })\stackrel{}{A}_{LW}(\tau ,\stackrel{}{\sigma })0`$, $`\stackrel{}{\pi }_{}(\tau ,\stackrel{}{\sigma })\stackrel{}{\pi }_{LW}(\tau ,\stackrel{}{\sigma })0`$. Let us remark that this could be done also in presence of an arbitrary incoming radiation $`\stackrel{}{A}_{(rad)}(\tau ,\stackrel{}{\sigma })`$, $`\stackrel{}{\pi }_{(rad)}(\tau ,\stackrel{}{\sigma })=\frac{}{\tau }\stackrel{}{A}_{(rad)}(\tau ,\stackrel{}{\sigma })`$ \[it is an arbitrary solution of the homogeneous wave equation and must not be interpreted as a pair of canonical variables\] by modifying the constraints to the form $`\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma })\stackrel{}{A}_{LW}(\tau ,\stackrel{}{\sigma })\stackrel{}{A}_{(rad)}(\tau ,\stackrel{}{\sigma })0`$, $`\stackrel{}{\pi }_{}(\tau ,\stackrel{}{\sigma })\stackrel{}{\pi }_{LW}(\tau ,\stackrel{}{\sigma })\stackrel{}{\pi }_{(rad)}(\tau ,\stackrel{}{\sigma })0`$.
The reduced phase space is obtained by means of the introduction of the Dirac brackets associated with these second class constraints. Since the old particle positions and momenta are no longer canonical in this reduced phase space, we had to find the new (Darboux) basis of particle canonical variables. The generators of the “internal” Poincaire‘ group inside the Wigner hyperplanes in the rest-frame instant form of dynamics can be reexpressed in terms of these new variables: the 3-momentum $`\stackrel{}{𝒫}_{(int)}`$ and the angular momentum $`\stackrel{}{𝒥}_{(int)}`$ become equal to those for $`N`$ free scalar particles (as expected in an instant form). The interaction dependent boosts $`\stackrel{}{𝒦}_{(int)}`$ are proportional to the “internal” canonical center of mass $`\stackrel{}{q}_+`$ inside the Wigner hyperplane: $`\stackrel{}{q}_+0`$ are the gauge-fixings to be be added to the rest-frame conditions $`\stackrel{}{𝒫}_{(int)}0`$, if one wishes to re-express the dynamics only in terms of particle “internal” relative variables. Also the energy-momentum tensor has been evalulated in the new canonical variables and there is a suggestion on how to find the Møller center of energy of a cluster of $`n`$ particles contained in the N particle isolated system.
The Hamiltonian in the rest frame frame instant form, generating the evolution in the rest-frame time of the decoupled “external” canonical center of mass, is the “internal” energy generator $`M=𝒫_{(int)}^\tau `$ (the invariant mass of the isolated $`N`$ particle system). The semiclassical Lienard-Wiechert solution implies the existence of interparticle action-at-a-distance potentials of two types: vector potentials minimally coupled to the Wigner spin 1 particle three-momentum under the square root associated with the kinetic energies; ii) a scalar potential (including the Coulomb potential) outside the square roots. In the semiclassical approximation all these potentials can be replaced by a unique scalar potential, which is the sum of the Coulomb potential and of a generalized Darwin one for arbitrary $`N`$. It is the (semiclassical) static and non-static complete potential corresponding to the one photon exchange tree Feynman diagrams of scalar electrodynamics and is a completely new result. The expression we find contains no $`N`$-body forces, being simply a sum of two particle interactions. This is a consequence of our use of Grassmann charges.
In the $`N=2`$ case we obtain a closed form of the solution by evaluating it in the rest frame after the gauge fixing $`\stackrel{}{q}_+0`$: the lowest order in $`1/c^2`$ contribution of the generalized Darwin potential agrees with the expression of the standard Darwin potential. We then show that in a semiclassical sense a special solution of the Hamilton equations is the Schild solution in which the two particles move in concentric circular orbits. We evaluate the frequency for equal masses
Future work will proceed along three parallel courses.
The first is the extention of this work to include semiclassical spinning particles. This will not only build upon the work here but also that of Ref.. Of particular concern there will be the issue of reproducing the correct spin-orbit and Darwin terms of the appropriate order (the two-body extentions of such terms in the one-body Dirac equation).
The second line is that of the quantization of our general Hamiltonian. Let us remark that quantization of the closed form Hamiltonian for $`N=2`$ in configuration space would involve not only the usual nonlocal operators for the kinetic energy but also non-local operators for the Darwin portions of the potential. For free scalar particles the positive energy wave equation $`i_\tau \varphi (\tau ,\stackrel{}{\sigma })=\sqrt{m^2+\mathrm{}}\varphi (\tau ,\stackrel{}{\sigma })`$ has been studied in Ref. by using pseudo-differential operators. Instead see Refs. for the difficulties in quantizing the equal mass two-body problem ($`H=2\sqrt{m^2+\stackrel{}{\kappa }^2}+\frac{\alpha }{\eta }`$) with only the Coulomb potential outside the square root. Note that the so-called spinless Salpeter equation would correspond to the quantization of our Hamiltonian with just the Coulomb interaction and at most the lowest order Darwin interaction. For the quantization when there are only scalar potentials inside the square root ($`\sqrt{m^2+V(|\stackrel{}{\eta }|)+\stackrel{}{\kappa }^2}`$) see Ref..
In our semiclassical approximation we have two options for the action-at-a-distance potentials: i) the vector ones $`\stackrel{}{V}_i`$ under the square roots and a scalar one $`U`$ outside; ii) a unique effective scalar potential sum of the Coulomb and Darwin ($`V_{DAR}=V_{LOD}+V_{HOD}`$) ones outside the square roots. It can be expected that the results of the quantization of the two options would produce inequivalent theories.
The third line is the development of the quantization of scalar electrodynamics on the Wigner hyperplanes in the rest-frame instant form. This would be a special Wigner-covariant instance of Tomonaga-Schwinger quantum field theory with a well defined covariant concept of “equal times”. To introduce a particle concept in such a quantum formulation, one will have to define Tomonaga-Schwinger asymptotic states and a reduction formalism. The natural candidates for the $`N`$-particle wave functions in such asymptotic states in the case of the Klein-Gordon field would be the wave functions corresponding to the quantization of $`N`$ positive energy scalar particles on the Wigner hyperplanes in the rest-frame instant form.
Moreover, in the rest-frame quantum field theory there will be new covariant “equal time” Green functions. This should allow the definition of a relativistic Schrödinger equation for bound states (replacing the Bethe-Salpeter equation for $`N`$=2 and avoiding by construction its problems with the spurious solutions, which are a byproduct of the use of asymptotic Fock states, since in a tensor product one cannot eliminate the possibility that an “in” particle be in the absolute future of another one). In the rest-frame quantum field theory it should also be possible to include bound states among the asymptotic Tomonaga-Schwinger states: they should be described by the quantization of isolated $`N`$ particle systems like the one studied in this paper.
ACKNOWLEDGMENT: L.Lusanna wishes to thank Prof. K.Kuchar for helpful discussions about the need of compatibility with the equations of motion of the second class constraints to be added by hand to select a symplectic subspace of solutions in the space of solutions of a given isolated system.
## A Exact Summation of the Transverse Vector Potential Series
In this Appendix we perform explicitly the summation of the vector potential below \[see Eq.(306} with $`\dot{\stackrel{}{\eta }}_i=\stackrel{}{\beta }_i`$ replaced by $`\stackrel{}{\kappa }_i/\sqrt{m_i^2+\stackrel{}{\kappa }_i^2}`$\]
$`\stackrel{}{A}_{}(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{1}{(2m)!}}{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m1}`$ (A2)
$`{\displaystyle \frac{1}{(2m+2)!}}\stackrel{}{}_\sigma ({\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m+1}]:=`$
$`:`$ $`=\stackrel{}{A}_1(\stackrel{}{\sigma },\tau )+\stackrel{}{A}_2(\stackrel{}{\sigma },\tau ).`$ (A3)
Using the result of Appendix C that $`(\stackrel{}{\kappa }_i\stackrel{}{}_\eta )^{2m}|\stackrel{}{\eta }|^{2m1}=[(2m1)!!]^2\frac{1}{|\stackrel{}{\eta }|}[\stackrel{}{\kappa }_i^2(\stackrel{}{\kappa }_i\frac{\stackrel{}{\eta }}{|\stackrel{}{\eta }|})^2]^m`$, we get \[$`\stackrel{}{}_\eta =/\stackrel{}{\eta }`$\]
$$\stackrel{}{A}_1(\tau ,\stackrel{}{\sigma })=\underset{i=1}{\overset{N}{}}\frac{Q_i}{4\pi }\underset{m=0}{\overset{\mathrm{}}{}}\left[\frac{[(2m1)!!]^2}{(2m)!|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}\frac{\stackrel{}{\kappa }_i}{(m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2})^{m+1/2}}(\stackrel{}{\kappa }_i^2(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2)^m\right].$$
(A4)
By using
$`{\displaystyle \frac{[(2m1)!!]^2}{(2m)!}}`$ $`=`$ $`{\displaystyle \frac{(2m)!}{(m!)^22^{2m}}}={\displaystyle \frac{(m1/2)!}{\sqrt{\pi }m!}}={\displaystyle \frac{\sqrt{\pi }()^{m1}(m1/2)}{(1/2m)!m!}}=`$ (A5)
$`=`$ $`{\displaystyle \frac{\sqrt{\pi }()^m}{(1/2m)!m!}}=()^m\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{m}}\right)`$ (A6)
we find that
$`\stackrel{}{A}_1(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}[{\displaystyle \frac{\stackrel{}{\kappa }_i}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}()^m\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{m}}\right)\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^m]=`$ (A7)
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{\kappa }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}{\displaystyle \frac{1}{\sqrt{m_i^2+(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}}}.`$ (A8)
For $`\stackrel{}{A}_2(\tau ,\stackrel{}{\sigma })`$ we need an expression for $`(\stackrel{}{\kappa }_i\stackrel{}{}_\eta )^{2m+1}|\stackrel{}{\eta }|^{2m+1}`$. One can show by an induction procedure that
$$(\stackrel{}{\kappa }_i\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_i|^{2m+1}=[(2m+1)!!]^2\underset{l=0}{\overset{m}{}}()^l\left(\genfrac{}{}{0pt}{}{m}{l}\right)(\stackrel{}{\kappa }^2)^{ml}\frac{(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^{2l+1}}{2l+1},$$
(A9)
and hence
$`\stackrel{}{A}_2(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{}_\sigma }{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{[(2m+1)!!]^2}{(2m+2)!}}\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^m\times `$ (A11)
$`{\displaystyle \underset{l=0}{\overset{m}{}}}()^l\left({\displaystyle \genfrac{}{}{0pt}{}{m}{l}}\right)(\stackrel{}{\kappa }_i^2)^l{\displaystyle \frac{(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^{2l+1}}{2l+1}}.`$
Now
$$\frac{[(2m+1)!!]^2}{(2m+2)!}=()^{m+1}\left(\genfrac{}{}{0pt}{}{1/2}{m+1}\right)$$
(A12)
so that
$`\stackrel{}{A}_2(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{}_\sigma }{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}()^{m+1}\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{m+1}}\right)\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^m\times `$ (A14)
$`{\displaystyle _0^{(\stackrel{}{\kappa }_i\frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})}}{\displaystyle \underset{l=0}{\overset{m}{}}}()^l\left({\displaystyle \genfrac{}{}{0pt}{}{m}{l}}\right)\left({\displaystyle \frac{w^2}{\stackrel{}{\kappa }_i^2}}\right)^ldw=`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{}_\sigma }{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}()^{m+1}\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{m+1}}\right)\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^m\times `$ (A16)
$`{\displaystyle _0^{(\stackrel{}{\kappa }_i\frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})}}\left(1{\displaystyle \frac{w^2}{\stackrel{}{\kappa }_i^2}}\right)^m𝑑w=`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{}_\sigma }{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}{\displaystyle _0^{(\stackrel{}{\kappa }_i\frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})}}\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2w^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^1\times `$ (A18)
$`{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}()^{m+1}\left({\displaystyle \genfrac{}{}{0pt}{}{1/2}{m+1}}\right)\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2w^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^{m+1}dw=`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}\stackrel{}{}_\sigma {\displaystyle _0^{(\stackrel{}{\kappa }_i\frac{(\stackrel{}{\sigma }\stackrel{}{\eta }_i)}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})}}\left({\displaystyle \frac{\stackrel{}{\kappa }_i^2w^2}{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}\right)^1\left[{\displaystyle \frac{1}{\sqrt{m_i^2+w^2}}}{\displaystyle \frac{1}{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}}\right]𝑑w.`$ (A19)
Now
$$\stackrel{}{}_\sigma \left(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}\right)=\frac{\stackrel{}{\kappa }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}\left(𝐈\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}\right).$$
(A20)
So we obtain
$`\stackrel{}{A}_2(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{Q_i}{4\pi }}{\displaystyle \frac{\stackrel{}{\kappa }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}(𝐈{\displaystyle \frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}}{\displaystyle \frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|}})\times `$ (A22)
$`({\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}{\sqrt{m_i^2+(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}}}1){\displaystyle \frac{\sqrt{m_i^2+\stackrel{}{\kappa }_{i}^{}{}_{}{}^{2}}}{\stackrel{}{\kappa }_i^2(\stackrel{}{\kappa }_i\frac{\stackrel{}{\sigma }\stackrel{}{\eta }_i}{|\stackrel{}{\sigma }\stackrel{}{\eta }_i|})^2}},`$
which when combined with the expression for $`\stackrel{}{A}_1(\tau ,\stackrel{}{\sigma })`$ yields the results given by Eq.(311).
## B Computation of Field Energy and Momentum Integrals
Here we carry out the details in the computation of the field energy and momentum for the case $`N=2`$. The general $`N`$ results obtained in the text are an immediate generalization. From Eq.(359) and Eq.(362) we find that
$`\stackrel{}{E}_S^2(\tau ,\sigma )`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2\times `$ (B8)
$`[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}][(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}]`$
$`{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n+2)!(2m)!}}\times `$
$`[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}][(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n+1}]`$
$`{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m+2)!(2n)!}}\times `$
$`[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+1})|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}][(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2})|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}]+`$
$`+{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m+2)!(2n+2)!}}\times `$
$`[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+2})|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n+1}][\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2})|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}],`$
$`\stackrel{}{B}_S^2(\tau ,\stackrel{}{\sigma })`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2`$ (B12)
$`[(\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}][(\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n})|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}]`$
$`{\displaystyle \frac{Q_1Q_2}{8\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}`$
$`[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m})|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}][(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n})|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}],`$
$`(\stackrel{}{E}_S(\tau ,\stackrel{}{\sigma })\times \stackrel{}{B}_S(\tau ,\stackrel{}{\sigma }))_k=`$ (B13)
$`=`$ $`{\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2\times `$ (B22)
$`[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}][\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}]`$
$`{\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}\dot{\stackrel{}{\eta }}_2\times `$
$`[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}][(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}]`$
$`{\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m+2)!}}\times `$
$`\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}][(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}]+`$
$`+{\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m+2)!}}\dot{\stackrel{}{\eta }}_2\times `$
$`[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}][\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}+`$
$`+(12).`$
Our aim here is to compute
$`{\displaystyle \frac{1}{2}}{\displaystyle d^3\sigma (\stackrel{}{E}_S^2+\stackrel{}{B}_S^2)(\tau ,\stackrel{}{\sigma })}`$ $`:`$ $`={\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}[(\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2)I_{1mn}`$ (B25)
$`{\displaystyle \frac{1}{(2n+1)(2n+2)}}I_{2mn}{\displaystyle \frac{1}{(2m+1)(2m+2)}}I_{3mn}+`$
$`+{\displaystyle \frac{1}{(2n+1)(2n+2)(2m+1)(2m+2)}}I_{4mn}+\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2I_{5mn}I_{6mn}],`$
$`{\displaystyle d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })}`$ $`:`$ $`={\displaystyle \frac{Q_1Q_2}{16\pi ^2}}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2n)!(2m)!}}[(\dot{\stackrel{}{\eta }}_1\dot{\stackrel{}{\eta }}_2\stackrel{}{I}_{7mn}\dot{\stackrel{}{\eta }}_2I_{8mn}`$ (B29)
$`{\displaystyle \frac{1}{(2m+1)(2m+2)}}\stackrel{}{I}_{9mn}+{\displaystyle \frac{\dot{\stackrel{}{\eta }}_2}{(2m+1)(2m+2)}}I_{10mn}]+`$
$`+(12),`$
which involve ten different integrals defined by
$`I_{1mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]\left[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]},`$ (B30)
$`I_{2mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]\left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n+1}\right]},`$ (B31)
$`I_{3mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]\left[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}\right]},`$ (B32)
$`I_{4mn}`$ $`=`$ $`{\displaystyle }d^3(_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}][\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n+1}],`$ (B33)
$`I_{5mn}`$ $`=`$ $`{\displaystyle d^3\left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]\left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]},`$ (B34)
$`I_{6mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]\left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]},`$ (B35)
$`\stackrel{}{I}_{7mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]\left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]},`$ (B36)
$`I_{8mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]\left[(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+1}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m1}\right]},`$ (B37)
$`\stackrel{}{I}_{9mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]\left[(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}\right]},`$ (B38)
$`I_{10mn}`$ $`=`$ $`{\displaystyle d^3\sigma \left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\sigma )^{2m+2}|\stackrel{}{\sigma }\stackrel{}{\eta }_1|^{2m+1}\right]\left[\stackrel{}{}_\sigma (\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\sigma )^{2n}|\stackrel{}{\sigma }\stackrel{}{\eta }_2|^{2n1}\right]}.`$ (B39)
The powers of $`\sigma `$ in each of the ten integrands is $`\sigma ^{2m12m1}\sigma ^{2n12n1}\sigma ^4`$ Therefore the integrals converge. Thus
$$I_i=\genfrac{}{}{0pt}{}{lim}{\mathrm{\Lambda }\mathrm{}}_0^\mathrm{\Lambda }\sigma ^2𝑑\sigma 𝑑\widehat{\mathrm{\Omega }}_\sigma ():=\genfrac{}{}{0pt}{}{lim}{\mathrm{\Lambda }\mathrm{}}_\mathrm{\Lambda }d^3\sigma ()=\genfrac{}{}{0pt}{}{lim}{\mathrm{\Lambda }\mathrm{}}I_i(\mathrm{\Lambda }).$$
(B40)
Since $`I_i(\mathrm{\Lambda })`$ is finite for finite $`\mathrm{\Lambda }`$, we can bring out the derivatives. Now perform integrations by parts, change to $`/\eta `$ from $`/\sigma `$, translate, and then use the fact that the integrals are finite so that they can be replaced by $`_\mathrm{\Lambda }`$. This gives ($`\stackrel{}{\eta }:=\stackrel{}{\eta }_{12}=\stackrel{}{\eta }_1\stackrel{}{\eta }_2`$)
$`I_{1mn}`$ $`=`$ $`(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2_\eta )^{2n+1}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n1},`$ (B41)
$`I_{2mn}`$ $`=`$ $`(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n+1},`$ (B42)
$`I_{3mn}`$ $`=`$ $`(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m+1}\sigma ^{2n1},`$ (B43)
$`I_{4mn}`$ $`=`$ $`\stackrel{}{}_\eta ^2(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m+1}\sigma ^{2n+1},`$ (B44)
$`I_{5mn}`$ $`=`$ $`\stackrel{}{}_\eta ^2(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n1},`$ (B45)
$`I_{6mn}`$ $`=`$ $`(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+1}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n1},`$ (B46)
$`\stackrel{}{I}_{7mn}`$ $`=`$ $`\stackrel{}{}_\eta (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n1},`$ (B47)
$`I_{8mn}`$ $`=`$ $`(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}\sigma ^{2n1},`$ (B48)
$`\stackrel{}{I}_{9mn}`$ $`=`$ $`\stackrel{}{}_\eta (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+1}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }+(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m+1}\sigma ^{2n1},`$ (B49)
$`I_{10}mn`$ $`=`$ $`\stackrel{}{}_\eta ^2(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}{\displaystyle _\mathrm{\Lambda }}d^3\sigma |\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m+1}\sigma ^{2n1}.`$ (B50)
The integrals that remain to be evaluated are each of the form
$`{\displaystyle _\mathrm{\Lambda }}d^3\sigma \sigma ^{2n1}|\stackrel{}{\sigma }(\stackrel{}{\eta }_1\stackrel{}{\eta }_2)|^{2m1}=2\pi {\displaystyle _0^\mathrm{\Lambda }}𝑑\sigma \sigma ^{2n+1}{\displaystyle _1^1}𝑑z(\sigma ^2+\eta ^22\eta \sigma z)^{m1/2}=`$ (B51)
$`=`$ $`{\displaystyle \frac{2\pi }{\eta (2m+1)}}{\displaystyle _0^\mathrm{\Lambda }}d\sigma \sigma ^{2n}(`$ (B52)
$`=`$ $`{\displaystyle \frac{2\pi \eta ^{2n+2m+1}}{(2m+1)}}{\displaystyle \underset{k=0}{\overset{2m+1}{}}}\left[{\displaystyle \frac{(\mathrm{\Lambda }/\eta )^{2n+k+1}}{2n+k+1}}({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{k}})[()^{1k}1]2{\displaystyle \frac{()^{1k}}{2n+k+1}}({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{k}})\right].`$ (B53)
It is of interest to show how the $`\mathrm{\Lambda }`$ dependent terms vanish after the $`\eta `$ derivatives have acted. First note that only the even $`k`$ terms of the $`\mathrm{\Lambda }`$ dependent sum survive. Of the ones that remain, the power of $`\eta `$ is $`\eta ^{2n+2m+12n2k1}=\eta ^{2m2k}=\eta ^{evenpower}`$. But since the power is even and since the number of $`\eta `$ derivatives always exceeds that even power, the derivative vanishes,. e.g.
$`{\displaystyle \frac{^3}{\eta _i\eta _j\eta _k}}\eta ^2={\displaystyle \frac{^2}{\eta _i\eta _j}}2\eta _k=0,`$ (B54)
$`{\displaystyle \frac{^5}{\eta _i\eta _j\eta _k\eta _l\eta _m}}\eta ^4=4{\displaystyle \frac{^4}{\eta _i\eta _j\eta _k\eta _l}}\eta ^2\eta _m=4{\displaystyle \frac{^3}{\eta _i\eta _j\eta _k}}(2\eta _l\eta _m+\eta ^2\delta _{ml})=0.`$ (B55)
Note further that the last sum is
$`2{\displaystyle \underset{k=0}{\overset{2m+1}{}}}()^k({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{k}}){\displaystyle _0^1}𝑑xx^{2n+k}=2{\displaystyle _0^1}𝑑xx^{2n}(1x)^{2m+1}=2B(2n+1,2m+2).`$ (B56)
Thus the portion of the integral that “survives” is
$$\mathrm{`}\mathrm{`}_\mathrm{\Lambda }d^3\sigma \sigma ^{2n1}|\stackrel{}{\sigma }\stackrel{}{\eta }|^{2m1}\mathrm{"}=\frac{4\pi (2n)!(2m)!}{(2n+2m+2)!}\eta ^{2n+2m+1},$$
(B58)
and our ten integrals appear as (using $`\stackrel{}{}^2\eta ^l=l\left(l1\right)\eta ^{l2}`$)
$`{\displaystyle \frac{I_{1mn}}{16\pi ^2(2n)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+2)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+1},`$ (B62)
$`{\displaystyle \frac{I_{2mn}}{16\pi ^2(2n+2)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+4)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3},`$
$`{\displaystyle \frac{I_{3mn}}{16\pi ^2(2n)!(2m+2)!}}={\displaystyle \frac{1}{4\pi (2n+2m+4)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3},`$
$`{\displaystyle \frac{I_{4mn}}{16\pi ^2(2n+2)!(2m+2)!}}={\displaystyle \frac{1}{4\pi (2n+2m+6)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+2}\stackrel{}{}_\eta ^2\eta ^{2n+2m+5}=`$
$`=`$ $`{\displaystyle \frac{I_{3mn}}{16\pi ^2(2n)!(2m+2)!}},`$ (B69)
$`{\displaystyle \frac{I_{5mn}}{16\pi ^2(2n)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+2)!}}\stackrel{}{}_\eta ^2(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}\eta ^{2n+2m+1},`$
$`{\displaystyle \frac{I_{6mn}}{16\pi ^2(2n)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+2)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+1},`$
$`{\displaystyle \frac{\stackrel{}{I}_{7mn}}{16\pi ^2(2n)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+2)!}}\stackrel{}{}_\eta (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+1}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}\eta ^{2n+2m+1},`$
$`{\displaystyle \frac{I_{8mn}}{16\pi ^2(2n)!(2m)!}}={\displaystyle \frac{1}{4\pi (2n+2m+2)!}}(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}\eta ^{2n+2m+1},`$
$`{\displaystyle \frac{\stackrel{}{I}_{9mn}}{16\pi ^2(2n)!(2m+2)!}}={\displaystyle \frac{1}{4\pi (2n+2m+4)!}}\stackrel{}{}_\eta (\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+3},`$
$`{\displaystyle \frac{I_{10mn}}{16\pi ^2(2n)!(2m+2)!}}={\displaystyle \frac{1}{4\pi (2n+2m+4)!}}\stackrel{}{}_\eta ^2(\dot{\stackrel{}{\eta }}_1\stackrel{}{}_\eta )^{2m+2}(\dot{\stackrel{}{\eta }}_2\stackrel{}{}_\eta )^{2n}\eta ^{2n+2m+3}=`$
$`=`$ $`{\displaystyle \frac{I_{8mn}}{16\pi ^2(2n)!(2m)!}}.`$ (B70)
Substitution into Eq.(LABEL:b4) and Eq.(B29) leads to Eq.(368) and Eq.(372)
## C The Evaluation of $`(\stackrel{}{\kappa }_1\stackrel{}{})^a(\stackrel{}{\kappa }_2\stackrel{}{})^b\eta ^{a+b1}`$
Consider the case with $`a=2m+1,b=2n+1`$,
$$(\stackrel{}{\kappa }_1\stackrel{}{})^{2m+1}(\stackrel{}{\kappa }_2\stackrel{}{})^{2n+1}\eta ^{2m+2n+1}.$$
(C1)
Let
$$\stackrel{}{\kappa }_1=\kappa _1\widehat{\kappa }_1;\stackrel{}{\kappa }_2=\kappa _2\widehat{\kappa }_2;\widehat{\kappa }_1^2=1=\widehat{\kappa }_2^2,$$
(C2)
and
$`\widehat{\kappa }_2`$ $`=`$ $`\alpha \widehat{\kappa }_1+\beta \widehat{\kappa _1}_{};\widehat{\kappa }_1\widehat{\kappa _1}_{}=0;\widehat{\kappa }_1\widehat{\kappa }_2=\alpha ;`$ (C3)
$`\alpha ^2+\beta ^2`$ $`=`$ $`1;\beta =(1(\widehat{\kappa }_1\widehat{\kappa }_2)^2)^{\frac{1}{2}},`$ (C4)
so
$`Eq.(C1)=\kappa _1^{2m+1}\kappa _2^{2n+1}{\displaystyle \underset{j=0}{\overset{2n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{j}}\right)\alpha ^j\beta ^{2n+1j}(\widehat{\kappa }_1\stackrel{}{})^{2m+1+j}(\widehat{\kappa _1}_{}\stackrel{}{})^{2n+1j}\eta ^{2m+n+1}.`$ (C5)
Orient our axes so that
$$\stackrel{}{\eta }=x\widehat{ı}+y\widehat{ȷ}+Z\widehat{k},$$
(C7)
with
$$Z=\stackrel{}{\eta }(\widehat{\kappa }_1\times \widehat{\kappa }_2),$$
(C8)
so that, the $`Z`$ direction is perpendicular to the plane containing $`\stackrel{}{\kappa }_1`$ and $`\stackrel{}{\kappa }_2`$. Further orient axes so that
$$\stackrel{}{\eta }=x\widehat{\kappa }_1+y\widehat{\kappa }_2+Z\widehat{\kappa }_1\times \widehat{\kappa }_2.$$
(C9)
Thus
$`Eq.(C4)=\kappa _1^{2m+1}\kappa _2^{2m+1}{\displaystyle \underset{j=0}{\overset{2n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{j}}\right)\alpha ^j\beta ^{2n+1j}(_x)^{2m+1+j}(_y)^{2n+1j}\eta ^{2m+n+1}.`$ (C10)
Consider just the portion
$$_x^{2m+1+j}_y^{2n+1j}\eta ^{2m+2n+1}.$$
(C12)
Let
$$\eta =(\rho ^2+Z^2)^{\frac{1}{2}};\rho ^2=x^2+y^2=zz^{};z=x+iy,z^{}=xiy.$$
(C13)
Thus
$$_x=\frac{}{z}+\frac{}{z^{}}=\frac{z}{x}\frac{}{z}+\frac{z^{}}{x}\frac{}{z^{}};_y=i(_z_z^{}),$$
(C14)
and
$`Eq.(C9)`$ $`=`$ $`i^{2n+1j}{\displaystyle \underset{h=0}{\overset{2m+1+j}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1+j}{h}}\right){\displaystyle \underset{k=0}{\overset{2n+1+j}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1j}{k}}\right)`$ (C16)
$`\times _z^{h+k}()^{2n+1yk}_z^{}^{2m+2n+2hk}(zz^{}+Z^2)^{m+n+\frac{1}{2}}.`$
Since $`_zz^{}=0=_z^{}z`$, the derivatives become relatively simple
$$_z^{h+k}(zz^{}+Z^2)^{m+m+\frac{1}{2}}=\frac{(m+n+\frac{1}{2})!}{(m+n+\frac{1}{2}hk)!}z^{h+k}(zz^{}+Z^2)^{m+n+\frac{1}{2}hk},$$
(C17)
and
$`_z^{}^{2m+2n+2hk}z^{h+k}(zz^{}+Z^2)^{m+n+\frac{1}{2}hk}=`$ (C18)
$`=`$ $`{\displaystyle \underset{i=0}{\overset{2m+2n+2hk}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+2n+2hk}{i}}\right)_z^{}^iz^{(h+k)}_z^{}^{2m+2n+2hki}(zz^{}+z^2)^{l+\frac{1}{2}hk}.`$ (C19)
Now
$$_z^{}^iz^{h+h}=\frac{(h+k)!}{(h+ki)!}z^{h+ki},$$
(C21)
Note that the factorial in denominator takes care of cutoff, so we can avoid having to worry about upper limit on sum.
$`_z^{}^{2(m+n+1)hki}(zz^{}+Z^2)^{l+\frac{1}{2}hk)}=`$ (C22)
$`=`$ $`{\displaystyle \frac{(m+n+\frac{1}{2}hk)!}{(imn\frac{3}{2})!}}z^{2(m+n+1)hki}(zz^{}+Z^2)^{imn\frac{3}{2}}.`$ (C23)
Thus combining factors
$`Eq.(C8)=\kappa _1^{2m+1}\kappa _2^{2n+1}(m+n+{\displaystyle \frac{1}{2}})!\eta ^{2m2n3}z^{2m+2n+1}(i)\beta ^{2n+1}()^{n+1}\times `$ (C24)
$`{\displaystyle \underset{j=0}{\overset{2n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{j}}\right)({\displaystyle \frac{i\alpha }{\beta }})^j{\displaystyle \underset{h=0}{\overset{2(m+n+1)+j}{}}}{\displaystyle \underset{k=0}{\overset{2n+1j}{}}}{\displaystyle \underset{i=0}{\overset{(2(m+n+1hk}{}}}\left({\displaystyle \frac{\eta ^2}{zz^{}}}\right)^i\left({\displaystyle \frac{z^{}}{z}}\right)^{h+k}()^k\times `$ (C25)
$`\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1+j}{h}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1j}{k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+2n+2hk}{i}}\right){\displaystyle \frac{1}{(imn\frac{3}{2})}}{\displaystyle \frac{(n+k)!}{(h+ki)!}}.`$ (C26)
Consider second two lines and let upper limits be $`\mathrm{}`$ and let the factorials set the limits. That becomes of form
$$\underset{j=0}{\overset{\mathrm{}}{}}\underset{i=0}{\overset{\mathrm{}}{}}\underset{h=0}{\overset{\mathrm{}}{}}\underset{k=0}{\overset{\mathrm{}}{}}f(j,h,k,h+k,i).$$
(C28)
Let $`q=2n+1j,p=h+k.`$ Then
$`Eq.(C18)`$ $`=`$ $`(2n+1)!\left({\displaystyle \frac{i\alpha }{\beta }}\right)^{2n+1}{\displaystyle \underset{q=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{\rho }{}}}\left({\displaystyle \frac{i\beta }{\alpha }}\right)^q\left({\displaystyle \frac{n^2}{zz^{}}}\right)^i\left({\displaystyle \frac{z^{}}{z}}\right)^p()^k`$ (C31)
$`\times {\displaystyle \frac{1}{q!}}{\displaystyle \frac{1}{(2n+1q)!}}{\displaystyle \frac{q!}{k!(qk)!}}{\displaystyle \frac{(2(n+m+1)q)!(2(m+m+1)p)!}{(pk)!(2(m+m+1)qp+k)!}}`$
$`\times {\displaystyle \frac{p!}{i!(2(m+n+1)pi)!}}{\displaystyle \frac{i}{(pi)!}}{\displaystyle \frac{1}{(imn\frac{3}{2})}}.`$
Use factorials to place upper bound on summation limits. Consider $`q`$ first; $`q0,qk,q2(m+n+1)p+k;q2n+1.`$ Let $`s=qk0`$. Upper bound appears ambiguous. One can make the replacement
$$\frac{(2(n+m+1)q)!}{(2n+1q)!}=\left(\frac{d}{dw}\right)^{2m+1}w^{2(n+m+1)q}_{w=1}=\left(\frac{d}{dw}^{2m+1}\right)w^{2(n+m+1)ks}_{w=1},$$
(C32)
and so our $`s`$ sum is restricted to
$$0s2(m+n+1)p.$$
(C33)
Note that the right hand side must be positive. So we can perform the $`q`$(or $`s`$) sum involved in (C31). It is (using $`1=\frac{(2(m+n+1)p)!}{(2(m+n+1)p)!}`$ )
$`{\displaystyle \underset{s=0}{\overset{2(m+n+1)p}{}}}w^{2(n+m+1)ks}\left({\displaystyle \frac{i\beta }{\alpha }}\right)^{k+s}{\displaystyle \frac{1}{s!}}{\displaystyle \frac{1}{(2(n+m+1)ps)!}}=`$ (C34)
$`=`$ $`w^{2(n+n+1)k}\left({\displaystyle \frac{i\beta }{\alpha }}\right)^k\left(1{\displaystyle \frac{i\beta }{\alpha w}}\right)^{2(m+n+1)p}{\displaystyle \frac{1}{(2(m+n+1)p)!}}=`$ (C35)
$`=`$ $`w^{pk}\left({\displaystyle \frac{i\beta }{\alpha }}\right)^k\left(w{\displaystyle \frac{i\beta }{\alpha }}\right)^{2(m+n+1)p}{\displaystyle \frac{1}{(2(m+n+1)p)!}}.`$ (C36)
So our expression (C31) reduces to (canceling the last factor)
$`Eq.(C19)=(2n+1)!\left({\displaystyle \frac{i\alpha }{\beta }}\right)^{2n+1}\left({\displaystyle \frac{d}{dw}}\right)^{2n+1}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{k=0}{\overset{p}{}}}w^{pk}\left({\displaystyle \frac{+i\beta }{\alpha }}\right)^k\left({\displaystyle \frac{wi\beta }{\alpha }}\right)^{2(m+m+1)p}`$ (C37)
$`\left({\displaystyle \frac{\eta ^2}{zz^{}}}\right)^2\left({\displaystyle \frac{z^{}}{z}}\right)^p{\displaystyle \frac{1}{k!}}{\displaystyle \frac{p!}{(pk)!}}{\displaystyle \frac{1}{i!}}{\displaystyle \frac{1}{(pi)!}}{\displaystyle \frac{1}{(2(m+n+1)pi)!}}{\displaystyle \frac{1}{(imn\frac{3}{2})!}}_{w=1}.`$ (C38)
We can perform the $`k`$ sum and note that the $`p`$ sum now has well defined limits since $`2(m+n+1)p10`$ and therefore $`ip2(m+n+1)i`$.
Thus
$`Eq.(C23)=(2n+1)!\left({\displaystyle \frac{i\alpha }{\beta }}\right)^{2n+1}\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}{\displaystyle \underset{i=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{p=i}{\overset{2(m+n+1)i}{}}}w^p\left(1+{\displaystyle \frac{i\beta }{\alpha \omega }}\right)^p`$ (C39)
$`\times \left(w{\displaystyle \frac{i\beta }{\alpha }}\right)^{2(m+n+1)}\left({\displaystyle \frac{\eta ^2}{zz^{}}}\right)^i\left({\displaystyle \frac{z^{}}{z}}\right)^p{\displaystyle \frac{1}{(2(m+n+1)pi)!}}{\displaystyle \frac{1}{i!}}{\displaystyle \frac{1}{(imn\frac{3}{2})!}}.`$ (C40)
Next, noting that $`w^p\left(1+\frac{i\beta }{\alpha w}\right)^p=(w+\frac{i\beta }{\alpha })^p`$ and that the $`p`$ sum produces $`\frac{1}{(2(m+n+1)2i)!}`$ we see that $`i`$ is restricted to $`0im+n+1`$.
The $`p`$ sum is best performed by changing variables to $`r=pi`$. Then $`0r2(m+n+1i)`$. Thus we obtain
$`Eq.(C24)=(2n+1)!\left({\displaystyle \frac{i\alpha }{\beta }}\right)^{2n+1}\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}{\displaystyle \underset{i=0}{\overset{m+n+1}{}}}{\displaystyle \frac{1}{i!}}{\displaystyle \frac{1}{(2(m+n+1i)!}}{\displaystyle \frac{1}{(imn\frac{3}{2})!}}\times `$ (C41)
$`\left({\displaystyle \frac{w+i\beta /\alpha )}{(wi\beta /\alpha )}}\right)^i\left({\displaystyle \frac{\eta ^2}{zz^{}}}\right)^i\left({\displaystyle \frac{z^{}}{z}}\right)^i\left[1+{\displaystyle \frac{z^{}(w+\frac{i\beta }{\alpha })}{z(wi\frac{\beta }{\alpha })}}\right]^{2(m+n+1i)}\left(w{\displaystyle \frac{i\beta }{\alpha }}\right)^{2(m+n+1)}.`$ (C42)
To perform the $`i`$ sum, notice that
$`{\displaystyle \frac{1}{(imn3/2)!}}{\displaystyle \frac{1}{(2(m+n+1i))!}}=`$ (C43)
$`=`$ $`{\displaystyle \frac{\pi ^{\frac{1}{2}}}{2^{(m+n+1i)}(m+n+1i)!(m+n+\frac{1}{2}i)!(1mn3/2)!}}=`$ (C44)
$`=`$ $`{\displaystyle \frac{\pi ^{\frac{1}{2}}}{2^{2(m+n+11)}}}{\displaystyle \frac{sin(\pi (imn\frac{1}{2})}{(m+n+1i)!}}={\displaystyle \frac{()^{imn+1}}{\pi ^{\frac{1}{2}}2^{(m+n+1i)}(m+n+1i)!}}.`$ (C45)
Further notice that the 2nd line of (C42) can be simplified to
$$\left(w^2+\beta ^2/\alpha ^2\right)^i(\eta ^2)^i\frac{[z(wi\beta /\alpha )+z^{}(w+\beta /\alpha )]^{2(m+n+1)1}}{z^{2(m+n+1)}}.$$
(C46)
So
$`Eq.(C25)`$ $`=`$ $`(2n+1)!({\displaystyle \frac{i\alpha }{\beta }})^{2n+1}\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}{\displaystyle \frac{()^{m+n+1}}{\pi ^{\frac{1}{2}}2^{2(m+n+1)}}}(z(w{\displaystyle \frac{1}{2}}\beta )+z^{}(w+{\displaystyle \frac{i\beta }{\alpha }}))^{(2m+n+1)}`$ (C48)
$`\times {\displaystyle \underset{i=0}{\overset{m+n+1}{}}}{\displaystyle \frac{1}{i!}}{\displaystyle \frac{2^{2i}}{(m+n+1i)!}}{\displaystyle \frac{\eta ^{2i}(w^2+\beta ^2/\alpha ^2)^i()^i}{[z(wi\beta /\alpha )+z^{}(w+i\beta /\alpha )]^{2i}}}.`$
The sum produces
$$\frac{1}{(m+n+1)!}\left[1\frac{4\eta ^2(w^2+\beta ^2/\alpha ^2)}{[z(wi\beta \alpha )+z^{}(w+i\beta /\alpha )]^2}\right]^{m+n+1}.$$
(C49)
Combine this with (C48) and then (LABEL:c17) gives
$`Eq.(C17)=\kappa _1^{2m+1}\kappa _2^{2n+1}{\displaystyle \frac{(m+n+\frac{1}{2})!}{\pi ^{\frac{1}{2}}2^{2(m+n+1)}}}\eta ^{2m2n3}()i^2\alpha ^{2n+1}(2n+1)!()^{m+n+1}\times `$ (C50)
$`\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}\left[{\displaystyle \frac{(z(wi\beta /\alpha )+z^{}(w+i\beta /\alpha ))^24\eta ^2(w^2+\beta ^2/\alpha ^2)}{(m+n+1)!}}\right]^{m+n+1}.`$ (C51)
Before expanding or simplifying consider case when
$`\stackrel{}{\kappa }_2=\stackrel{}{\kappa }_1=\stackrel{}{\kappa }\beta =0.`$
Then
$`Eq.(C30)`$ $`=`$ $`()^{m+n}{\displaystyle \frac{\kappa ^{2m+2n+2}}{\pi ^{\frac{1}{2}}2^{2(m+n+1)}}}(m+n+{\displaystyle \frac{1}{2}})!\eta ^{2m2n3}(2n+1)!\times `$ (C53)
$`\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}{\displaystyle \frac{(w^2(4x^2\eta ^2))^{m+n+1}}{(m+n+1)!}}.`$
Use
$$\frac{d^{2m+1}}{dw^{2m+1}}w^{2m+2n+1}_{w=1}=\frac{(2(m+n+1))!}{(2n+1)!}.$$
(C54)
Also $`(m+n+\frac{1}{2})!=\frac{(2(m+n+1))!}{(m+n+1)!}\frac{\pi ^{\frac{1}{2}}}{2^{2m+2n+1}}`$. Thus
$`Eq.(C30)`$ $`=`$ $`\kappa ^{2m+2n+2}{\displaystyle \frac{[(2(m+n+1))!]^2}{[(m+n+1)!]^22^{2m+2n+1}}}\left(1{\displaystyle \frac{x^2}{\eta ^2}}\right)^{m+n+1}{\displaystyle \frac{1}{\eta }}=`$ (C55)
$`=`$ $`[(2m+2n+1)!!]^2{\displaystyle \frac{1}{\eta }}(\kappa ^2(\kappa \widehat{\eta })^2)^{m+n+1}.`$ (C56)
Return to (C51) and consider $`z(wi\beta /\alpha )+z^{}(w+i\beta /\alpha )=2(xw+y\beta /\alpha ).`$ In (C51) this factor becomes
$`[]^{m+n+1}`$ $`=`$ $`2^{(m+n+1)}()^{m+n+1}[\eta ^2(w^2+\beta ^2/\alpha ^2)x^2w^22xyw\beta /\alpha y^2\beta ^2/\alpha ^2]^{m+n+1}=`$ (C57)
$`=`$ $`2^{2(m+n+1)}()^{m+n+1}[w^2(\eta ^2x^2)+\beta ^2/\alpha ^2(\eta ^2y^2)2xyw\beta /\alpha ]^{m+n+1}=`$ (C58)
$`=`$ $`2^{2(m+n+1)}()^{m+n+1}[\eta ^2(w^2+\beta ^2/\alpha ^2)(xw+y\beta /\alpha )^2]^{m+n+1}.`$ (C59)
Thus
$`Eq.(C30)`$ $`=`$ $`\kappa _1^{2m+1}\kappa _2^{2n+1}{\displaystyle \frac{[2(m+n+1)]!}{2^{2m+2n+2}[(n+m+1)!]^2}}{\displaystyle \frac{\alpha ^{2n+1}(2n+1)!}{\eta }}\times `$ (C61)
$`\left({\displaystyle \frac{d}{dw}}\right)^{2m+1}\left((w^2+\beta ^2/\alpha ^2)({\displaystyle \frac{xw}{\eta }}+{\displaystyle \frac{y}{\eta }}{\displaystyle \frac{\beta }{\alpha }})^2\right)^{m+n+1}|_{w=1}.`$
The derivative is of the form
$$\left(\frac{d}{dw}\right)^{2m+1}\left(aw^2+bw+c\right)^{m+n+1}=\left(\frac{d}{dw}\right)^{2m+1}F(Q(w))=\left(\frac{d}{dw}\right)^{2m+1}f(w).$$
(C62)
This radical is
$$\omega ^2\left(1\frac{x^2}{\eta ^2}\right)\frac{2xy\beta w}{\eta ^2\alpha }+\frac{\beta ^2}{d^2}\left(1\frac{y^2}{\eta ^2}\right),$$
(C63)
with roots
$$w=\frac{\frac{xy\beta }{\eta ^2\alpha }\pm \sqrt{\frac{x^2y^2\beta ^2}{\eta ^4\alpha ^2}\frac{\beta ^2}{\alpha ^2}\left(1\frac{y^2}{\eta ^2}\right)\left(1\frac{x^2}{\eta ^2}\right)}}{\left(1\frac{x^2}{\eta ^2}\right)}.$$
(C64)
Since the argument of the root is
$$\frac{\beta ^2}{\alpha ^2}\left(1\frac{x^2}{\eta ^2}\frac{z^2}{\eta ^2}\right)<0,$$
(C65)
the roots are complex. Thus
$`aw^2+bw+c=a(w\gamma i\delta )(w\gamma +i\delta ),`$
where
$$\gamma =\frac{xy\beta }{\alpha (\eta ^2x^2)};\delta =\frac{\beta }{\alpha }\frac{Z\eta }{(\eta ^2x^2)}.$$
(C66)
Consider the derivative:
$`{\displaystyle \frac{d^{2m+1}}{dw^{2m+1}}}(w\gamma i\delta )^{m+n+1}(w\gamma +i\delta )^{m+n+1}=`$ (C67)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{2m+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right)\left({\displaystyle \frac{d}{dw}}\right)^j(w\gamma i\delta )^{m+n+1}\left({\displaystyle \frac{d}{dw}}\right)^{2m+1j}(w\gamma +i\delta )^{m+n+1}=`$ (C68)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{2m+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right){\displaystyle \frac{[(m+n+1)!]^2}{(m+n+1j)!(n+jm)!}}`$ (C70)
$`(w\gamma i\delta )^{m+n+1j}(w\gamma +i\delta )^{n+jm}.`$
The factorial further restricts the sum to
$`{\displaystyle \underset{j=max(0,nm)}{\overset{min(2m+1,m+n+1)}{}}}.`$
Consider the case when $`\beta =0=\delta =\gamma .`$ Then as $`w=1`$
$`Eq.(C41)`$ $`=`$ $`{\displaystyle \underset{j=max(0,nm)}{\overset{min(2m+1,n+m+1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right){\displaystyle \frac{(m+n+1)!}{(m+n+1)j)!}}{\displaystyle \frac{(m+n+1)!}{(n+jm)!}}=`$ (C71)
$`=`$ $`{\displaystyle \frac{[(m+n+1)!]^2}{(2n+1)!}}{\displaystyle \underset{j=max(0,nm)}{\overset{min(2m+1,n+m+1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n+jm}}\right).`$ (C72)
In order to perform this sum, consider the related product (for $`s=1)`$
$`{\displaystyle \underset{j=0}{\overset{2m+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right)s^j{\displaystyle \underset{i=mn}{\overset{n+m+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n+jm}}\right)s^i=`$ (C73)
$`=`$ $`(\mathrm{let}l=imn)={\displaystyle \underset{j=0}{\overset{2m+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right)s^j{\displaystyle \underset{l=0}{\overset{2n+1}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{l}}\right)s^{\mathrm{}}s^{2n+m}=`$ (C74)
$`=`$ $`s^{(n+m)}(1+s)^{2(m+n+1)}={\displaystyle \underset{h=0}{\overset{2(m+n+1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2(m+n+1)}{h}}\right)s^hs^{n+m}.`$ (C75)
But consider the product of the two sums. Let $`h=j+\mathrm{}`$ so that $`0h2(m+n+1)`$. Let $`l=hj`$. Thus, since $`l0=>hj`$ and since $`2m+1j`$, we have $`jmin(2m+1,h)`$. Now we also have $`2n+1l`$ or $`jh2n1`$; thus $`jmax(0,h2n1)`$ so
$$\underset{h=0}{\overset{2(m+n+1)}{}}\left(\genfrac{}{}{0pt}{}{2(m+n+1)}{h}\right)s^h=\underset{h=0}{\overset{2(m+n+1)}{}}s^h\underset{j=max(0,h2n1)}{\overset{min(2m+1,h)}{}}\left(\genfrac{}{}{0pt}{}{2m+1}{j}\right)\left(\genfrac{}{}{0pt}{}{2n+1}{hj}\right).$$
(C76)
Consider the term $`h=m+n+1`$ in the sum. Then
$$\frac{(2(m+n+1)!}{[(m+n+1)!]^2}=\underset{j=max(0,mn)}{\overset{min(2m+1,m+n+1}{}}\left(\genfrac{}{}{0pt}{}{2m+1}{j}\right)\left(\genfrac{}{}{0pt}{}{2n+1}{n+jm}\right),$$
(C77)
where we have used $`\left(\genfrac{}{}{0pt}{}{2n+1}{hj}\right)=\left(\genfrac{}{}{0pt}{}{2n+1}{n+m+1j}\right)=\left(\genfrac{}{}{0pt}{}{2n+1}{n+jm}\right)`$. Thus
$`Eq.(C42)={\displaystyle \frac{(2(m+n+1))!}{(2n+1)!}}.`$
Note that
$`Eq.(C35)|_{\beta =0}`$ $`=`$ $`\kappa ^{2(m+n+1)}{\displaystyle \frac{(2(m+n+1)!)^2}{2^{2(m+n+1)}[(n+m+1)!]^2}}{\displaystyle \frac{1}{\eta }}\left(1{\displaystyle \frac{x^2}{\eta ^2}}\right)^{m+n+1}=`$ (C78)
$`=`$ $`([2m+2n+1)!!]^2{\displaystyle \frac{1}{\eta }}(\kappa ^2(\kappa \eta )^2)^{m+n+1},`$ (C79)
which checks again with Eq.(C33). Thus, with this we write again
$`Eq.(C35)`$ $`=`$ $`\kappa _1^{2m+1}\kappa _2^{2n+1}{\displaystyle \frac{(2(m+n+1))!}{2^{2m+2+2}[(n+m+1)!]^2}}{\displaystyle \frac{\alpha ^{2n+1}}{\eta }}(1{\displaystyle \frac{x^2}{\eta ^2}})^{m+n+1}\times `$ (C82)
$`{\displaystyle \underset{j=max(0,mn)}{\overset{min(2m+1,n+m+1)}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{j}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{n+jm}}\right)`$
$`(w\gamma i\delta )^{m+n+1j}(w\gamma +i\delta )^{n+jm}|_{w=1}.`$
To simplify, let $`j=kn`$. Also, note
$`(w\gamma i\delta )^{m+n+1j}(w\gamma +i\delta )^{n+jm}|_{w=1}=`$ (C83)
$`=`$ $`\left[{\displaystyle \frac{1}{\alpha (\eta ^2x^2)}}\right]^{2n+1}[\alpha (\eta ^2x^2)\beta (xy+i\eta Z)]^{m+2n+1k}`$ (C85)
$`[\alpha (\eta ^2x^2)\beta (xyi\eta z)]^{km}.`$
Thus Eq.(C82) becomes
$`{\displaystyle \frac{[(2m+2n+1)!!]^2}{(2(m+n+1))!}}{\displaystyle \frac{\kappa _1^{2m+1}\kappa _2^{2n+1}}{\eta }}{\displaystyle \frac{\left(1\frac{x^2}{\eta ^2}\right)}{\eta ^{4n+2}}}[(m+n+1))!]^2{\displaystyle }_{k=max(m,n)}^{min(2m+n+1,2n+m+1)}\times `$ (C86)
$`\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{kn}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{km}}\right)(\sigma ^2\alpha \beta (xy+i\eta Z))^{m+2n+1k}(\sigma ^2\alpha \beta (xyi\eta z))^{km}.`$ (C87)
where
$$\sigma ^2=\eta ^2x^2.$$
(C88)
If $`mn`$ then sum becomes $`(l:=km)`$
$$=\underset{l=0}{\overset{2m+1}{}}\left(\genfrac{}{}{0pt}{}{2n+1}{l}\right)\left(\genfrac{}{}{0pt}{}{2m+1}{\mathrm{}+m+n}\right)(\sigma ^2\alpha (xy+i\eta Z)\beta )^{2n+1l}(\sigma ^2\alpha (xyi\eta Z)\beta )^l.$$
(C89)
Because of the symmetry we have
$$=\underset{l=0}{\overset{\eta }{}}\left(\genfrac{}{}{0pt}{}{2n+1}{l}\right)\left(\genfrac{}{}{0pt}{}{2m+1}{\mathrm{}+mn}\right)[(\sigma ^2\alpha (xy+i\eta Z)\beta )^{2n+1l}(\sigma ^2\alpha (xyi\eta Z)\beta )^l+c.c.].$$
(C90)
The bracket $`[]`$ in (C90) above is of the form
$`[]`$ $`=`$ $`(uiv)^{2n+1l}(u+iv)^l+c.c.=(\mathrm{let}u+iv=re^{i\varphi })=`$ (C91)
$`=`$ $`r^{2n+1}(e^{i(2(ln)1)\varphi }+c.c.)=2\gamma ^{2n+1}\mathrm{cos}(2(nl)+1)\varphi ,`$ (C92)
in which
$`r`$ $`=`$ $`\sqrt{(\sigma ^2\alpha xy\beta )^2+\eta ^2Z^2\beta ^2},`$ (C93)
$`\varphi `$ $`=`$ $`tan^1{\displaystyle \frac{\beta \eta z}{\sigma ^2\alpha xy\beta }}.`$ (C94)
We further need
$`cos((2(nl)+1)\varphi )`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{nl}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2(nl)+1}{2k}}\right)()^k(sin\varphi )^{2k}(cos\varphi )^{2(nlk)+1}=`$ (C95)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{nl}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2(nl)+1)}{2k}}\right)()^k{\displaystyle \underset{j=0}{\overset{k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right)()^j(cos\varphi )^{2(nlk+j)+1},`$ (C96)
where
$$cos\varphi =\frac{\sigma ^2\alpha xy\beta }{r}.$$
(C97)
Recall the original variables
$`\eta `$ $`=`$ $`|\stackrel{}{\eta }|;\eta ^2=\rho ^2+Z^2=\sigma ^2+x^2=x^2+y^2+Z^2,`$ (C98)
$`\alpha `$ $`=`$ $`\widehat{\kappa }_1\widehat{\kappa }_2;\beta =\sqrt{1(\widehat{\kappa }_1\widehat{\kappa }_2)^2},`$ (C99)
$`x`$ $`=`$ $`\widehat{\kappa }_1\stackrel{}{\eta };\beta y=\widehat{\kappa }_1\stackrel{}{\eta }=(\widehat{\kappa }_2\alpha \widehat{\kappa }_1)\stackrel{}{\eta }.`$ (C100)
Thus
$`\sigma ^2\alpha xy\beta `$ $`=`$ $`\alpha (\eta ^2(\stackrel{}{\eta }\widehat{\kappa }_1)^2)\stackrel{}{\eta }\widehat{}\kappa _1(\widehat{\kappa }_2\alpha \widehat{\kappa }_1)\stackrel{}{\eta }=`$ (C101)
$`=`$ $`\eta ^2(\widehat{\kappa }_1\widehat{\kappa }_2\widehat{\kappa }_1\widehat{\eta }\widehat{\kappa }_2\widehat{\eta }),`$ (C102)
$`\eta ^2Z^2\beta ^2`$ $`=`$ $`\eta ^4(1(\widehat{\eta }\widehat{\kappa }_1)^2(\widehat{\eta }\widehat{\kappa }_2)^2(\widehat{\kappa }_1\widehat{\kappa }_2)^22\widehat{\kappa }_1\widehat{\kappa }_2\widehat{\kappa }\widehat{\eta }\widehat{\kappa }_2\widehat{\eta })),`$ (C103)
and
$$r^2=\eta ^4(1(\widehat{\kappa }_1\widehat{\eta })^2)(1(\widehat{\kappa }_2\widehat{\eta })^2).$$
(C104)
Thus
$$cos\varphi =\frac{\stackrel{}{\kappa }_1\stackrel{}{\kappa }_2\stackrel{}{\kappa }_1\widehat{\eta }\stackrel{}{\kappa }_2\stackrel{}{\eta }}{\sqrt{\kappa _1^2(\widehat{\kappa }_1\stackrel{}{\eta })^2}\sqrt{\stackrel{}{\kappa }_2^2(\kappa _2\widehat{\eta })^2}},$$
(C105)
and
$`{\displaystyle \frac{r^{2n+1}}{\eta ^{4n+2}}}\left(1{\displaystyle \frac{x^2}{\eta ^2}}\right)^{mn}\kappa _1^{2m+1}\kappa _2^{2n+1}=`$ (C106)
$`=`$ $`(1(\widehat{\kappa }_1\widehat{\eta })^2)^n(1(\widehat{\kappa }_2\widehat{\eta })^2)^n(\kappa _1^2(\kappa _1\widehat{\eta })^2)^{1/2}(\kappa _2^2(\widehat{\kappa }_2\widehat{\eta })^2)^{1/2},`$ (C107)
$`(\kappa _1^2(\kappa _1\widehat{\eta })^2)^{mn}\kappa _1^{2n}\kappa _2^{2n}=`$ (C108)
$`=`$ $`(\kappa _1^2(\kappa _1\stackrel{}{\eta })^2)^{\frac{1}{2}}(\kappa _2^2(\stackrel{}{\kappa }_2\widehat{\eta })^2)^{1/2}(\kappa _2^2(\stackrel{}{\kappa }_2\widehat{\eta })^2)^n(\kappa _1^2\stackrel{}{\kappa }_1\eta )^2)^m.`$ (C109)
So finally factoring out $`cos\varphi `$ we obtain
$`(\stackrel{}{\kappa }_1\stackrel{}{})^{2m+1}(\stackrel{}{\kappa }_2\stackrel{}{})^{2n+1}\eta ^{2m+2n+1}=`$ (C110)
$`=`$ $`2{\displaystyle \frac{[(2m+2n+1)!!]^2}{(2(m+n+1))!}}[(n+m+1))!]^2`$ (C114)
$`\times {\displaystyle \frac{(\stackrel{}{\kappa }_1\stackrel{}{\kappa }_2\stackrel{}{\kappa }_1\widehat{\eta }\stackrel{}{\kappa }_2\widehat{\eta })(\kappa _1^2(\stackrel{}{\kappa }_1\widehat{\eta }))^2)(\kappa _2^2\stackrel{}{\kappa }_2\widehat{\eta }))^2)^n}{\eta }}`$
$`\times {\displaystyle \underset{l=0}{\overset{n}{}}}{\displaystyle \underset{k=0}{\overset{nl}{}}}{\displaystyle \underset{j=0}{\overset{2k}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+1}{l}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2m+1}{l+mn}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2(n\mathrm{})+1}{2k}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{2k}{j}}\right)()^{k+j}(cos^2\varphi )^{nl+jk}.`$
Check once again the limit $`\stackrel{}{\kappa }_1=\stackrel{}{\kappa }_2`$ where Eq.(63) becomes
$$Eq.(C63)=2\frac{[(2m+2n+1)!!]^2}{(2(m+n+1))!}[(n+m+1)!]^2\frac{(\kappa ^2(\stackrel{}{\kappa }\widehat{\eta })^2)}{\eta }^{m+n+1},$$
(C115)
where the triple sum $``$ collapses to a single sum:
$`{\displaystyle (cos\varphi =1)}j=k=0,`$
so that
$$=\underset{l=0}{\overset{n}{}}\left(\genfrac{}{}{0pt}{}{2n+1}{l}\right)\left(\genfrac{}{}{0pt}{}{2m+1}{l+mn}\right)=\frac{1}{2}\frac{(2(m+n+1))!}{(n+m+1)!]^2}.$$
(C116)
## D Derivation of the invariant mass $`M`$ using the new Dirac brackets
Kerner has shown that it is possible to develope a single-time Hamiltonian formulation of Wheeler-Feynman dynamics. His idea, basically is to replace the infinitude of “field” coordinates by an infinity of mechnical ones and then with the high order of the equations of motion replaced by higher powers of the momentum in the interaction. In Ref., Crater and Yang give a modification of his approach to obtain a Hamiltonian expression for both scalar and vector interactions through order $`1/c^4`$. The approach taken in this Appendix is similar to that given in with two important distinctions: 1) terms of all order in $`1/c^2`$ are included and 2) the effects of the new Dirac brackets are included. The net result is an expression for $`M`$ that agrees exactly with the results obtained in Eq.(LABEL:VI14). This result would not be obtained without the use of these brackets in working our Hamilton’s equations.
We begin with the Lagrangian expression for the invariant mass \[see Eq.(229) with $`\stackrel{}{\lambda }(\tau )=0`$; $`h_1`$ is defined in Eq.(368)\]
$`E_{rel}`$ $`=`$ $`h(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })={\displaystyle \frac{m_1}{\sqrt{1\dot{\stackrel{}{\eta }_1}^2}}}+{\displaystyle \frac{m_2}{\sqrt{1\dot{\stackrel{}{\eta }_2}^2}}}+{\displaystyle \frac{Q_1Q_2}{4\pi |\stackrel{}{\eta }|}}+Q_1Q_2h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta }):=`$ (D1)
$`:`$ $`=h_0(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })+Q_1Q_2h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta }).`$ (D2)
In order to find the Hamiltonian $`H(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2;\stackrel{}{\eta })`$ from $`h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })`$ we must demand that Hamilton’s equation be satisfied . We use the Dirac bracket since we have used the constraint as a strong condition on the dynamical variables. Thus we begin with
$`\dot{\stackrel{}{\eta }_i}`$ $`=`$ $`\{\stackrel{}{\eta }_i,H\}^{}=\{\stackrel{}{\eta }_i,H\}`$ (D5)
$`[{\displaystyle }d^3\sigma \{\stackrel{}{\eta }_i,{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i)\}\{{\displaystyle \underset{j}{}}Q_j\stackrel{}{\mathrm{\Pi }}_{Sj}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j),H\}`$
$`\{\stackrel{}{\eta }_i,{\displaystyle \underset{j}{}}Q_j\stackrel{}{\mathrm{\Pi }}_{Sj}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_j,\stackrel{}{\kappa }_j)\}\{{\displaystyle \underset{i}{}}Q_i\stackrel{}{A}_{Si}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_i,\stackrel{}{\kappa }_i),H\}],`$
in which
$$H=H_0+Q_1Q_2H_1(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2;\stackrel{}{\eta }),$$
(D6)
and
$$H_0=\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}+\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}.$$
(D7)
Substituting this Hamiltonian into the above bracket and using Grassmann truncation yields
$`\dot{\stackrel{}{\eta }_1}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_1}}`$ (D10)
$`[{\displaystyle }d^3\sigma \{\stackrel{}{\eta }_1,Q_1\stackrel{}{A}_{S1}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_1,\stackrel{}{\kappa }_1)\}\{Q_2\stackrel{}{\mathrm{\Pi }}_{S2}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_2,\stackrel{}{\kappa }_2),\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}\}`$
$`\{\stackrel{}{\eta }_1,Q_1\stackrel{}{\mathrm{\Pi }}_{S1}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_1,\stackrel{}{\kappa }_1)\}\{Q_2\stackrel{}{A}_{S2}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_2,\stackrel{}{\kappa }_2),\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}\}=`$
$`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_1}}{\displaystyle \frac{}{\stackrel{}{\kappa }_1}}\left({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{}{\stackrel{}{\eta }_2}}\right)𝒦_{12},`$ (D11)
$`\dot{\stackrel{}{\eta }_2}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_2}}`$ (D14)
$`[{\displaystyle }d^3\sigma \{\stackrel{}{\eta }_2,Q_2\stackrel{}{A}_{S2}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_2,\stackrel{}{\kappa }_2)\}\{Q_1\stackrel{}{\mathrm{\Pi }}_{S1}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_1,\stackrel{}{\kappa }_1),\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}\}`$
$`\{\stackrel{}{\eta }_2,Q_2\stackrel{}{\mathrm{\Pi }}_{S2}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_2,\stackrel{}{\kappa }_2)\}\{Q_1\stackrel{}{A}_{S1}(\tau ,\stackrel{}{\sigma }\stackrel{}{\eta }_1,\stackrel{}{\kappa }_1),\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}\}=`$
$`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_2}}+{\displaystyle \frac{}{\stackrel{}{\kappa }_2}}\left({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{}{\stackrel{}{\eta }_2}}\right)𝒦_{12}.`$ (D15)
But
$$d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })=\frac{}{\stackrel{}{\eta }_2}𝒦_{12}=\frac{}{\stackrel{}{\eta }_1}𝒦_{12}:=\frac{}{\stackrel{}{\eta }}𝒦_{12}.$$
(D16)
Hence
$`\dot{\stackrel{}{\eta }_1}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_1}}{\displaystyle \frac{}{\stackrel{}{\kappa }_1}}\left({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })}\right),`$ (D17)
$`\dot{\stackrel{}{\eta }_2}`$ $`=`$ $`{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}+Q_1Q_2{\displaystyle \frac{H_1}{\stackrel{}{\kappa }_2}}{\displaystyle \frac{}{\stackrel{}{\kappa }_2}}\left({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })}\right).`$ (D18)
Substituting this into $`h_0`$ we find that
$`h_0`$ $`=`$ $`H_0+{\displaystyle \frac{Q_1Q_2}{4\pi \eta }}+`$ (D21)
$`+{\displaystyle \frac{(m_1^2+\stackrel{}{\kappa }_1^2)}{m_1^2}}\stackrel{}{\kappa }_1{\displaystyle \frac{}{\stackrel{}{\kappa }_1}}(Q_1Q_2H_1{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle }d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })+`$
$`+{\displaystyle \frac{(m_2^2+\stackrel{}{\kappa }_2^2)}{m_2^2}}\stackrel{}{\kappa }_2{\displaystyle \frac{}{\stackrel{}{\kappa }_2}}(Q_1Q_2H_1{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle }d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })=`$
$`=`$ $`HQ_1Q_2h_1=H_0+Q_1Q_2H_1Q_1Q_2h_1.`$ (D22)
Thus letting
$$H_1=\frac{1}{4\pi \eta }+\stackrel{~}{H}_1,$$
(D23)
and using \[see Eqs.(368) and (372) for the expressions of $`h_1`$ and $`\stackrel{}{h}_1`$\]
$`h_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })`$ $`=`$ $`h_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })+O(Q_1Q_2),`$ (D24)
$`\stackrel{}{h}_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta })`$ $`=`$ $`\stackrel{}{h}_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })+O(Q_1Q_2),`$ (D25)
where
$$d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma }):=Q_1Q_2\stackrel{}{h}_1(\dot{\stackrel{}{\eta }_1},\dot{\stackrel{}{\eta }_2},\stackrel{}{\eta }),$$
(D26)
then we obtain the following differential equation for $`\stackrel{~}{H}_1`$
$`\stackrel{~}{H}_1{\displaystyle \frac{\stackrel{~}{H}_1}{\stackrel{}{\kappa }_1}}\stackrel{}{\kappa }_1{\displaystyle \frac{(m_1^2+\stackrel{}{\kappa }_1^2)}{m_1^2}}{\displaystyle \frac{\stackrel{~}{H}_1}{\stackrel{}{\kappa }_2}}\stackrel{}{\kappa }_2{\displaystyle \frac{(m_2^2+\stackrel{}{\kappa }_2^2)}{m_2^2}}=`$ (D27)
$`=`$ $`h_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })`$ (D28)
$``$ $`{\displaystyle \frac{(m_1^2+\stackrel{}{\kappa }_1^2)}{m_1^2}}\stackrel{}{\kappa }_1{\displaystyle \frac{}{\stackrel{}{\kappa }_1}}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{h}_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta }))`$ (D29)
$``$ $`{\displaystyle \frac{(m_2^2+\stackrel{}{\kappa }_2^2)}{m_2^2}}\stackrel{}{\kappa }_2{\displaystyle \frac{}{\stackrel{}{\kappa }_2}}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{h}_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })+O(Q_1Q_2).`$ (D30)
The $`O(Q_1Q_2)`$ term gives a vanishing contribution as both sides are multiplied by $`Q_1Q_2`$. Using this and $`\stackrel{}{}_\eta ^2|\stackrel{}{\eta }|^l=l(l1)|\stackrel{}{\eta }|^{l2}`$ we obtain, in addition to the expression given for $`h_1`$ in Eq.(401),
$`{\displaystyle d^3\sigma (\stackrel{}{E}_S\times \stackrel{}{B}_S)(\tau ,\stackrel{}{\sigma })}=Q_1Q_2\stackrel{}{h}_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })=`$ (D31)
$`=`$ $`{\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(\stackrel{}{}_\eta [{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\times `$ (D33)
$`{\displaystyle \frac{\left((\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n}+(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}\right)\eta ^{2n+2m+1}}{(2n+2m+2)!}}`$
$``$ $`{\displaystyle \frac{\left((\frac{\stackrel{}{\kappa }_1\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}{}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}+(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\right)\eta ^{2n+2m+3}}{(2n+2m+4)!}}]]).`$ (D34)
The differential equation for $`\stackrel{~}{H}_1`$ is of the form
$$\stackrel{~}{H}_1𝒪\stackrel{~}{H}_1=h_1(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}},\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}},\stackrel{}{\eta })\stackrel{}{𝒪}\stackrel{}{h}_1(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}},\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}},\stackrel{}{\eta }),$$
(D36)
where we define the linear operator $`𝒪`$ by
$$𝒪=\frac{(m_1^2+\stackrel{}{\kappa }_1^2)}{m_1^2}\stackrel{}{\kappa }_1\frac{}{\stackrel{}{\kappa }_1}+\frac{(m_2^2+\stackrel{}{\kappa }_2^2)}{m_2^2}\stackrel{}{\kappa }_2\frac{}{\stackrel{}{\kappa }_2},$$
(D37)
and the linear operator $`\stackrel{}{𝒪}`$ by
$$\stackrel{}{𝒪}=\frac{(m_1^2+\stackrel{}{\kappa }_1^2)}{m_1^2}\stackrel{}{\kappa }_1\frac{}{\stackrel{}{\kappa }_1}\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}+\frac{(m_2^2+\stackrel{}{\kappa }_2^2)}{m_2^2}\stackrel{}{\kappa }_2\frac{}{\stackrel{}{\kappa }_2}\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}.$$
(D38)
In order to solve for $`\stackrel{~}{H}_1`$ we need first to work out the right hand side of its differential equation. First note
$$\stackrel{}{𝒪}\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}=(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}+\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}})\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}},$$
(D39)
while
$`\stackrel{}{𝒪}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2n+1}=`$ (D40)
$`=`$ $`[(2n+1){\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}+(2m+1){\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}]\times `$ (D42)
$`({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2n+1}.`$
So Eq.(372) imply
$`\stackrel{}{𝒪}Q_1Q_2\stackrel{}{h}_1({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}},{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}},\stackrel{}{\eta })={\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}([(2n+1)({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^2+`$ (D43)
$`+(2n+2m+4)({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )+(2m+1)({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^2]\times `$ (D44)
$`\times [{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n}}{(2n+2m+2)!}}\eta ^{2n+2m+1}`$ (D45)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}}{(2n+2m+4)!}}\eta ^{2n+2m+3}]).`$ (D46)
In analogy to our decomposition of the field energy integral we decompose this into single and double sum pieces giving
$`\stackrel{}{𝒪}Q_1Q_2\stackrel{}{h}_1={\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{(2m)!}}[{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\times `$ (D47)
$`(({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m}+({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m}))\eta ^{2m1}`$ (D48)
$`{\displaystyle \frac{1}{(2m+2)!}}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta ({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1}+`$ (D49)
$`+{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta ({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1})\eta ^{2m+1}]+`$ (D50)
$`+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}((2n+2m+4)[{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}`$ (D51)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}}{(2n+2m+2)!}}\eta ^{2n+2m+1}`$ (D52)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}}{(2n+2m+4)!}}\eta ^{2n+2m+3}]+`$ (D53)
$`+(2n+2m+6)[{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}}{(2n+2m+4)!}}\eta ^{2n+2m+3}`$ (D54)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+3}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+3}}{(2n+2m+6)!}}\eta ^{2n+2m+5}])\}.`$ (D55)
Combining like terms we obtain
$`h_1\stackrel{}{𝒪}\stackrel{}{h}_1={\displaystyle \frac{Q_1Q_2}{4\pi }}\{({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta }}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}){\displaystyle \frac{\eta }{2}})+`$ (D56)
$`+{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}((2n+2m+3)[{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\times `$ (D57)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}}{(2n+2m+2)!}}\eta ^{2n+2m+1}`$ (D58)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}}{(2n+2m+4)!}}\eta ^{2n+2m+3}]`$ (D59)
$`(2n+2m+5)[{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}}{(2n+2m+4)!}}\eta ^{2n+2m+3}`$ (D60)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+3}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+3}}{(2n+2m+6)!}}\eta ^{2n+2m+5}])\}.`$ (D61)
Based on the above expression we assume a particular solution of the form
$`Q_1Q_2H_1(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,\stackrel{}{\eta })=k{\displaystyle \frac{Q_1Q_2}{4\pi }}\times `$ (D62)
$`\left({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta }}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta ){\displaystyle \frac{\eta }{2}}\right)+`$ (D63)
$`+{\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[a_{mn}{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\times `$ (D64)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+1}}{(2n+2m+2)!}}`$ (D65)
$`b_{mn}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}+`$ (D66)
$`+c_{mn}{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}`$ (D67)
$`d_{mn}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+3}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+3}\eta ^{2n+2m+5}}{(2n+2m+6)!}}].`$ (D68)
Using
$$𝒪\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}=2\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}},$$
(D69)
and
$`𝒪({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}})^{2n+1}=`$ (D70)
$`=`$ $`(2n+2m+2)({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2m+1}({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )^{2n+1},`$ (D71)
the left hand side of the differential equation for $`\stackrel{~}{H}_1(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,\stackrel{}{\eta })`$ becomes
$`(1𝒪)Q_1Q_2\stackrel{~}{H}_1(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,\stackrel{}{\eta })=k{\displaystyle \frac{Q_1Q_2}{8\pi }}\times `$ (D73)
$`\left({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta }}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )\eta \right)=`$
$`=`$ $`{\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[a_{mn}(2n+2m+3){\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\times `$ (D79)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+1}}{(2n+2m+2)!}}+`$
$`+b_{mn}(2n+2m+3){\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}`$
$`c_{mn}(2n+2m+5){\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\times `$
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}+`$
$`+d_{mn}(2n+2m+5)({\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+3}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+3}\eta ^{2n+2m+5}}{(2n+2m+6)!}}].`$
Comparing the two sides of the equation leads to
$$k=1;a_{mn}=b_{mn}=c_{mn}=d_{mn}=1,$$
(D80)
and thus
$`Q_1Q_2\stackrel{~}{H}_1(\stackrel{}{\kappa }_1,\stackrel{}{\kappa }_2,\stackrel{}{\eta })={\displaystyle \frac{Q_1Q_2}{4\pi }}\times `$ (D81)
$`\left({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{1}{\eta }}({\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta )({\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}\stackrel{}{}_\eta ){\displaystyle \frac{\eta }{2}}\right)+`$ (D82)
$`+{\displaystyle \frac{Q_1Q_2}{4\pi }}{\displaystyle \underset{m=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}[{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+1}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+1}\eta ^{2n+2m+1}}{(2n+2m+2)!}}`$ (D83)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}+`$ (D84)
$`+{\displaystyle \frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}}{\displaystyle \frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}}{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+2}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+2}\eta ^{2n+2m+3}}{(2n+2m+4)!}}`$ (D85)
$`{\displaystyle \frac{(\frac{\stackrel{}{\kappa }_1}{\sqrt{m_1^2+\stackrel{}{\kappa }_{1}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2m+3}(\frac{\stackrel{}{\kappa }_2}{\sqrt{m_2^2+\stackrel{}{\kappa }_{2}^{}{}_{}{}^{2}}}\stackrel{}{}_\eta )^{2n+3}\eta ^{2n+2m+5}}{(2n+2m+6)!}}],`$ (D86)
which agrees exactly with the Darwin portion of $`M`$ obtained earlier Eq.(LABEL:VI14).
## E Schild-Like Solution for the Two-Body Problem in the Case of Equal Masses
Here we present the semiclassical Hamilton equations for the equal mass case restricted to circular orbits (we suppress the tilde notation). First note that the Hamiltonian Eq.(LABEL:VI37)
$`H=2\sqrt{m^2+\stackrel{}{\kappa }^2}+{\displaystyle \frac{Q_1Q_2}{4\pi \eta }}+{\displaystyle \frac{Q_1Q_2}{8\pi \eta }}\times `$ (E1)
$`{\displaystyle \frac{[m^2(3\stackrel{}{\kappa }^2+(\stackrel{}{\kappa }\widehat{\eta })^2]2\stackrel{}{\kappa }^2[\stackrel{}{\kappa }^23(\stackrel{}{\kappa }\widehat{\eta })^2]\sqrt{\frac{m^2+\stackrel{}{\kappa }^2}{m^2+(\stackrel{}{\kappa }\widehat{\eta })^2}}2[\stackrel{}{\kappa }^2+(\stackrel{}{\kappa }\widehat{\eta })^2][m^2+(\stackrel{}{\kappa }\widehat{\eta })^2]}{(m^2+\stackrel{}{\kappa }^2)[m^2+(\stackrel{}{\kappa }\widehat{\eta })^2]}}.`$ (E2)
is of the form
$$H=2\sqrt{m^2+\stackrel{}{\kappa }^2}+\frac{Q_1Q_2}{4\pi \eta }+\frac{Q_1Q_2}{8\pi \eta }f(\stackrel{}{\kappa }^2,(\stackrel{}{\kappa }\widehat{\eta })^2).$$
(E3)
Thus, Hamilton’s equations are \[$`f_{,1}=f/\stackrel{}{\kappa }^2`$, $`f_{,2}=f/(\stackrel{}{\kappa }\widehat{\eta })^2`$\]
$`\dot{\stackrel{}{\kappa }}\stackrel{}{=}{\displaystyle \frac{H}{\stackrel{}{\eta }}}`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{4\pi \eta ^3}}\stackrel{}{\eta }{\displaystyle \frac{Q_1Q_2}{8\pi \eta ^3}}\stackrel{}{\eta }f(\stackrel{}{\kappa }^2,(\stackrel{}{\kappa }\widehat{\eta })^2)+{\displaystyle \frac{Q_1Q_2}{4\pi \eta }}f,_2(\stackrel{}{\kappa }^2,(\stackrel{}{\kappa }\widehat{\eta })^2)\stackrel{}{\kappa }(\stackrel{}{\kappa }\widehat{\eta }),`$ (E4)
$`\dot{\stackrel{}{\eta }}`$ $`\stackrel{}{=}{\displaystyle \frac{H}{\stackrel{}{\kappa }}}={\displaystyle \frac{2\stackrel{}{\kappa }}{\sqrt{\stackrel{}{\kappa }^2+m^2}}}+{\displaystyle \frac{Q_1Q_2}{4\pi \eta }}[f,_1(\stackrel{}{\kappa }^2,(\stackrel{}{\kappa }\widehat{\eta })^2)\stackrel{}{\kappa }+f,_2(\stackrel{}{\kappa }^2,(\stackrel{}{\kappa }\widehat{\eta })^2)\stackrel{}{\eta }(\stackrel{}{\kappa }\widehat{\eta })].`$ (E5)
From this we can see that circular orbits defined by $`\dot{\stackrel{}{\eta }}\stackrel{}{\eta }=0`$ are implied by $`\stackrel{}{\kappa }\stackrel{}{\eta }=0`$. This furthermore implies that $`\dot{\stackrel{}{\kappa }}\stackrel{}{\kappa }=0.`$ Thus not only is $`\stackrel{}{\eta }^2=const.`$ but also $`\stackrel{}{\kappa }^2=const.`$ Imposing these conditions on the above Hamilton equations we can simplify our equations above to
$`\dot{\stackrel{}{\kappa }}`$ $`\stackrel{}{=}`$ $`{\displaystyle \frac{H}{\stackrel{}{\eta }}}={\displaystyle \frac{Q_1Q_2}{4\pi \eta ^3}}\stackrel{}{\eta }{\displaystyle \frac{Q_1Q_2}{8\pi \eta ^3}}\stackrel{}{\eta }f(\stackrel{}{\kappa }^2,0)=B\stackrel{}{\eta },`$ (E6)
$`\dot{\stackrel{}{\eta }}`$ $`\stackrel{}{=}`$ $`{\displaystyle \frac{H}{\stackrel{}{\kappa }}}={\displaystyle \frac{2\stackrel{}{\kappa }}{\sqrt{\stackrel{}{\kappa }^2+m^2}}}+{\displaystyle \frac{Q_1Q_2}{4\pi \eta }}f,_1(\stackrel{}{\kappa }^2,0)\stackrel{}{\kappa }=A\stackrel{}{\kappa }.`$ (E7)
with $`B`$ and $`A`$ constants. Combine the two equations and we find
$$\dot{\stackrel{}{\kappa }}=\frac{\ddot{\stackrel{}{\eta }}}{A}=B\stackrel{}{\eta },$$
(E8)
so that
$$\ddot{\stackrel{}{\eta }}=AB\stackrel{}{\eta }:=\mathrm{\Omega }^2\stackrel{}{\eta },$$
(E9)
with
$`\mathrm{\Omega }^2`$ $`=`$ $`{\displaystyle \frac{Q_1Q_2}{4\pi \eta ^3\sqrt{\stackrel{}{\kappa }^2+m^2}}}[1+f(\stackrel{}{\kappa }^2,0)]=`$ (E10)
$`=`$ $`{\displaystyle \frac{Q_1Q_2}{4\pi \eta ^3\sqrt{\stackrel{}{\kappa }^2+m^2}}}[1\stackrel{}{\kappa }^2{\displaystyle \frac{2\stackrel{}{\kappa }^2\sqrt{\stackrel{}{\kappa }^2+m^2}m^3}{2m^3(m^2+\stackrel{}{\kappa }^2)}}].`$ (E11)
The frequency is defined by the initial data $`\eta `$ and $`|\stackrel{}{\kappa }|`$ and is real for $`Q_1Q_2<0`$ for $`0\stackrel{}{\kappa }^2\stackrel{}{\kappa }_{\mathrm{max}}^2`$ in which $`\stackrel{}{\kappa }_{\mathrm{max}}^2`$ is the value at which $`\mathrm{\Omega }^2=0`$. Let us remark that at the semiclassical level $`Q_i`$ are Grassmann variables: therefore $`\mathrm{\Omega }\sqrt{Q_1Q_2}`$ is to be interpreted as an even algebraic object satisfying $`\mathrm{\Omega }^4=0`$, $`Q_1\mathrm{\Omega }^2=Q_2\mathrm{\Omega }^2=0`$.
We also find
$$H=2\sqrt{\stackrel{}{\kappa }^2+m^2}[1\frac{1}{4}\mathrm{\Omega }^2\stackrel{}{\eta }^2].$$
(E12)
Our solution to Eq.(E9) is
$`\stackrel{}{\eta }(\tau )`$ $`=`$ $`\stackrel{}{\alpha }\mathrm{cos}\mathrm{\Omega }\tau +\stackrel{}{\beta }\mathrm{sin}\mathrm{\Omega }\tau =`$ (E13)
$`=`$ $`\stackrel{}{\alpha }(1{\displaystyle \frac{1}{2!}}\mathrm{\Omega }^2\tau ^2)+\stackrel{}{\beta }(\mathrm{\Omega }\tau {\displaystyle \frac{1}{3!}}\mathrm{\Omega }^3\tau ^3),`$ (E14)
and thus
$$\dot{\stackrel{}{\eta }}(\tau )=\stackrel{}{\alpha }\mathrm{\Omega }^2\tau +\stackrel{}{\beta }(\mathrm{\Omega }\frac{1}{2}\mathrm{\Omega }^3\tau ^2).$$
(E15)
But $`\dot{\stackrel{}{\eta }}\stackrel{}{\eta }=0.`$ This implies that
$`(\stackrel{}{\beta }^2\stackrel{}{\alpha }^2)\mathrm{\Omega }^2\tau +\stackrel{}{\alpha }\stackrel{}{\beta }\mathrm{\Omega }(12\mathrm{\Omega }^2\tau ^2)=0,`$
so that
$$\stackrel{}{\beta }^2=\stackrel{}{\alpha }^2,\stackrel{}{\alpha }\stackrel{}{\beta }=0,$$
(E16)
and therefore
$$\stackrel{}{\alpha }^2=\stackrel{}{\eta }^2.$$
(E17)
Otherwise the vectors $`\stackrel{}{\alpha }`$ and $`\stackrel{}{\beta }`$ are arbitrary. |
warning/0001/cond-mat0001159.html | ar5iv | text | # Equilibrium and adhesion of Nb/sapphire: the effect of oxygen partial pressure
## I Introduction
Oxide-metal interfaces continue to be studied intensively because of the many ways in which they are of commercial and scientific importance. Applications range from the nanoscale in microelectronics packaging to the macroscale engineering of thermal barrier coatings or the formation of protective scales. The science of these interfaces has been addressed in volumes of conference papers and reviews . There are also reviews in the literature which specifically address the theoretical questions about the nature of the bonding at these interfaces, such as: what determines the site-preference of metal atoms on the oxide surface; whether the bonding can be thought of as predominantly covalent or metallic and how to quantify these concepts; whether a simple classical image model can be used to interpret the bonding; what is the strength of adhesion of metal to oxide. The basis for answering these questions is to have reliable calculations of the electronic structure and total energy of particular surfaces and interfaces. Such calculations came of age over the past ten years or so with the use of first-principles methods. These mainly apply density functional theory (DFT) and the local density approximation (LDA) , which are the basis of the calculations we shall report here. Hartree-Fock calculations are also feasible and have been applied effectively to the Ag/MgO interface , although they tend to be more expensive than DFT for larger systems. Since the reviews cited, there have been numerous applications of DFT to study bonding in the initial stages of deposition of metal on oxide, with cluster or multilayer geometries, notably on MgO , but to a lesser extent on more complex oxides such as TiO<sub>2</sub> , MgAl<sub>2</sub>O<sub>4</sub> and $`\alpha `$Al<sub>2</sub>O<sub>3</sub> .
The bonding of Nb to $`\alpha `$Al<sub>2</sub>O<sub>3</sub> has long been a subject for experimental work, because besides its relevance to electronic components it offers practical advantages for sample preparation: the two materials bond strongly (anomalously strongly according to a recent study ), do not react chemically and have similar coefficients of thermal expansion. In the orientation Nb(111)/Al<sub>2</sub>O<sub>3</sub> (0001) there is a lattice mismatch of $`<2\%`$, allowing the preparation of a nearly coherent interface (using molecular beam epitaxy), the atomic structure of which has been studied by high resolution transmission electron microscopy (HRTEM) and analysed in detail . This interface was the subject of first-principles calculations which used periodic boundary conditions, making the reasonable assumption that the effect of misfit dislocation can be neglected . Our recent work analysed the nature of the bonding in detail by calculating Mulliken populations and bond orders, concluding that the bonding across the interface is strongly ionic. The work of separation $`W_{sep}`$ of the interface was calculated, and found to be very high: of order 10 Jm<sup>-2</sup> when niobium was bonded to the oxygen terminated Al<sub>2</sub>O<sub>3</sub> surface. Lower energy pathways for the cleavage of this interface would be within the Nb metal or the oxide itself. Two other interfaces were studied corresponding to the two other possible terminations of bulk Al<sub>2</sub>O<sub>3</sub> (0001); namely the stoichiometric, aluminium termination (one layer of aluminium) and the aluminium-rich termination (two layers of aluminium). HRTEM could not distinguish between the stoichiometric termination and the oxygen termination, however evidence from electron energy loss spectroscopy favoured the oxygen termination.
We point out here that as far as we know the question of which termination is more stable has not yet been addressed in all the theoretical work which has been published so far on any oxide-metal interfaces. The structural predictions have been confined to the question of the relative displacement of the crystals, parallel and perpendicular to the interface, and the local relaxations of atoms at the interface, as well as the energy needed to separate the crystals $`W_{sep}`$. This has been done for interfaces with different terminations or local stoichiometry; all calculations were carried out with atoms at rest ($`T=0K`$) and minima in the total energy were located as a function of atomic positions. However, the question as to whether the oxygen terminated or the aluminium terminated interface is more stable was not discussed. It is well known that this question can only be answered with respect to the chemical potentials of the species in the environment with which the interfaces are in equilibrium, which is normally characterised by temperature and partial pressure of oxygen , and the difficulty of relating the quantities accessible to a first principles calculation to these parameters may have been a reason for leaving this question to one side.
The main purpose of our present paper is to show how in fact we are already able to make predictions of the stability of different interfaces when they differ not only in structure but also in composition. With certain simplifying assumptions we show how this can now be done with little more effort than the calculations which need to be done to calculate the work of separation, and we present first results for the Nb/Al<sub>2</sub>O<sub>3</sub> interface. The ingredients of the theory are the works of adhesion and surface energies. For these we draw upon the results reported briefly in , supplemented by some further calculations to discuss the case of oxygen on the Nb surface. The basic theory is outlined in Section II. We derive the equations for a general $`A_mO_n`$ oxide in contact with a metal $`B`$; it would be a short step to generalise them still further to an interface between arbitrary compounds. Essentially the same thermodynamics was applied by Wang et al in calculations of the surface energy of oxides with different terminations, over a range of chemical potentials of oxygen; our theory makes the further connection to the temperature and in particular the pressure of oxygen, which are the parameters directly under the control of the experimentalist. A detailed study of the Al<sub>2</sub>O<sub>3</sub> surface will be reported elsewhere .
The plan of the rest of the paper is as follows. Sections III-V cover aspects of our total energy calculations which were not dealt with, or dealt with only briefly, in our Letter . In Section III our method of total energy calculation is summarised. In Section IV we describe the atomic relaxations parallel to the interfaces which are generally not commented upon. Although this structural aspect is not central to the thrust of our paper, it turned out that lateral relaxations also have a strong part to play in determining the interplanar relaxations and energies reported previously, and we therefore describe them for completeness. In Section V we describe and comment on the results for the work of separation on different planes and with different terminations of the interface. Our calculated interfacial free energies are presented in Section VI and we conclude in Section VII.
## II Principles of Calculating Interfacial and Surface Energies
Let us consider the interface between metal $`B`$ and an oxide of metal $`A`$ in equilibrium at temperature and pressure $`(T,P)`$. The stoichiometric composition of the $`A`$ oxide is $`A_mO_n`$. We obtain the definitions of interfacial quantities by referring to the contents of a periodically repeated supercell of area $`S`$ parallel to the interfaces which it may contain. All extensive thermodynamic quantitites in the following will refer to the contents of such a supercell. The interfacial energy per unit area, counting the two interfaces within each supercell, is given by:
$`\gamma _{int}`$ $`=`$ $`(G_{int}(T,P)N_A\mu _A(T,P)N_O\mu _O(T,P)`$ (2)
$`N_B\mu _B(T,P))/2S,`$
where $`G_{int}`$ is the Gibbs energy of the contents of a supercell containing two interfaces, $`\mu _A`$, $`\mu _B`$ and $`\mu _O`$ are the chemical potentials of the three components, and $`N_A`$, $`N_B`$ and $`N_O`$ are the numbers of atoms of the three components within the supercell. The denominator $`2S`$ occurs because there are two interfaces in the supercell, as required by periodic boundary conditions. Chemical potentials here are per atom rather than per mole, which would be the usual convention for macroscopic thermodynamics. Special cases of Eqn. (2) are when either the metal $`B`$ or the oxide is absent from the supercell, in which cases we recover expressions for the surface energies of the oxide $`\gamma _{AO}`$ and the metal $`\gamma _B`$ respectively:
$$\gamma _{AO}=(G_{SAO}(T,P)N_A\mu _A(T,P)N_O\mu _O(T,P))/2S,$$
(3)
$$\gamma _B=(G_{SB}(T,P)N_B\mu _B(T,P))/2S.$$
(4)
The quantities $`G_{SAO}`$ and $`G_{SB}`$ are the Gibbs energies of slabs of oxide and of metal, with free surfaces separated in their respective supercells by an adequate thickness of vacuum. We have assumed in Eqn.(4) that the metal surface is clean; this will suffice for a calculation of the energy of the $`\mathrm{𝑖𝑛𝑡𝑒𝑟𝑓𝑎𝑐𝑒}`$ discussed below. However, we can easily consider for example adding a monolayer of oxygen to the metal surface in the calculation of $`G_{SB}(T,P)`$. Since there is no separate oxide phase, the number of oxygen atoms in the system, $`N_O`$, now resides on the metal surface. The contribution $`N_O\mu _O(T,P)`$ must be subtracted in the calculation of the corresponding surface energy just as in the calculation of the interface energy. We shall in fact make this calculation of an ‘oxidised’ Nb surface in the course of obtaining the work of separation of an interface by a pathway which leaves oxygen on the exposed metal surface.
The motivation for calculating $`\gamma _{int}`$ is as follows. An interface between two crystals requires five parameters for its macroscopic specification, for example three to specify the relative crystallographic orientation of the materials and two more to specify the orientation of the interface. We note in passing that a free surface in contact with vapour or liquid only requires two parameters to specify its crystallographic orientation. There is always a large set of hypothetical interfaces which have the same five macroscopic parameters but which differ in their atomic structure and local stoichiometry. The member of this set which minimises $`\gamma _{int}`$ for given chemical potentials is the equilibrium interface. So provided we know the chemical potentials of the components, we could in principle predict the atomistic structure of the equilibrium interface, including its local stoichiometry, by evaluating $`\gamma _{int}`$ for each member of the set. In practice, of course, we can only calculate $`\gamma _{int}`$ for a small subset of the entire set, and rely on our experience and intuition, together with experimental information, to ensure that we have not omitted an important structure. Prior to the present work we and others have calculated total energies for a number of structures at 0K, in which the atomic positions are relaxed by energy minimisation. The equations to be derived below show how to go the two important steps further, namely to correct the 0K information to finite temperature and to take account of local stoichiometry.
The relationship between local stoichiometry and $`\gamma _{int}`$ is well known in thermodynamics as the Gibbs adsorption equation, in which local stoichiometry is measured in terms of excesses $`\mathrm{\Gamma }_i`$ of one or more components labelled $`i`$. Our final version of Eqn. (2) will be in terms of the excess of oxygen at the interface with respect to the metal $`A`$, per unit surface area, which is defined as:
$$\mathrm{\Gamma }_O=(N_O\frac{n}{m}N_A)/2S.$$
(5)
This choice of component $`i`$ is arbitrary; we could equally well work in terms of $`\mathrm{\Gamma }_{Al}`$, because they are related through
$$m\mathrm{\Gamma }_O+n\mathrm{\Gamma }_{Al}=0.$$
(6)
We note three further points in connection with excesses. Firstly, each further component in the system would introduce another excess, each excess being referred to the same designated component. Secondly, a stoichiometric interface is by definition one for which all the excesses vanish. Finally, since one of the phases is the pure metal $`B`$, there can be no excess of the metal $`B`$. In particular, if $`B`$ were in fact also $`Al`$, thereby reducing the number of components to two, the interface could not be described as having an excess of $`O`$ or $`Al`$. For more discussion of the thermodynamics of excess quantities the reader is referred to .
A difficulty up to now has been to calculate the chemical potentials involved in these equations and, more specifically, to relate them to given experimental conditions. In the following we show how Eqn. (2) can be reformulated to relate $`\gamma _{int}`$ to the partial pressure of oxygen $`P_{O_2}`$.
First we define the work of separation $`W_{sep}`$ of the interface. It does not refer to chemical equilibrium states and therefore does not involve chemical potentials of the separate components:
$`W_{sep}`$ $`=`$ $`(G_{SAO}+G_{SB}G_{int})/2S`$ (7)
$`=`$ $`\gamma _{AO}+\gamma _B\gamma _{int}.`$ (8)
For brevity we do not explicitly indicate the temperature and pressure dependence of all the terms unless it needs to be emphasised. An important point to note about this quantity, which makes it relatively straightforward to calculate, is that the separate slabs of metal and oxide have exactly the same composition as the two slabs which are joined to form an interface. This would not in general be the case if these slabs and the interface were in equilibrium with a given environment (constant $`\mu _i`$), because one would expect for example some loss or gain of surface oxygen or metal from the oxide to the vapour phase when the surfaces are created. If the interface as well as the exposed surfaces are the ones which are in chemical equilibrium (which we denote by superfix $`eq`$), an equation similar to (8) defines the work of adhesion:
$$W_{ad}=\gamma _{AO}^{eq}+\gamma _B^{eq}\gamma _{int}^{eq}.$$
(9)
This is the quantity of relevance to contact angles and wetting for example, and unlike $`W_{sep}`$ it is not obtainable by a simple comparison of three total energies.
Calculations of $`W_{sep}`$ for Nb/Al<sub>2</sub>O<sub>3</sub> were reported in our Letter and these have been extended here, as described in the following sections. $`W_{sep}`$ is probably more relevant than $`W_{ad}`$ in formulating a fracture criterion, when internal surfaces are formed which are not in equilibrium, but in order to predict the equilibrium structure of interfaces we also need to be able to evaluate Eqns.(2)-(4).
We now introduce the quantity $`g_{AO}`$, the Gibbs energy per formula unit of bulk $`A_mO_n`$ in equilibrium with metal $`A`$ and oxygen in vapour form:
$$g_{AO}(T,P)=m\mu _A(T,P)+n\mu _O(T,P),$$
(10)
so that
$$G_{AO}=(N_A/m)g_{AO}$$
(11)
is the Gibbs energy of a stoichiometric cell containing $`N_A`$ atoms of $`A`$. Inserting (10) and (5) into (3) gives the surface energy of the oxide in a form which makes the effect of the excess oxygen explicit:
$$\gamma _{AO}=(G_{SAO}G_{AO})/2S\mathrm{\Gamma }_O\mu _O.$$
(12)
Consider now how to go about calculating the two surface energies from (4) and (12), which we will eventually combine with $`W_{sep}`$ in (8) to give us $`\gamma _{int}`$. The Gibbs energy of all slabs can be calculated at $`T=0K`$ and $`P=0`$ from first principles, it is just the total energy. If the slabs are bulk pure material, their Gibbs energy can be corrected to temperature T by using experimental specific heat data. On the other hand when the slabs are separated in the supercell by a layer of vacuum to represent free surfaces, there is no such experimental data and the correction to finite $`T`$ could be done by calculating the phonon spectrum and using the quasiharmonic approximation for the free energy. This has been done previously for classical ionic models by Taylor and coworkers, in order to obtain the temperature dependence of their surface energy, but we have not yet made the equivalent calculation with our ab initio code. For a metal slab (Ag), the quasiharmonic free energy based on ab initio phonon frequencies was recently calculated by Xie and coworkers. In the case of the pure metal slab, the chemical potential $`\mu _B`$ is the Gibbs energy per atom of a bulk slab. The surface energy of $`B`$ is therefore obtained from the results of two supercell total energy calculations in the standard way. The main present issue, which is less familiar in the context of total energy calculations, is how to calculate the significant term due to the chemical potential of oxygen, which must be included when the surface of the oxide is non-stoichiometric ($`\mathrm{\Gamma }_O0`$).
The chemical potential of oxygen is well described in terms of its partial pressure $`P_{O_2}`$ by the standard ideal gas expression
$$\mu _O=\mu _O^0+\frac{1}{2}kT\mathrm{log}(P_{O_2}/P^0).$$
(13)
In Eqn. (13), $`\mu _O^0`$ is the oxygen chemical potential in its standard state (STP) at $`T^0`$=298.15K, $`P^0`$=1at. Chemists would set $`\mu _O^0`$ to zero by definition, but we cannot do that since our zero of energy is already defined as the energy of separated ions and electrons at $`T=0`$K. On the other hand the energy of oxygen molecules is not something we want to calculate, since there are well known problems in using density functional theory for this system. Fortunately, we can circumvent the problem by using a thermodynamic cycle. From the defining equation for the standard Gibbs energy of formation $`\mathrm{\Delta }G_{AO}^0`$:
$$g_{AO}^0=m\mu _A^0+n\mu _O^0+\mathrm{\Delta }G_{AO}^0,$$
(14)
we obtain the troublesome oxygen chemical potential at STP in terms of $`g_{AO}^0`$ and $`\mathrm{\Delta }G_{AO}^0`$. The quantities $`g_{AO}^0`$ and $`\mu _A^0`$ are things we can calculate accurately, and we can look up $`\mathrm{\Delta }G_{AO}^0`$ in tables of thermodynamic data.
Inserting $`\mu _O^0`$ from (14) into (13) and (13) into (12) gives us our final expression for the surface energy of the oxide:
$`\gamma _{AO}`$ $`=`$ $`(G_{SAO}(T,P){\displaystyle \frac{N_A}{m}}g_{AO}(T,P))/2S`$ (17)
$`\mathrm{\Gamma }_O(g_{AO}^0m\mu _A^0\mathrm{\Delta }G_{AO}^0)/n`$
$`\mathrm{\Gamma }_O{\displaystyle \frac{1}{2}}kT\mathrm{log}(P_{O_2}/P^0).`$
from which we obtain the final expression for the interfacial energy by subsituting (17) into (8):
$`\gamma _{int}`$ $`=`$ $`\gamma _B(T,P)W_{sep}(T,P)`$ (21)
$`+(G_{SAO}(T,P){\displaystyle \frac{N_A}{m}}g_{AO}(T,P))/2S`$
$`\mathrm{\Gamma }_O(g_{AO}^0m\mu _A^0\mathrm{\Delta }G_{AO}^0)/n`$
$`\mathrm{\Gamma }_O{\displaystyle \frac{1}{2}}kT\mathrm{log}(P_{O_2}/P^0).`$
The quantities $`g_{AO}^0`$ and $`\mu _A^0`$ entering the third line of (21) are well described by $`T=0`$K quantities which we calculate. It can be verified that correcting them to standard state has a negligible effect on the surface energy.
The minimum physically meaningful value of $`P_{O_2}`$, which we denote $`P_{O_2}^{min}`$, is set by the condition that if $`P_{O_2}P_{O_2}^{min}`$ the oxide would spontaneously decompose into metal and oxygen. Neglecting the small variation in solid energies with temperature by comparison with $`\mathrm{\Delta }G_{AO}^0`$ this condition is:
$$\mathrm{log}(P_{O_2}^{min}/P^0)=\frac{2}{nkT}\mathrm{\Delta }G_{AO}^0.$$
(22)
Similarly, the maximum physically meaningful value of $`P_{O_2}`$ is defined by the lowest standard Gibbs energy of formation of a metal B oxide $`\mathrm{\Delta }G_{BO}^O`$:
$$\mathrm{log}(P_{O_2}^{max}/P^0)=\frac{2}{n^{}kT}\mathrm{\Delta }G_{BO}^0$$
(23)
where the first oxide to form would have the stoichiometry $`B_m^{}O_n^{}`$. The thermodynamic data used here are summarised in Table I.
## III Method of total energy calculation
For the interface calculations we use the total energy plane wave pseudopotential method based on Lanczos diagonalization of the Kohn-Sham density matrix. The supercell has the form of a rhombohedral prism and in the stoichiometric slab it contains 45 atoms: 14 Al, 21 O and 10 Nb atoms (see Fig.1). By stripping off the outer plane of Al from each interface we obtain an interface which is O-terminated with an O excess $`\mathrm{\Gamma }_OS=+1.5`$ atoms per surface unit cell. By adding the surface plane of Al atoms to the neutral interface we make an oxygen poor interface, with the negative O excess $`\mathrm{\Gamma }_OS=1.5`$. The total energy of the contents of a supercell is minimized with respect to the atomic coordinates by the quasi-Newton method with Hessian updated using the Broyden-Fletcher-Goldfarb-Shano (BFGS) method. The pseudopotential for Nb was of Troullier-Martins form , with $`s`$ and $`d`$ non-locality. The pseudopotential for O was of optimised form , with $`s`$ non-locality. The pseudopotential for Al was of Gonze type with $`s`$ non-locality.
All calculations were made with two $`k`$-points in the irreducible wedge of the Brillouin zone, and with a plane-wave cut-off of 40 Ry. The effect of increasing the plane wave cutoff from 40 to 60Ry was to reduce $`W_{sep}`$ by 3.3% for the Nb/Al interface, which we take as a satisfactory indication of the basis set convergence. For the neutral 45 atom interface we have made test calculations with six $`k`$-points which results in a decrease of total energy by about 4mRy and very small ($`<10^3`$ nm ) changes of relaxed positions of atoms compared with two $`k`$-point calculations. The effect of increasing the $`k`$-point sampling from 2 points to 9 is to change $`W_{sep}`$ by less than 1%.
By doubling the original unit cell in the x-y plane we obtained a 180 atom cell, with which we recalculated the wavefunctions at the gamma point with the previously relaxed atomic coordinates. The gamma point wavefunctions in this cell were used for Mulliken population analysis which was made by projecting the optimized wave functions onto the pseudoatomic orbitals $`|\varphi _{i\alpha }>`$ (i labels site, $`\alpha `$ \- orbitals) according to the procedure suggested in. The “spillage” of each occupied orbital $`\psi `$ was less than 1.5%.
## IV Relaxation of the interface
The slab with which the Nb(111)/Al<sub>2</sub>O<sub>3</sub>(0001) interface was modelled is shown in Fig.1. The interlayer relaxation of the interface has been reported previously, and we refer to that paper for results. Here we mention a feature which has not previously been discussed, namely the relaxations parallel to the interface, which we refer to as in-plane relaxations. It has been found that to make a calculation of the interlayer relaxation of the alumina surface one needs to take into account the in-plane relaxations of the oxygen atoms, which were neglected in some earlier work. The present results show that in-plane relaxation of the oxygen ions is a general feature of the structure near the interface. The geometry of these relaxations is described by two parameters; the rotation $`\alpha `$ and the bond length extension $`\delta r`$ of the equilateral triangle of oxygen atoms in a plane, illustrated in Fig.2. The rotation and dilation of this equilateral triangle does not break any symmetry, preserving for example the three fold axis about the centre of the triangle concerned.
The calculated values of $`\alpha `$ and $`\delta r`$ are shown in Table II for the terminating and second layer oxygen planes (denote by subscripts 1 and 2) in five cases. The first two cases, labelled N(b)/A(O) and N(b)/A(Al) are the O and Al terminated bulk Nb/Al<sub>2</sub>O<sub>3</sub> interfaces. Cases three and four are pure alumina surfaces, labelled A(O) and A(Al) to indicate that they are oxygen and aluminium terminated respectively. Case five, labelled N(m)/A(O) is a Nb monolayer on an oxygen terminated alumina surface. One can see that the in-plane relaxation is a feature of all the systems studied. From the evidence of the first two layers, the rotation of O-triangles and the increase of the O-O bond lengths appears to be localised near the surface of the alumina. The Nb monolayer on the O-terminated surface of alumina shows this effect most strongly, which is quite surprising, since the interplanar relaxation in this case is much less than that of the Al layer for which the Nb substitutes. There is experimental confirmation of the effect, obtained by small angle X-ray diffraction, in the case of the A(Al) surface. In this case the experimental results are $`\delta r_1=4.5\%`$ and $`\alpha _1`$=$`6.7^{}`$, compared with our calculated results of $`\delta r_1=3.2\%`$ and $`\alpha _1`$=$`3.1^{}`$. The agreement is only qualitative.
## V Work of separation
All our results for the calculated work of separation $`W_{sep}`$ of different interfaces and cleavage planes are shown in Table III. The column of ‘unrelaxed’ results refers to values obtained by assuming bulk unrelaxed atomic positions both at the interface and for the free surfaces. The interplanar spacing between O and Nb across the interface in this case was simply taken as the bulk spacing between O and Al planes. The results in the ‘relaxed’ column are calculated with atomic positions relaxed both before and after cleavage.
The effect of relaxations on $`W_{sep}`$ naturally depends on the interface. It is most pronounced when an Al terminated Al<sub>2</sub>O<sub>3</sub> surface is exposed, because of the large relaxation of this surface, which lowers its surface energy by about 1.5 Jm<sup>-2</sup>. On the other hand, if the relaxation of the interface dominates the energy balance in Eqn.(8) then the relaxed value of $`W_{sep}`$can even be larger than the unrelaxed value, as in the case of cleavage between Nb and O at the N(b)/A(O) interface.
The lowest value of $`W_{sep}`$, 2.7 Jm<sup>-2</sup>, is found for the cleavage of bulk Nb from the stoichiometric Al-terminated alumina. The highest values are found for the cleavage of bulk Nb from the O-terminated alumina surface. Indeed we can deduce from Table III that this interface would be unlikely to separate between Nb and O planes, but would prefer to separate inside the Nb, leaving a monolayer of Nb on the surface, or even between O and Al, leaving a monolayer of O on the Nb surface.
The highest value (relaxed) of all in Table I is for the cleavage of pure Al<sub>2</sub>O<sub>3</sub> between O and Al planes. Experimentally, $`\alpha `$Al<sub>2</sub>O<sub>3</sub> does not cleave on the basal plane at all, but its lowest energy cleavage on this plane would clearly be between Al planes. This is what one expects on the basis of charge neutrality arguments, because by cleaving between Al planes two identical, neutral surfaces are created. On the other hand by cleaving between O and Al, different surfaces are created which, in order to be neutral, require the oxygen or aluminium at the surface to be in an unfavourable valence state, hence this is a final state of especially high energy.
The above interpretations are supported by the Mulliken populations shown in Table IV for the three interfaces between bulk Nb and Al<sub>2</sub>O<sub>3</sub>. We make the usual caveat here that Mulliken charges do not have absolute significance, since they depend on the choice of basis set, but they are nevertheless a useful indicator of trends in ionicity or covalency. The interfaces in Table III are labelled by their oxygen excesses, to highlight certain trends with the stoichiometry of the interface. Bulk oxygen carries a Mulliken charge of -1, and for the oxygen plane nearest the interface this value is reduced to -0.99, -0.93 and -0.86 in turn as the excess of oxygen at the interface is increased from negative to positive. The change is rather modest, indicating that oxygen does not readily alter its valence state. The charge on the interfacial oxygen is provided by the terminating layer of Al in the case of Al termination, or in the case of the oxygen terminated interface the electrons are provided mainly by the first two layers of Nb.
## VI Interfacial free energy and oxygen pressure
Five surface energies are shown in Fig.3a as a function of $`P_{O_2}`$. The x-axis is appropriate to a temperature of 1500K; to obtain the results at temperature $`T`$ the numbers on the x-axis should be multiplied by $`1500/T`$. The O-rich ($`\mathrm{\Gamma }_OS=1.5`$) and O-poor ($`\mathrm{\Gamma }_OS=1.5`$) alumina surfaces have negative and positive slopes respectively, while the stoichiometric Al<sub>2</sub>O<sub>3</sub> and pure Nb(111) surface energies are constant, and by chance nearly equal. The most negatively sloping surface energy we have plotted here refers to the Nb(111) surface with an attached monolayer of oxygen. It becomes negative at an oxygen pressure inside the regime of stability of NbO.
The interfacial free energy from Eqn.(2) is shown as a function of $`P_{O_2}`$ in Fig.3c for three interfaces, O-rich, stoichiometric and O-poor ($`\mathrm{\Gamma }_OS=1.5,0,1.5`$). To generate the interfacial free energies one has to subtract the work of separation shown for convenience in Fig.3b, from the sum of the equilibrium surface energies of the two corresponding free surfaces. With increasing $`P_{O_2}`$ the O-rich interface becomes increasingly stable, the Al-rich interface less stable and the free energy of the stoichiometric interface remains constant, exactly parallel to the behaviour of the free surfaces. The interfaces can only be in thermodynamic equilibrium in the range of oxygen pressure which is indicated on the figure. At 1500K this range is as given in Table I; at values of $`P_{O_2}`$ above this range, the Nb would oxidise to NbO, and below it the alumina would decompose.
The work of adhesion at a given $`P_{O_2}`$ can also be estimated from the results on this graph using Eqn.(9). The free surface and interface energies should be those with lowest free energy at the given oxygen pressure, and these can be read off from Figs.3a and 3c. The result is plotted in Fig.3d.
## VII Discussion and Conclusions
We have made a careful distinction between work of separation, a mechanically defined quantity, and work of adhesion, a thermodynamic quantity, focusing on how to go about calculating these quantities within an atomistic model. We particularly consider the interface between a metal and an oxide, since it is of practical importance and since oxygen is a troublesome component for which to calculate the chemical potential, a key quantity in interfacial energies. A useful practical equation for the free energy of an interface involving oxygen has been derived, namely Eqn.(9), which gets around the previous difficulty by using a thermodynamic cycle to express the result in terms of quantities which can be readily calculated, namely the total energies of slabs, and quantities which can be obtained from tables, namely the standard Gibbs energy of formation of the oxide and the Gibbs energies of the bulk materials relative to their $`T=0`$K values.
To illustrate and apply the method we have made a number of first principles calculations for Nb(111)/Al<sub>2</sub>O<sub>3</sub> (0001) interfaces, oxygen rich, oxygen poor and stoichiometric, and for several surfaces. We fully relax the atomic positions in supercells using a plane wave, pseudopotential methodology. The relaxations are significant, and in all cases they involve in-plane as well as interlayer relaxations of the oxygen ions. Results on the work of separation of these interfaces were given in a Letter recently, and we have extended them to include the possibility of a cleavage of the O-terminated interface which leaves the Nb coated with oxygen. This turns out indeed to be a lower energy mode of separation (4.9 Jm<sup>-2</sup>) than the alternative which leaves an oxygen rich Al<sub>2</sub>O<sub>3</sub> surface behind (9.8 Jm<sup>-2</sup>), because the favourable degree of ionicity of oxygen is thereby preserved as it is in both bulk alumina and its stoichiometric surface. Considering further the strongly bound O-terminated Nb/Al<sub>2</sub>O<sub>3</sub> interface, it turns out that the hypothetical processes of (i) cleavage within bulk Nb (4.2 Jm<sup>-2</sup>), or (ii) leaving a monolayer of Nb on the oxide surface (3.8 Jm<sup>-2</sup>), or even (iii) cleavage within bulk Al<sub>2</sub>O<sub>3</sub> (3.9 Jm<sup>-2</sup>) are all marginally of lower energy than the cleavage which takes oxygen with the niobium.
By combining the results of our calculations with thermodynamic data we obtain surface energies and interfacial energies as a function of oxygen partial pressure and temperature. An approximation we make here is in omitting the temperature dependence of the solid state free energy, but we include the $`kT\mathrm{log}(P_{O_2}/P^0)`$ term which describes the temperature dependence of the oxygen chemical potential; this is also the term which describes the dependence of all the interfacial and surface free energies on oxygen pressure. It is clear how a more accurate calculation could be made by implementing the quasiharmonic approximation to correct solid surface free energies, and it will probably become a routine matter to include such a correction in future work. Another approximation is made by considering only a small set of possible interface and surface compositions which we think are representative. Nevertheless, despite the present simplifications, some clear results have emerged.
Of the free surfaces of Al<sub>2</sub>O<sub>3</sub> , the stoichiometric one, terminated by a single layer of Al, is the most stable over the whole range of oxygen partial pressure up to over one atmosphere. It may be that a treatment of the temperature dependence of the energy of the slabs could modify the upper and lower bounds on pressure somewhat. Correction of the LDA error is also likely to lower surface energies by 10-20% (I. G. Batyrev, unpublished). For example, work of C. E. Sims et al with classical potentials indicates that the surface energies of Al-terminated Al<sub>2</sub>O<sub>3</sub> can be reduced by up to 0.2-0.3 Jm<sup>-2</sup> at 1500-2000K. We expect an oxygen terminated surface to be stable at a pressure not too far from atmospheric, but we cannot unfortunately be more quantitative in the prediction at this stage. At very low oxygen pressures it is also reasonable that the experimentally observed Al-rich ($`\sqrt{31}\times \sqrt{31}`$) structure is stable; we cannot model a supercell of the size needed to calculate this. Instead we modelled a much simpler Al-rich interface, which is predicted to become the most stable one just above $`P_{O_2}^{min}`$ where Al<sub>2</sub>O<sub>3</sub> decomposes. Since the experimental $`\sqrt{31}\times \sqrt{31}`$ is a very Al-rich surface ($`\mathrm{\Gamma }_OS=7.5`$ in the present notation), the slope of its surface energy versus $`\mathrm{log}(P_{O_2}/P^0)`$ is correspondingly very steep and positive, and it must intersect all the other surface energies just above $`P_{O_2}^{min}`$.
The Nb free surface should obviously become unstable with respect to some adsorption of oxygen when $`P_{O_2}>P_{O_2}^{max}`$, the pressure at which NbO begins to form. The particular configuration and concentration of an oxygen monolayer which we have calculated is not likely to be the optimum configuration of the first oxygen covered Nb(111) surface, but it does become more stable than the free surface at pressures somewhat above $`P_{O_2}^{max}`$ (Fig.3a).
A significant new result is the theoretical analysis of the thermodynamic stability of the O-terminated interface, the strong bonding of which we discussed above. No interface is thermodynamically stable above the (very low) oxygen pressure at which NbO forms, but over most of the range below this the O-terminated interface is less stable than the Al-terminated one (Fig.3c), despite its strong bonding. In fact at the very lowest pressure of oxygen, as would pertain in the presence of pure aluminium, our prediction is of an Al-enriched interface. The experimental indications from EELS show no evidence for Al-Nb bonding, and suggest rather the existence of the O-terminated interface. According to our analysis this could only be marginally in thermodynamic equilibrium if the oxygen pressure is being ‘buffered’ by Nb/NbO and lies close to $`P_{O_2}^{max}`$, which does not seem unreasonable.
We have not included in our comparison interfaces with a different macroscopic orientation such as the Nb(110)/Al<sub>2</sub>O<sub>3</sub> (0001) interface . Although this interface is believed to be thermodynamically more stable than the Nb(111)/Al<sub>2</sub>O<sub>3</sub> (0001) interface, the kinetic barrier to changing the macroscopic orientation is presumably much greater than the barriers to changing the local interface structure.
Although the formalism has been developed for describing metal-oxide bonding, there are obvious applications to systems in which water or other substances may contaminate surfaces or interfaces. The comparison of the energetics of interfaces with differing amounts of segregation follows the same lines. The application of the present formalism using a thermodynamic cycle to avoid the most difficult calculations may be fruitful in other situations in the field of interface chemistry.
###### Acknowledgements.
We thank J. Hutter for technical help with the calculations. This work has been supported by the UK Engineering and Physical Sciences Research Council under grants No. GR/L08380 and GR/M01753, and by the European Communities HCM Network “Electronic Structure Calculations of Materials Properties and Processes for Industry and Basic Science” under grant No. ERBFMRXCT980178. The Centre for Supercomputing in Ireland is gratefully acknowledged for computer resources. |
warning/0001/cond-mat0001028.html | ar5iv | text | # Itinerant ferromagnetism in half-metallic CoS2
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## Abstract
We have investigated electronic and magnetic properties of the pyrite-type CoS<sub>2</sub> using the linearized muffin-tin orbital (LMTO) band method. We have obtained the ferromagnetic ground state with nearly half-metallic nature. The half-metallic stability is studied by using the fixed spin moment method. The non-negligible orbital magnetic moment of Co $`3d`$ electrons is obtained as $`\mu _L=0.06\mu _B`$ in the local spin density approximation (LSDA). The calculated ratio of the orbital to spin angular momenta $`L_z/S_z=0.15`$ is consistent with experiment. The effect of the Coulomb correlation between Co $`3d`$ electrons is also explored with the LSDA + $`U`$ method. The Coulomb correlation at Co sites is not so large, $`U1`$ eV, and so CoS<sub>2</sub> is possibly categorized as an itinerant ferromagnet. It is found that the observed electronic and magnetic behaviors of CoS<sub>2</sub> can be described better by the LSDA than by the LSDA + $`U`$.
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Transition metal disulfides of MS<sub>2</sub> (M = Fe, Co, Ni, etc.) with the pyrite structure shows a variety of physical properties. Electronic and magnetic properties of the MS<sub>2</sub> has been qualitatively explained by the successive filling of the $`e_g`$ band of metal ions. In the NaCl-type structure of pyrite MS<sub>2</sub>, divalent metal ions of M<sup>2+</sup> are located at the center of octahedron formed by S<sub>2</sub><sup>2-</sup> dimers. Hence, the M-$`3d`$ bands are split into $`t_{2g}`$ and $`e_g`$ subbands due to large crystal field. The M-$`3d`$ electron configuration varies from $`t_{2g}^6e_g^0`$, $`t_{2g}^6e_g^1`$, to $`t_{2g}^6e_g^2`$ for FeS<sub>2</sub>, CoS<sub>2</sub>, and NiS<sub>2</sub>. All of them are in the low-spin states, $`S=0`$, $`S=1/2`$, and $`S=1`$ for FeS<sub>2</sub>, CoS<sub>2</sub>, and NiS<sub>2</sub>, respectively. With fully occupied $`t_{2g}^6`$ and completely empty $`e_g^0`$ bands, FeS<sub>2</sub> is a semiconductor with an energy gap of $`E_g0.8`$ eV. CoS<sub>2</sub> is a ferromagnetic metal with $`T_C120`$ K. NiS<sub>2</sub> is an antiferromagnetic insulator ($`T_N40`$ K) with a half-filled $`e_g^2`$ band due to the large on-site Coulomb correlation.
Despite extensive studies on these systems, the electronic and magnetic properties of CoS<sub>2</sub> are not fully understood yet. The issues for CoS<sub>2</sub> to be resolved are $`(i)`$ whether it is a strongly correlated electron system, $`(ii)`$ whether it is classified as an itinerant electron ferromagnet, and $`(iii)`$ whether it is a half-metal or not. It is considered that Co $`3d`$ states in CoS<sub>2</sub> are just at the boundary of the localized and delocalized regime. Magnetic measurement using the polarized neutron diffraction indicates that most part of the magnetic moment is localized at Co sites. The analysis of photoemission spectroscopy (PES) data for CoS<sub>2</sub> in terms of the cluster model calculations yields a large Coulomb correlation parameter, $`U=3.04.2`$ eV. According to the magnetic circular dichroism (MCD) measurements, a non-negligible orbital magnetic moment is observed for Co $`3d`$ electrons, $`L_z/S_z=0.18`$ and $`0.14`$. Such an unquenched orbital moment is exceptional, considering the low-spin state of Co $`3d`$ electrons in CoS<sub>2</sub> with a large crystal field splitting of $`10D_q2.5`$ eV. On the other hand, transport properties of CoS<sub>2</sub> are known to be well described by the itinerant electron model. Further, the dynamical susceptibility measured by the inelastic magnetic neutron scattering and optical properties are analyzed based on the itinerant band model.
Using the linearized atomic orbitals (LCAO) band method, Zhao et al. has found that CoS<sub>2</sub> is close to half-metallic in the sense that the occupation number of the minority spin $`e_g`$ band is very small compared to that of the majority spin $`e_g`$ band. In fact, recent experimental reflectivity data have been consistently interpreted under the assumption of half-metallic CoS<sub>2</sub>. In contrast, the linearized muffin-tin orbital (LMTO) band calculation by Yamada et al. gives rise to the normal-metallic ferromagnetic ground state with a partially filled minority spin $`e_g`$ band.
To clarify the above mentioned issues and to resolve the differences between existing band calculations, we have reexamined electronic and magnetic properties of CoS<sub>2</sub> using the LMTO band method. Both the local spin density approximation (LSDA) and the LSDA + $`U`$ approximation are employed to explore the effect of the Coulomb correlation interaction between Co $`3d`$ electrons. The von Barth-Hedin form of the exchange-correlation potential is utilized. Further, the fixed spin moment method is used to check the stability of half-metallic state. We have also estimated the orbital contribution to the total magnetic moment in CoS<sub>2</sub>, using the fully relativistic calculations in which the spin-orbit coupling is simultaneously included in the self-consistent variational loop. Two main parameters in the LSDA + $`U`$ method are the Coulomb $`U`$ and the exchange $`J`$ interactions. The relations of these parameters to the Slater integrals are given by $`U=F^0`$ and $`J=(F^2+F^4)/14`$. The ratio of $`F^4/F^2`$ is known to be constant around 0.625 for most $`3d`$ transition metal atoms. We have used parameter values of $`U=0.04.0`$ eV with a fixed $`J=0.89`$ eV.
The crystal structure of pyrite CoS<sub>2</sub> is cubic with the lattice constant of $`a=5.528`$ Å and the space group is $`T_h^6(Pa\overline{3})`$. There are twenty-four symmetry operations. The four Co atoms are located at positions $`(4a)`$ : (0,0,0), (1/2,0,1/2), (0,1/2,1/2), and (1/2,1/2,0). The eight S atoms are at positions $`(8c)`$ : $`(u,u,u)`$, $`(u+1/2,u,\overline{u}+1/2)`$, $`(u,\overline{u}+1/2,u+1/2)`$, and $`(\overline{u}+1/2,u+1/2,u)`$ where we have used the position parameter of $`u=\pm 0.389`$. To improve the packing ratio, twenty-four empty spheres are introduced at positions $`(24d)`$. For the Brillouin zone (BZ) integration, the tetrahedron method is adopted with 80 k-points sampling in the irreducible BZ wedge. The LMTO basis functions are included up to $`l=2`$ for Co, and $`l=1`$ for S and empty spheres.
In Fig. 1, we show the LSDA density of states (DOS) obtained at the experimental lattice constant. Because of strong covalent bonding nature of S-S dimers, each S band can be identified with molecular orbitals. The two peaks in the highest binding energy side are the bonding S $`3s\sigma `$ and the antibonding S $`3s\sigma ^{}`$ bands, respectively. The broad band of about 6 eV width between $`4`$ and $`10`$ eV consists mainly of S $`3p`$, which is a mixture of $`3p\sigma `$, $`3p\pi `$, and $`3p\pi ^{}`$ bands. The antibonding S $`3p\sigma ^{}`$ intra-dimer band is located at about $`2.5`$ eV above the Fermi energy $`E_\mathrm{F}`$. The spin split bands of Co $`t_{2g}`$ are located between $`1`$ and $`3`$ eV below $`E_\mathrm{F}`$ while Co $`e_g`$ band are near $`E_\mathrm{F}`$. It is seen that Co $`3d`$ bands are formed in the energy range between S $`3p`$ and S $`3p\sigma ^{}`$ bands. The $`t_{2g}`$ band is nearly dispersionless with very narrow band width, whereas the $`e_g`$ band near $`E_\mathrm{F}`$ is dispersive with relatively large band width of about $`2.5`$ eV due to strong hybridization with S $`3p\sigma ^{}`$ band. Nearly half-metallic nature and the exchange splitting $`\mathrm{\Delta }_X1.0`$ eV are prominent for $`e_g`$ band.
The present results are slightly different from those of existing calculations. In Ref. REFERENCES, S $`3p\sigma ^{}`$ band is at $`1.5`$ eV above $`E_\mathrm{F}`$ which is lower than the present result by about $`1.0`$ eV. Whereas, in Ref. REFERENCES, the states of S $`3p\sigma ^{}`$ band is smeared out by too strong hybridization with Co $`e_g`$ bands. The bremsstrahlung isochromat spectrum (BIS) shows clear double peak structure at about $`1.0`$ and $`2.5`$ eV above $`E_\mathrm{F}`$, which is consistent better with the present result. That is, the $`1.0`$ eV peak corresponds to the minority spin $`e_g`$ band and the $`2.5`$ eV peak to S $`3p\sigma ^{}`$ band. The occupation of the minority spin $`e_g`$ band is only $`0.04`$ electrons/Co, reflecting that CoS<sub>2</sub> is almost half-metallic. This feature is in agreement with that obtained by Zhao et al.. In the analysis of the reflectivity data in terms of the half-metallic CoS<sub>2</sub>, the minority spin $`e_g`$ band is assigned at about 1.5 eV above $`E_\mathrm{F}`$. This is close to but a bit higher than the present LSDA result.
Figure 2 is the LSDA total energy as a function of the lattice constant. The ferromagnetic ground state is lower in energy than the paramagnetic state even far below the experimental lattice constant of $`a=5.528`$ Å represented by the arrow. The minimum of the total energy is located near the experimental lattice constant for both the paramagnetic and the ferromagnetic states, and thus the LSDA seems to describes well the cohesive bonding properties of CoS<sub>2</sub>. The ferromagnetic ground states is lower in energy by about 10 mRy than the paramagnetic state at the experimental lattice constant. The energy difference between the ferromagnetic and the paramagnetic states decreases with decreasing the lattice constant. This is contrary to the result of Yamada et al. in which the paramagnetic is stable near the experimental lattice constant and the ferromagnetic state becomes stable only for $`a>5.57`$ Å. As they pointed out, it seems to be ascribed to small number of k-points they used. On the basis of the total energy curve shown in Fig. 2, the observed metamagnetic transition behavior under the high magnetic filed, the high external pressure, and the chemical pressure by substituting Se in place of S can be qualitatively understood
By using the fixed spin moment method within the LSDA, the total energy is evaluated with varying the magnetic moment in CoS<sub>2</sub> (Fig. 3). For all the fixed spin moment, $`0.00\mu _B<M4.00\mu _B`$, the ferromagnetic state is stable as compared to the paramagnetic state. The spin direction of S atoms is always parallel to that of Co. The ground state is found for $`M=3.66\mu _B`$ (denoted by the arrow) which is a normal ferromagnetic metal and lower in energy by about 10 mRy than the paramagnetic state. It is noticeable that the half-metallic state with exact integer magnetic moment of $`M=4.00\mu _B`$ is very close in energy to the normal-metallic ground state. The total energy difference between the normal-metallic ground state and the half-metallic state is less than 0.1 mRy because of quite small electron occupation of the minority spin $`e_g`$ band in the ground state. Thus it is expected that, with only small number of hole doping, the half-metallic state can be realized in CoS<sub>2</sub>.
In general, a polycrystalline half-metallic system is expected to show large magnetoresistance (MR) behavior. It arises from the activated electron tunneling mechanism at the grain boundary with the applied external magnetic field . To our knowledge, there has been no report available on the MR measurement for CoS<sub>2</sub>. Therefore, the MR measurement on polycrystalline CoS<sub>2</sub> and Fe<sub>x</sub>Co<sub>1-x</sub>S<sub>2</sub> is desirable to exploit the correlation between the half-metallic nature and the MR.
We have applied the LSDA and the LSDA+$`U`$ methods to explore the Coulomb correlation effect on the magnetic properties. The observed magnetic moments per Co atom are in the range of $`0.840.91\mu _B`$. The orbital magnetic moment in CoS<sub>2</sub> is obtained by including the spin-orbit coupling in the total Hamiltonian. In the LSDA + $`U`$ calculations, we have used fixed $`J(=0.89`$ eV$`)`$ which is insensitive to local environment. In the LSDA, we have obtained the spin and orbital magnetic moments per Co atom, $`\mu _S=0.80\mu _B`$ and $`\mu _L=0.06\mu _B`$, respectively (Table I). This produces the ratio of $`L_z/S_z=0.15`$, which is close to the MCD measurement of $`L_z/S_z=0.18`$ and $`0.14`$. The total magnetic moment per Co atom $`\mu _{tot}=0.86\mu _B`$ in the LSDA is also in agreement with the experiment.
Both the spin and orbital magnetic moments increase with increasing the Coulomb correlation parameter $`U`$. For $`U3.0`$ eV, the minority spin $`e_g`$ band moves to higher energy states and CoS<sub>2</sub> becomes completely half-metallic. For $`U=4.0`$ eV, the minority spin $`e_g`$ band is formed at about 2.0 eV above $`E_\mathrm{F}`$. However, the discrepancy of the total magnetic moment between the experiment and the theory becomes significant for $`U2.0`$ eV. Moreover, the ratio of $`L_z/S_z=0.31`$ for $`U=2.0`$ eV is too large compared to the experimental values. This finding suggests that the size of the on-site Coulomb correlation in CoS<sub>2</sub> is to be rather small $`U1`$ eV, and the LSDA gives a better description of Co $`3d`$ states in CoS<sub>2</sub>. Note that the LDA results for FeS<sub>2</sub> show an excellent agreement with PES and BIS data. In addition, the weak satellite feature in the PES of CoS<sub>2</sub> seems to be consistent with the present results of rather small Coulomb correlation $`U`$.
Finally, let us remark on the electronic structure of CoS<sub>2</sub> in its paramagnetic phase ($`T>T_C`$). As mentioned in the introduction, it is an unresolved issue. The question is whether the exchange splitting in Co $`3d`$ bands persists above $`T_C`$ or not. The present study reveals that the width of Co $`e_g`$ band is larger than both the Coulomb correlation interaction $`U`$ and the exchange splitting $`\mathrm{\Delta }_X`$. It thus suggests that the Co $`3d`$ electrons near $`E_\mathrm{F}`$ will behave as itinerant to manifest Stoner-type magnetic phenomena at the finite temperature. Indeed, there are such evidences from several experiments. Resonance PES data for Co $`3d`$ bands taken at temperatures across the magnetic transition indicate a slight but noticeable spectral change between the paramagnetic and the ferromagnetic phases, reflecting the long range spin exchange effect. Optical spectra above $`T_C`$ are described well using the itinerant Stoner model. Also inelastic magnetic neutron scattering data above $`T_C`$ exhibit the Stoner excitations stemming from the itinerant electrons. On the other hand, based on the PES studies, it is claimed that the electronic structure in the paramagnetic phase should be described by the local band picture . Hence further studies are necessary to ascertain this point.
In conclusion, we have investigated the electronic and magnetic properties of the pyrite-type CoS<sub>2</sub> using the LMTO band method. We have found the following results: $`(i)`$ CoS<sub>2</sub> does not belong to the strongly correlated electron system, because the effect of the Coulomb correlation at Co sites is rather small, $`U1`$ eV, $`(ii)`$ CoS<sub>2</sub> can be categorized as an itinerant electron ferromagnet, since the band width of Co $`e_g`$ state is larger than $`U`$ and the exchange splitting $`\mathrm{\Delta }_X`$, $`(iii)`$ CoS<sub>2</sub> is nearly half-metallic in that the half-metallic state is very close in energy to the normal-metallic ground state. A non-negligible orbital magnetic moment of Co $`3d`$ electrons is obtained as $`\mu _L=0.06\mu _B`$ in the LSDA, which is consistent with experiment. Accordingly, the electronic and magnetic properties of CoS<sub>2</sub> are described better with the LSDA than with the LSDA+$`U`$.
Acknowledgements$``$ The authors would like to thank Jin Ho Park for helpful discussions. This work was supported by the KOSEF (1999-2-114-002-5) and in part by the Korean MOST-FORT fund. |
warning/0001/hep-th0001042.html | ar5iv | text | # References
IFT/32/99
Running Couplings in Hamiltonians
December, 1999
Stanisław D. Głazek
Institute of Theoretical Physics, Warsaw University
ul. Hoża 69, 00-681 Warsaw
Abstract
We describe key elements of the perturbative similarity renormalization group procedure for Hamiltonians using two, third-order examples: $`\varphi ^3`$ interaction term in the Hamiltonian of scalar field theory in 6 dimensions and triple-gluon vertex counterterm in the Hamiltonian of QCD in 4 dimensions. These examples provide insight into asymptotic freedom in Hamiltonian approach to quantum field theory. The renormalization group procedure also suggests how one may obtain ultraviolet-finite effective Schrödinger equations that correspond to the asymptotically free theories, including transition from quark and gluon to hadronic degrees of freedom in case of strong interactions. The dynamics is invariant under boosts and allows simultaneous analysis of bound state structure in the rest and infinite momentum frames.
This article grew out of an invited talk at The Workshop on Light-Cone QCD and Nonperturbative Hadron Physics, Centre for The Subatomic Structure of Matter and The National Institute of Theoretical Physics, University of Adelaide, Adelaide, Australia, December 13-23, 1999.
1. INTRODUCTION
Canonical Hamiltonians of quantum field theories can be written in terms of creation and annihilation operators. Let us denote those operators by $`q`$. In this notation, for example, a triple-gluon vertex in the light-front QCD Hamiltonian has the structure
$$H_Y=\underset{123}{}[123]\delta (1+23)gY_{123}q_1^{}q_2^{}q_3+h.c.,$$
$`(1.1)`$
where symbols $`1`$, $`2`$ and $`3`$ denote colors, spins or momenta of gluons, $`[123]`$ is a shorthand notation for the integration measure over the gluon three-momenta, $`\delta (1+23)`$ is the Dirac $`\delta `$-function of three-momentum conservation, and $`Y_{123}`$ is a function of the gluon quantum numbers, implied by QCD. $`g`$ denotes a bare canonical coupling constant, which is expected from Lagrangian approach to require ultraviolet renormalization. However, the Hamiltonian lacks many of the Lagrangian density symmetries that are employed in the perturbative Lagrangian renormalization procedure and, in all light-front Hamiltonians, transverse and longitudinal directions are treated differently. For that reason, and to introduce a method for renormalization of Hamiltonians such as (1.1), we need to review the origin of the ultraviolet divergences.
One can evaluate matrix elements of $`H_Y`$ between states of the form $`|ijk\mathrm{}=q_i^{}q_j^{}q_k^{}\mathrm{}|0`$. Consider $`12|H_Y|3`$, which is proportional to $`Y_{123}`$. The trouble is that $`Y_{123}`$ does not vanish when the relative motion of particles 1 and 2 becomes very energetic. In other words, the interaction Hamiltonian directly couples states of small kinetic energy to states of arbitrarily high kinetic energy. For example, when we square $`H_Y`$ and evaluate $`_iH_Y|ii|H_Y`$, the range of $`|i`$s to sum over, on the energy scale, is infinite. The sum diverges, and the square of $`H_Y`$ does not exist. If one tries to find eigenstates of the Hamiltonian, the divergence will dominate finite terms. Also, $`\mathrm{exp}iHt`$ diverges and no conclusions can be drawn from knowing $`H`$ as it stands. The problem is worse than the inifinite energy range of interaction implies by itself: the function $`Y_{123}`$ in QCD grows when the energy difference between kinetic energies of particles 1 and 2, and particle 3 grows. The larger are the kinetic energies of the intermediate particles, the more important become the interactions, and we are sent into an abyss of Fock space states without bounds. All physically relevant local quantum field theories have this trouble.
To see the essence of the problem, imagine a matrix of the Hamiltonian matrix elements in the basis of eigenstates of certain $`H_0`$, $`H_0|i=E_{0i}|i`$, $`H_{ij}i|H|j`$. Our problem is that the
corners marked $`lowhigh`$ and $`highlow`$ on the above figure, contain too large matrix elements for the Hamiltonian matrix to have eigenvalues that are independent of the matrix boundaries. In order to understand the boundary-dependence problem of the eigenvalues, we put an upper bound, denoted by $`\mathrm{\Delta }`$, on the basis states energies and we work out what happens with eigenstates of Hamiltonian matrices whose size is limited by the conditions $`E_{0i}\mathrm{\Delta }`$ and $`E_{0j}\mathrm{\Delta }`$, when we change $`\mathrm{\Delta }`$. In particular, we ask what properties must $`H_Y`$ have for the spectrum of $`H_\mathrm{\Delta }`$ to have a limit when $`\mathrm{\Delta }\mathrm{}`$.
Wilson asked this question in case of a model Hamiltonian with big energy gaps between successive energy scales that were included in his calculation, and he studied the influence of coupling between small and large energy states on the lowest eigenvalues. His method for dealing with the $`\mathrm{\Delta }`$-dependence of the spectrum in the multiscale eigenvalue problem (in general, the problem of dependence on regularization, of any kind), was based on an iterative procedure. Initially, one solves the highest energy part of dynamics and focuses on its lowest eigenvalue levels whose dynamics, in turn, is dominated by states with energies lower by one energy gap. Then, one solves this next lower energy scale dynamical problem, and one repeats the process many times. Starting from the energy scale $`\mathrm{\Delta }`$, one eventually arrives at a finite scale $`\lambda `$. This is schematically indicated on the following figure.
Wilson’s approach
In this figure, we see a new small matrix of size $`\lambda `$, denoted by $`H_\lambda `$. This matrix is calculated using an operation $`R_\lambda `$, which is constructed in the sequence of steps lowering the cutoff from $`\mathrm{\Delta }`$ to $`\lambda `$. The construction is based on the principle that the smallest eigenvalues of the small matrix, $`H_\lambda `$, should be the same as the smallest eigenvalues of the big matrix $`H_\mathrm{\Delta }`$. The algebraic derivation of the small Hamiltonian $`H_\lambda `$ is designed to guarantee the equality of the smallest eigenvalues. The crux is that if all matrix elements of $`H_\lambda `$ are independent of $`\mathrm{\Delta }`$ (in general, independent of the regularization one uses to define $`H_\mathrm{\Delta }`$) then, the eigenvalues of $`H_\lambda `$ must be independent of the regularization. Therefore, if we know what to do with the regularization dependence of matrix elements of $`H_\lambda `$, then we know how to go about regularization dependence of the spectrum of $`H_\mathrm{\Delta }`$.
Note, that $`\lambda `$ is finite and can be chosen arbitrarily, as long as $`\mathrm{\Delta }\mathrm{}`$. The set of transformations that connect Hamiltonians with different values of $`\lambda `$ is called renormalization group. Physical results should be independent of $`\lambda `$, by construction. It is clear that $`H_\lambda `$ cannot be equal to merely $`P_\lambda H_\mathrm{\Delta }P_\lambda `$, where $`P_\lambda `$ denotes a projection operator that projects on the space of states with energies $`E_{0i}<\lambda `$. Some additional terms must be included, which reproduce dynamical effects from above $`\lambda `$. Similarly, we do not expect that $`H_\mathrm{\Delta }`$ is merely $`P_\mathrm{\Delta }H_{canonical}P_\mathrm{\Delta }`$, where $`H_{canonical}`$ is built from $`H_0`$ and terms such as (1.1). By the same argument as for $`H_\lambda `$, some additional terms are required in $`H_\mathrm{\Delta }`$, to include effects from above the cutoff $`\mathrm{\Delta }`$. We will call those terms counterterms. The problem is how to find them. According to Wilson, they are found from the condition that all matrix elements of $`H_\lambda `$ become independent of the regularization in the limit $`\mathrm{\Delta }\mathrm{}`$, for all finite values of $`\lambda `$. One must take care not only of the diverging (i.e. $`\mathrm{\Delta }`$-dependent) regularization dependence, but also of the finite regularization effects. The coupling constant $`g`$ in Eq. (1.1) is changed in the renormalization process as a result of introducing counterterms.
The problem with the transformation $`R_\lambda `$ is that the large energy gaps are absent in physically relevant cases, and perturbation theory based on energy scales alone fails. One can try to take advantage of a small coupling constant but we know that naive perturbative expansion does not work in degenerate cases. For example, one may think of effects familiar from elementary degenerate perturbation theory for eigenvalues of a Hamiltonian matrix with a few rows and columns. We know that perturbation theory cannot work unless one properly chooses the initial basis states in the degenerate subspace - eigenstates of the interaction matrix (if it is a few rows and columns). Only in that basis the perturbative limit $`g0`$ exists. Otherwise perturbation theory produces vanishing energy denominators that lead to diverging terms and the calculation is misleading. In case of quantum field theories of interest, the situation is much more involved than in simple matrix case due to multiply degenerated continuous spectra of $`H_0`$. Moreover, in asymptotically free theories, we expect that the interaction strength grows when we go from $`H_\mathrm{\Delta }`$ to $`H_\lambda `$ and small energy denominators are certainly expected to produce large effects. Thus, the degeneracy of spectra and strength of couplings do not allow us to do a precise analysis of $`R_\lambda `$ and $`H_\lambda `$. So, the operation $`R_\lambda `$ is of limited applicability in the canonical approach to quantum field theory.
2. SIMILARITY FOR HAMILTONIANS
There is an alternative approach (see the figure below). Instead of calculating a small Hamiltonian matrix, we can also calculate a narrow matrix. Namely, a similar (in the sense of algebraic similarity) matrix that has the same eigenvalues but whose matrix elements $`H_{\lambda ij}`$ vanish if $`|E_{0i}E_{0j}|>\lambda `$ (or another condition of “narrowness” is satisfied - hermitian matrices can be diagonalized and, therefore, partial diagonalization to a narrow matrix should be possible). The choice of near-diagonal form is motivated by the following property of near-diagonal matrices: when we act with them on a state of some finite energy, a single action of the matrix can rise the energy by at most $`\lambda `$, i.e. by its width on the energy scale. In perturbation theory for eigenvalues of a near-diagonal matrix, corrections will not be sensitive to the cutoff $`\mathrm{\Delta }`$ up to the order $`n\mathrm{\Delta }/(2\lambda )`$, since one has to go up in energies and come back to the initial energy range through action of interactions, and at least $`n`$ are needed to reach the boundary starting from $`\lambda `$. In similarity, the crux is that if matrix elements of the narrow matrix of finite width $`\lambda `$ are independent of regularization when $`\mathrm{\Delta }\mathrm{}`$, then the spectrum of $`H_\lambda `$ will be independent of regularization to all orders of perturbation theory.
Similarity approach
The transformation $`S_\lambda `$ is called similarity transformation. The effective Hamiltonian matrices with various widths $`\lambda `$ are connected by transformations that are called similarity renormalization group transformations for Hamiltonians. We limit our discussion to unitary transformations $`S_\lambda `$. The key feature of the similarity approach is that perturbative construction of $`S_\lambda `$ avoids small energy denominators entirely - they are limited from below by the width $`\lambda `$. In turn, the perturbatively calculated narrow Hamiltonians can be diagonalized numerically, which is the ultimate way to find solutions to complex non-perturbative problems of the original theory. One could ask, why don’t we go all the way to $`\lambda =0`$, which would mean complete diagonalization through $`S_\lambda `$ with $`\lambda =0`$? This is impossible in perturbation theory. We can trust perturbation theory for calculating $`H_\lambda `$ only for not too small values of $`\lambda `$.
In other words, the perturbative similarity transformation $`S_\lambda `$ involves energy changes that are at least as large as $`\lambda `$ and the problem with large effects in perturbative evaluation of effective Hamiltonians is overcome. But that does not mean we eliminated any of the nonperturbative effects. They are still hidden in the narrow effective Hamiltonian, as much as they were in the initial one. The only thing we accomplish through similarity, is the elimination of direct couplings between states of interest to us and very high energy states. This is a prerequisite that we need to define an ultraviolet finite, nonperturbative Hamiltonian eigenvalue problem in quantum field theory.
One can apply various perturbative procedures for calculating $`S_\lambda `$ and $`H_\lambda `$. Wegner invented a beautifully simple scheme for evaluating near diagonal Hamiltonians in solid state physics. It was shown that Wegner’s equation can be employed in the renormalization group scheme. A number of variations reported in the literature in different areas, is growing.
The bottom line is that when solving for the spectrum of a narrow Hamiltonian, we do not have to diagonalize the whole matrix of size $`\mathrm{\Delta }\mathrm{}`$. We can select a window, as is illustrated in the next figure. Diagonalization of that small window is much simpler than diagonalization of the whole matrix and the window eigenvalues match the whole matrix eigenvalues in the middle range of window energies. The reason is that the wave functions have width comparable to $`\lambda `$ on the energy scale, which is indicated on the figure (see also Appendix B in Ref. ). One cannot expect wave functions of eigenstates of the initial Hamiltonian matrix, $`H_\mathrm{\Delta }`$, to be dominated by some small range of energies. In contrast, the eigenstates of matrix $`H_\lambda `$ are expected to have wave functions that have this property. Recent calculations of quarkonium and glueball spectra in light-front QCD exploit this feature.
Window Hamiltonians
Ref. outlines the similarity procedure starting from an initial Hamiltonian, through evaluation of the narrow Hamiltonians (having found the necessary counterterms in the initial Hamiltonian) in perturbation theory, to diagonalization of a small window to get the bound state energy, using a matrix example. The matrix model is asymptotically free and has a bound state. The coupling constant is a function of the effective Hamiltonian width $`\lambda `$, we say it “runs”. For example, it may equal about 0.06 at $`\lambda =65`$ TeV and about 1 at 1 GeV. Still, the window Hamiltonian of a few GeV size can be calculated in second order perturbation theory and the window bound state eigenvalue deviates from exact solution by only 10%. Once we understand that example, we can return to quantum field theories with interactions of the form (1.1) and attempt a calculation of the corresponding “window” Hamiltonians, with running couplings. Calculations can be carried out using the notion of effective particles.
3. SIMILARITY FOR PARTICLES
When we deal with huge Hamiltonian matrices of quantum field theory the number of states is as big as we let it be and the number of matrix elements we have to think about becomes very quickly incomprehensible. We have to reduce the amount of information that we need to know at the beginning. Imagine we would know matrix elements of the interaction $`\alpha /r`$ in atomic basis functions, numerically, but we would not know that they all correspond to the Coulomb force. It would be very hard to relate what happens in one atomic system to what happens in another one. Therefore, when we aim at universal calculations of effective Hamiltonians in theories that contain interactions such as (1.1), we may proceed to a new version of similarity transformation, which avoids dealing directly with Hamiltonian matrix elements in a particular basis and, instead, operates at the level of field operators.
Let us introduce a unitary transformation $`𝒰_\lambda `$ that transforms field operators (denoted here by $`\varphi `$, independently of their spin or other quantum numbers they carry),
$$\varphi _\lambda (x)=𝒰_\lambda \varphi _{\mathrm{}}(x)𝒰_\lambda ^{}.$$
$`(3.1)`$
$`\varphi _{\mathrm{}}(x)`$ denotes a bare quantum field operator that, at any prescribed time, can be expanded into creation and annihilation operators for bare particles in a canonical fashion that we do not need to define here very precisely. $`\varphi _\lambda (x)`$ denotes an operator that is built in exactly the same way from creation and annihilation operators for effective (dressed) particles. This kind of transformation is motivated by physics of hadrons, whose structure can be explained in a constituent quark model. Dressed particles in a given theory interact differently than the bare ones. Namely, bare ones have interactions like (1.1), while the effective ones can only exchange momentum transfers that are limited by $`\lambda `$. This is secured by the construction of $`𝒰_\lambda `$, to be explained below. Therefore, the effective particle wave functions of eigenstates of the Hamiltonian may quickly fall off when momenta or number of the effective particles deviate from the physically dominant values. This is why one can hope to obtain a constituent picture of hadrons in QCD using similarity for particles. More generally, the expected convergence of eigenstate expansion in effective particle basis in Fock space opens a door to studies of few-body systems in quantum field theory.
In order to set up equations that will allow us to calculate Hamiltonians for effective particles, let’s rewrite Eq. (3.1) in terms of the creation and annihilation operators,
$$q_\lambda =𝒰_\lambda q_{\mathrm{}}𝒰_\lambda ^{}.$$
$`(3.2)`$
All we need to do next is: take the bare Hamiltonian of our theory, as it is given initially in terms of $`q_{\mathrm{}}`$, calculate counterterms it needs to contain in addition to the canonical terms, obtain this way our initial $`H_{\mathrm{}}(q_{\mathrm{}})`$, and rewrite it in terms of $`q_\lambda `$. $`𝒰_\lambda `$ is secured to be unitary by construction. The whole point of the construction is that the resulting $`H_\lambda (q_\lambda )`$ is to contain only such interaction terms that, when we evaluate their matrix elements between Fock basis states of effective particles, the resulting effective Hamiltonian matrix is narrow, of width $`\lambda `$, as in the similarity procedure for Hamiltonian matrices we discussed in previous Section. It will not be necessary to go into details here. Only a brief outline of the scheme follows.
Since rewriting the Hamiltonian in different degrees of freedom does not change the operator itself, we have $`H_\lambda (q_\lambda )=H_{\mathrm{}}(q_{\mathrm{}})`$. One may think about $`H_\lambda (q_\lambda )`$ as a QCD Hamiltonian written in terms of constituent quarks and gluons, and about $`H_{\mathrm{}}(q_{\mathrm{}})`$ as the same QCD Hamiltonian written in terms of canonical quarks and gluons, associated with partons, or current quarks (to make the connection between hadronic rest frame constituents and partons in the infinite momentum frame, we have to use the light-front form of Hamiltonian dynamics, see for an outline of light-front QCD in the context of renormalization group procedure for Hamiltonians).
Applying the transformation $`𝒰_\lambda `$, one obtains $`_\lambda H_\lambda (q_{\mathrm{}})=𝒰_\lambda ^{}H_{\mathrm{}}(q_{\mathrm{}})𝒰_\lambda `$. This relation means that the operator $`_\lambda `$ has the same coefficient functions in front of products of $`q_{\mathrm{}}`$ as the effective Hamiltonian $`H_\lambda `$ has in front of the unitarily equivalent products of $`q_\lambda `$. Differentiating $`_\lambda `$ one obtains
$$\frac{d}{d\lambda }_\lambda =[𝒯_\lambda ,_\lambda ],$$
$`(3.3)`$
where the generator $`𝒯_\lambda `$ is related to $`𝒰_\lambda `$ by
$$𝒯_\lambda =𝒰_\lambda ^{}\frac{d}{d\lambda }𝒰_\lambda .$$
$`(3.4)`$
The script letters are introduced to indicate that the operators can be conveniently thought about as expanded into sums of products of operators $`q_{\mathrm{}}`$. The latter are independent of $`\lambda `$ and are not differentiated in Eqs. (3.3) and (3.4). In other words, Eqs. (3.3) and (3.4) describe only the flow of coefficients in front of the creation and annihilation operators. Effective Hamiltonians are obtained from $`_\lambda `$ using $`H_\lambda (q_\lambda )=𝒰_\lambda _\lambda 𝒰_\lambda ^{}`$.
The key element now is how one defines $`𝒯_\lambda `$. This is the domain of similarity for effective particles. In its essence , one studies what one has to do to get the narrow Hamiltonian matrices as a result of the procedure, and these studies tell us what to put for $`𝒯_\lambda `$. There exist infinitely many choices. The one that the present author used to get results described in the next Sections is of the following form
$$[𝒯_\lambda ,_{0\lambda }]=\frac{d}{d\lambda }(1F_\lambda )[𝒢_\lambda ].$$
$`(3.5)`$
The symbols $`𝒢`$ and $`F`$ require explanation. The effective Hamiltonian $`H_\lambda `$ contains form factors of width $`\lambda `$ in all its vertices. If we denote an operator without the form factors by $`G_\lambda `$, our Hamiltonian takes the form $`H_\lambda =F_\lambda [G_\lambda ]`$, where the operator $`F_\lambda `$ inserts the form factors. With these form factors, momentum transfers in interactions between effective particles are guaranteed to be at most of the order of $`\lambda `$. $`𝒢_\lambda =𝒰_\lambda ^{}G_\lambda 𝒰_\lambda `$. We divide $`𝒢_\lambda `$ into two parts, a part that is bilinear in $`q_{\mathrm{}}`$, and an interaction part that would vanish if the coupling constant were equal 0, so that $`𝒢_\lambda =𝒢_0+𝒢_{I\lambda }`$. The operator $`𝒢_{I\lambda }`$ satisfies the following differential equation as a consequence of Eqs. (3.2)-(3.5),
$$\frac{d}{d\lambda }𝒢_{I\lambda }=[f𝒢_I,\left\{\frac{d}{d\lambda }(1f)𝒢_I\right\}_{𝒢_0}].$$
$`(3.6)`$
We dropped the subscript $`\lambda `$ on the right-hand side for clarity. $`f`$ denotes the similarity form factor introduced by $`F_\lambda `$ and the curly bracket with the subscript $`𝒢_0`$ denotes a solution for $`𝒯_\lambda `$ resulting from Eq. (3.5).
4. ASYMPTOTIC FREEDOM IN SCALAR THEORY
Since the interaction term (1.1) is only a part of the QCD Hamiltonian and the function $`Y_{123}`$ depends on spins and momenta of gluons, let us first discuss the case of scalar field with classical Lagrangian density
$$=\frac{1}{2}(_\mu \varphi ^\mu \varphi \mu ^2\varphi ^2)\frac{g}{3!}\varphi ^3.$$
$`(4.1)`$
In this case, the interaction term in the corresponding Hamiltonian is of the form (1.1), but $`Y_{123}=1/2`$ and calculations are much simpler than in QCD. Our goal is to describe results for the light-front Hamiltonian for effective bosons calculated in perturbation theory up to third power in $`g`$. Although our presentation is based on Ref. that uses plain expansion in powers of $`g`$, the reader may also wish to compare our results with Ref. , where a different scheme is used, including transverse locality and coupling coherence.
The light-front Hamiltonian corresponding to the Lagrangian density (4.1) reads
$$H_{\mathrm{}}=[k]\frac{k^{\mathrm{\hspace{0.17em}2}}+\mu ^2}{k^+}a_\mathrm{}k^{}a_\mathrm{}k+$$
$$+\frac{g}{2}[k_1k_2k_3]\mathrm{\hspace{0.17em}2}(2\pi )^5\delta ^5(k_1+k_2k_3)(a_{\mathrm{}k_1}^{}a_{\mathrm{}k_2}^{}a_{\mathrm{}k_3}+a_{\mathrm{}k_3}^{}a_{\mathrm{}k_2}a_{\mathrm{}k_1})r_\mathrm{\Delta }+X_\mathrm{\Delta },$$
$`(4.2)`$
where $`r_\mathrm{\Delta }`$ is a smooth regularization factor and $`X_\mathrm{\Delta }`$ denotes counterterms (derivable in perturbation theory). In $`n`$ dimensions, $`[k]`$ means $`\theta (k^+)dk^+d^{n2}k^{}/(2k^+(2\pi )^{n1})`$. We choose
$$r_\mathrm{\Delta }=\mathrm{exp}\frac{(\eta _1+\eta _2)\kappa _{12}^{\mathrm{\hspace{0.17em}2}}}{\mathrm{\Delta }^2},$$
$`(4.3)`$
where $`x_1=k_1^+/k_3^+`$ and $`\kappa _{12}^{}=k_1^{}x_1k_3^{}`$, $`\eta _i=\eta (x_i)`$, and $`\eta `$ is a useful function of its argument. A natural choice is $`\eta (x)=1`$, for it is simple. Leaving $`\eta `$ unspecified will help us identify finite regularization effects.
The similarity form factor for an operator containing $`u`$ creation operators and $`v`$ annihilation operators is defined by
$$f_\lambda (u,v)=\mathrm{exp}[(_u^2_v^2)^2/\lambda ^4].$$
$`(4.4)`$
The script notation for invariant masses means $`_u^2=(k_1+\mathrm{}+k_u)^2`$, where the minus components of the momentum four-vectors are given by $`k_i^{}=(k_i^{\mathrm{\hspace{0.17em}2}}+\mu ^2)/k_i^+`$ for $`i=1,\mathrm{},u`$, and similarly for $`v`$.
Equation (3.6) can now be solved order by order using expansion in powers of $`g`$. Firstly, one obtains the counterterms $`X_\mathrm{\Delta }`$ as the initial conditions at $`\lambda =\mathrm{}`$ that render regularization independent finite $`\lambda `$ Hamiltonians. To order $`g^3`$, the regularization dependence of $`H_\lambda `$ lets us identify two counterterms: the mass counterterm
$$\beta _{\mathrm{}11}=[k]\frac{\delta \mu _{\mathrm{}}^2}{k^+}a_\mathrm{}k^{}a_\mathrm{}k,$$
$`(4.5)`$
and the vertex counterterm
$$\gamma _{\mathrm{}21}=[k_1k_2k_3]\mathrm{\hspace{0.17em}2}(2\pi )^5\delta ^5(k_1+k_2k_3)\gamma _{\mathrm{}}(k_1,k_2,k_3)a_{\mathrm{}k_1}^{}a_{\mathrm{}k_2}^{}a_{\mathrm{}k_3}r_\mathrm{\Delta }.$$
$`(4.6)`$
Without loss of generality, we assume that some gedanken experimental data require the mass squared parameter in effective Hamiltonian with $`\lambda =\lambda _0`$ to be equal $`\mu ^2+\delta \mu _0^2`$. This means that when one calculates observables using the effective Hamiltonian, $`\mu _{\lambda _0}^2`$ must equal $`\mu ^2+\delta \mu _0^2`$ to fit the data. This condition, by tracing the renormalization group equation for $`H_\lambda `$ back to $`\lambda =\mathrm{}`$, tells us that
$$\delta \mu _{\mathrm{}}^2=\delta \mu _0^2\left(\frac{g}{2}\right)^2\frac{1}{2(2\pi )^5}_0^1\frac{dx}{x(1x)}d^4\kappa ^{}\frac{2}{^2\mu ^2}\left[f_{\lambda _0}^2(^2,\mu ^2)1\right]r_{\mathrm{\Delta }\beta }.$$
$`(4.7)`$
The script $``$ denotes invariant mass, $`^2=(\kappa ^{\mathrm{\hspace{0.17em}2}}+\mu ^2)/x(1x)`$, and the regularization factor is
$$r_{\mathrm{\Delta }\beta }=\mathrm{exp}\left\{2[\eta (x)+\eta (1x)]\kappa ^{\mathrm{\hspace{0.17em}2}}/\mathrm{\Delta }^2\right\}.$$
$`(4.8)`$
Integration gives two diverging terms, one proportional to $`\mathrm{\Delta }^2`$ and another one proportional to $`\mathrm{log}\mathrm{\Delta }`$. The remaining finite part depends on our choice of the function $`\eta `$. For example, evaluating the integral for $`\eta (x)=1/x`$ one obtains
$$\delta \mu _{\mathrm{}}^2=g^2\frac{1}{(4\pi )^3}\left[\frac{1}{24}\mathrm{\Delta }^2\mu ^2\frac{5}{6}\mathrm{log}\frac{\mathrm{\Delta }}{\mu }+\mu _\eta ^2\right],$$
$`(4.9)`$
where $`\mu _\eta `$ has a finite limit when $`\mathrm{\Delta }\mathrm{}`$. The logarithmically divergent part is independent of the function $`\eta `$ and agrees with results for the Lagrangian mass squared counterterm obtained using Feynman diagrams and dimensional regularization in the following sense: when one changes $`\mathrm{\Delta }`$ to $`\mathrm{\Delta }^{}`$ the logarithmic part of the counterterm changes with $`\mathrm{\Delta }`$ as the mass squared changes as a function of the renormalization scale in Eq. (7.1.22) in .
The vertex counterterm is defined by the requirement that the effective vertex in the Hamiltonian $`H_{\lambda _0}`$ is free from regularization dependence for arbitrary finite values of $`\lambda _0`$. The one loop regularization sensitive contributions to the effective vertex function are given by
$`\gamma _{\mathrm{}}(k_1,k_2,k_3)+\left({\displaystyle \frac{g}{2}}\right)^3{\displaystyle \frac{\pi ^2}{2(2\pi )^5}}\times `$
$`\times [{\displaystyle \frac{1}{2}}[{\displaystyle _{x_1}^1}{\displaystyle \frac{dx}{x(1x)(xx_1)}}{\displaystyle _0^{\mathrm{}}}\kappa ^2d\kappa ^2\mathrm{\hspace{0.17em}\hspace{0.17em}8}{\displaystyle \frac{xx_1}{xx_2^4}}\mathrm{exp}\left({\displaystyle \frac{c_\eta \kappa ^2}{\mathrm{\Delta }^2}}\right)+(x_1x_2)]+`$
$`+{\displaystyle _0^1}{\displaystyle \frac{dx}{x(1x)}}{\displaystyle _0^{\mathrm{}}}\kappa ^2d\kappa ^2{\displaystyle \frac{3}{^4}}\mathrm{exp}\left({\displaystyle \frac{d_\eta \kappa ^2}{\mathrm{\Delta }^2}}\right)],`$
$`(4.10)`$
where
$$c_\eta =\eta (x)+\eta (1x)+\left\{\eta (x_1/x)+\eta [(xx_1)/x]\right\}(x_1/x)^2+\eta [(xx_1)/x_2]+\eta [(1x)/x_2]$$
$`(4.11)`$
and
$$d_\eta =2[\eta (x)+\eta (1x)].$$
$`(4.12)`$
The counterterm function $`\gamma _{\mathrm{}}(k_1,k_2,k_3)`$ must remove the regularization dependence from the above expression. The regularization effects are independent of $`\kappa _{12}^{}`$. Dropping all parts that are independent of regularization, we conclude that
$`\gamma _{\mathrm{}}(k_1,k_2,k_3)+\left({\displaystyle \frac{g}{2}}\right)^3{\displaystyle \frac{1}{(4\pi )^3}}\times `$
$`[3\mathrm{log}{\displaystyle \frac{\mathrm{\Delta }}{\mu }}\mathrm{\hspace{0.17em}4}[{\displaystyle _{x_1}^1}dx{\displaystyle \frac{1x}{x_2}}\mathrm{log}c_\eta +(x_1x_2)]+\mathrm{\hspace{0.17em}3}{\displaystyle _0^1}dxx(1x)\mathrm{log}d_\eta ].`$
$`(4.12)`$
must be independent of regularization. We see that the diverging regularization dependence of the interaction vertex, i.e. the term proportional to $`\mathrm{log}\mathrm{\Delta }`$, is independent of the particle momenta and one can remove the divergence by introducing a $`\gamma _{\mathrm{}}(k_1,k_2,k_3)`$ that is equivalent to changing the initial coupling constant $`g`$ in Eq. (4.2). Thus, no diverging $`x`$-dependent counterterms are required - a different situation than in . However, it is also visible that the vertex contains a finite regularization dependent part that is a function of $`x_1`$. The function depends on our choice for $`\eta `$. For example, if $`\eta =1`$ one has $`c_\eta =4+2(x_1/x)^2`$ and $`d_\eta =4`$. The resulting integral is a function of $`x_1`$, and needs to be subtracted. But this would not assure us that the whole ultraviolet regularization dependence is removed, because we work with a specific functional form of the regulating function (4.3).
Since the whole regularization effect is independent of $`\lambda `$ and $`\kappa _{12}^{}`$, it can be completely removed from the effective interaction by subtracting its value for $`\kappa _{12}^{}=0`$ at an arbitrarily chosen finite $`\lambda _0`$. However, one has to add back the finite regularization independent part of the effective vertex, which is a function of $`x_1`$, denoted below by $`\gamma _0(x_1)`$. The function $`\gamma _0(x_1)`$ is necessary to recover Poincaré symmetry of observables, because our regularization spoiled the symmetry. The symmetry may be restored once counterterms remove the regularization effects, but one is not allowed to change terms independent of the regularization, which were given by the initial covariant Lagrangian density unambiguously. Therefore, the function $`\gamma _0(x_1)`$ must be reinserted. This function is not altered when $`\lambda `$ changes and could be considered marginal in analogy with usual renormalization group analysis. The ultimate adjustment of the function $`\gamma _0(x_1)`$ requires 4th order calculations. For there exists in $`\varphi ^3`$ theory no 3rd order scattering amplitude one could use to find out what function $`\gamma _0(x_1)`$ renders Poincaré symmetry of scattering observables with our choice of $`r_\mathrm{\Delta }`$ in Eq. (4.2). However, it should be pointed out that the function does not influence the way the 3rd order running coupling constant in effective Hamiltonians depends on $`\lambda `$.
Thus, in Eq. (4.6), the counterterm function $`\gamma _{\mathrm{}}(k_1,k_2,k_3)\gamma _{\mathrm{}}(x_1,\kappa _{12}^{})`$, which removes the regularization dependence from the effective vertex reads
$$\gamma _{\mathrm{}}(x_1,\kappa _{12}^{})=\gamma _{\lambda _0}(x_1,0^{})+\gamma _0(x_1).$$
$`(4.13)`$
This result is used to define the new regularization dependent coupling constant $`g_\mathrm{\Delta }`$ in the initial Hamiltonian in Eq. (4.2). We select a convenient value of $`x_1=x_0`$ and obtain
$$\frac{g_\mathrm{\Delta }}{2}=\frac{g}{2}+\gamma _{\mathrm{}}(x_0,0^{})=\frac{g}{2}\gamma _{\lambda _0}(x_0,0^{})+\gamma _0,$$
$`(4.14)`$
where $`\gamma _0\gamma _0(x_0)`$. We see that the initial coupling $`g`$ is replaced by a new $`\mathrm{\Delta }`$-dependent quantity
$$g_\mathrm{\Delta }=g\left[1g^2\frac{3}{4(4\pi )^3}\mathrm{log}\frac{\mathrm{\Delta }}{m_0}\right]+o(g^5),$$
$`(4.15)`$
with certain free constant $`m_0`$. Thus, the theory exhibits asymptotic freedom in 3rd order terms. Our result agrees with literature, say Eq. (7.1.26) from , in the sense that when we change $`\mathrm{\Delta }`$, the change required in the coupling constant in the initial Hamiltonian for obtaining $`\mathrm{\Delta }`$-independent effective Hamiltonians matches the change implied by Feynman diagrams and dimensional regularization.
Having established the structure of counterterms we can proceed to evaluation of the finite similarity flow of effective Hamiltonians towards small widths $`\lambda `$. The effective kinetic energy term in narrow Hamiltonians is
$$H_{\lambda 11}=[k]\frac{k^{\mathrm{\hspace{0.17em}2}}+\mu _\lambda ^2}{k^+}a_{\lambda k}^{}a_{\lambda k},$$
$`(4.16)`$
where
$`\mu _\lambda ^2=\mu ^2+\delta \mu _\lambda ^2=`$
$`=\mu ^2+\delta \mu _0^2+\left({\displaystyle \frac{g}{2}}\right)^2{\displaystyle \frac{1}{2(2\pi )^5}}{\displaystyle _0^1}{\displaystyle \frac{dx}{x(1x)}}{\displaystyle d^4\kappa ^{}\frac{2}{^2\mu ^2}\left[f_\lambda ^2(^2,\mu ^2)f_{\lambda _0}^2(^2,\mu ^2)\right]}.`$
$`(4.17)`$
The above result is particularly simple for $`\mu =0`$ and in that case it reads ($`\delta \mu _0^2`$ is proportional to $`g^2`$)
$$\mu _\lambda ^2=\delta \mu _0^2+g^2\frac{1}{(4\pi )^3}\frac{1}{24}\sqrt{\frac{\pi }{2}}(\lambda ^2\lambda _0^2).$$
$`(4.18)`$
Logarithmic dependence on $`\lambda `$ arises for $`\mu >0`$. The value of $`\delta \mu _0^2`$ could be found, for example, by solving a single physical boson eigenvalue problem, expressing the physical boson mass in terms of $`\delta \mu _0^2`$ and adjusting the latter to obtain the gedanken experimental mass value for bosons. Note a change in the mass function of cutoff parameter, from the case of the mass counterterm, dependent on $`\mathrm{\Delta }`$, to the case of running mass term, dependent on the width $`\lambda `$ (independent of $`\mathrm{\Delta }`$). The change corresponds to a transition from the initial side of a fixed point (bare canonical Hamiltonian with regularization) to the other side of the fixed point (renormalization group trajectory of effective Hamiltonians in the similarity procedure, cf. ).
The effective vertex reads
$$H_{\lambda 21}=[k_1k_2k_3]\mathrm{\hspace{0.17em}2}(2\pi )^5\delta ^5(k_1+k_2k_3)f_\lambda [(k_1+k_2)^2,k_3^2]V_\lambda (x_1,\kappa _{12}^{})a_{\lambda k_1}^{}a_{\lambda k_2}^{}a_{\lambda k_3},$$
$`(4.19)`$
where $`V_\lambda (x_1,\kappa _{12}^{})`$ is the effective vertex function and $`f_\lambda `$ is the similarity vertex form factor. The vertex function is given by an integral over loop variables $`x`$ and $`\kappa ^{}`$ of a known function.
We define the running coupling constant as the value of $`2V_\lambda (x_1,\kappa _{12}^{})`$ at a chosen configuration of momentum variables, denoted by $`(x_{10},\kappa _{120}^{})`$. In other words, $`g_\lambda =2V_\lambda (x_{10},\kappa _{120}^{})`$. A possible choice for massless bosons is $`x_{10}=0`$ and $`\kappa _{120}^{}=0`$. This is a natural definition, analogous to the standard Thomson limit in the case of electron charge in QED. This choice greatly simplifies the integrand, giving $`V_\lambda (0,0^{})`$, so that the result can be fully produced here ($`g_0`$ is the value of $`g_{\lambda _0}`$ required by phenomenology done using $`H_{\lambda _0}`$)
$`g_\lambda =g_0+g_0^3{\displaystyle \frac{1}{24}}{\displaystyle \frac{1}{(4\pi )^3}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dz}{z}}`$
$`\left[2(f_\lambda f_\lambda ^3)2(f_0f_0^3)+20(f_\lambda ^3f_\lambda ^2)20(f_0^3f_0^2)+9(f_0^2f_\lambda ^2)\right],`$
$`(4.20)`$
where $`f_\lambda =\mathrm{exp}z^2/\lambda ^4`$ and $`f_0=\mathrm{exp}z^2/\lambda _0^4`$. A straightforward integration gives
$$g_\lambda =g_0g_0^3\frac{3}{4(4\pi )^3}\mathrm{log}\frac{\lambda }{\lambda _0},$$
$`(4.21)`$
which exhibits asymptotic freedom. Differentiating with respect to $`\lambda `$ and keeping terms up to order $`g_\lambda ^3`$ one obtains
$$\frac{d}{d\lambda }g_\lambda =g_\lambda ^3\frac{3}{256\pi ^3}\frac{1}{\lambda }.$$
$`(4.22)`$
This equation demonstrates the same $`\beta `$ function for coupling constants in effective Hamiltonians as obtained in Lagrangian approaches using Feynman diagrams and dimensional regularization, when one identifies the renormalization scale with the Hamiltonian width $`\lambda `$. This is encouraging but one needs to remember that for comparison of perturbative scattering amplitudes in Hamiltonian and Lagrangian approaches it is necessary to make additional calculations and at least of fourth order in $`g`$. Beyond model matrix studies such as in , 4th order similarity calculations have so far been carried out only in a simplified Yukawa model by Masłowski and Wiȩckowski (the latter model calculations should be helpful in setting up a light-front theory of nucleons and pions).
Integrating Eq. (4.10), one obtains ($`\alpha =g^2/4\pi `$)
$$\alpha _\lambda =\frac{\alpha _0}{1+\alpha _0(3/32\pi ^2)\mathrm{log}\lambda /\lambda _0},$$
$`(4.23)`$
which shows our result for a boost invariant running coupling constant in effective Hamiltonians. Our procedure explains how the running coupling constant can be included in quantum mechanics of effective particles, which is given by the Schrödinger equation with the corresponding Hamiltonian $`H_\lambda `$. One can evaluate matrix elements of the small width Hamiltonian in a limited subspace of Fock states built of the effective particles. The form factor in the interaction vertex (4.19) secures a small range of the interactions on the energy scale and one can expect a rapid convergence of wave functions in the effective particle basis.
Note three characteristic features of the Hamiltonian calculation. (1) No field renormalization constant appeared, since the similarity transformation did not eliminate (or integrated out) any degrees of freedom. (2) No vacuum effect played any role, since we used the light-front form of dynamics. Extensive literature concerning the vacuum issue can be traced through reference . $`\varphi ^3`$ theory is unstable due to a possibility that the field $`\varphi `$ takes an infinitely large, negative value. It would be interesting to check if the perturbatively evaluated effective Hamiltonians of small widths have any tendency to develop eigenstates that deviate in that direction. (3) The Hamiltonian structure is invariant with respect to boosts, including boosts from the rest frame of any bound state to the infinite momentum frame. This suggests the approach outlined above should be tried in QCD, and in effective theories of strong interactions in nuclear physics, to connect low energy observables, such as binding energies, radii or magnetic moments of bound states, with high energy ones, such as parton distributions, form factors or jets.
5. QCD GLUON VERTEX COUNTERTERM
We come back to Eq. (1.1) in QCD and repeat the same analysis as we did for the scalar theory. Most of the procedure remains the same, but an important complication arises. The vertex function in the canonical Hamiltonian has now the form ($`c`$ refers to color and $`ϵ`$ to polarization of gluons)
$$Y_{123}=if^{c_1c_2c_3}[ϵ_1^{}ϵ_2^{}ϵ_3^{}\kappa _{12}^{}ϵ_1^{}ϵ_3^{}ϵ_2^{}\kappa _{12}^{}/x_2ϵ_2^{}ϵ_3^{}ϵ_1^{}\kappa _{12}^{}/x_1],$$
$`(5.1)`$
in which the characteristic factors of $`\kappa ^{}/x`$ tend to infinity when $`x0`$. Ultraviolet coupling constant divergences in $`\varphi ^3`$ theory in 6 dimensions resulted from transverse momentum integration $`d^4\kappa ^{}/\kappa ^4`$, where $`1/\kappa ^4`$ came from the two denominators of third order perturbation theory. In QCD in 4 dimensions, we have instead $`d^2\kappa ^{}(\kappa ^{}/x)^2/\kappa ^4`$. Therefore, in QCD (more generally, in gauge theories), we have to introduce a separate regularization of small $`x`$ behavior of interaction vertices in the Hamiltonians. For example, in the QCD counterpart of Eq. (4.2), we have to insert a factor, denoted by $`r_\delta `$, in addition to $`r_\mathrm{\Delta }`$, that will effectively cut-off the region where one of the gluons 1 or 2 carries a smaller fraction of $`k_3^+`$ than the size of a small parameter $`\delta `$.
The small $`x`$ regularization function $`r_\delta `$ may appear to be only a technical detail. But it was pointed out by Perry that singularities at small $`k^+`$ may be related to effective confining potentials in $`H_\lambda `$, calculable already in second order perturbation theory. In short, the canonical light-front QCD Hamiltonian contains terms that are singular at small $`k^+`$ and the singularity contributes to the effective Hamiltonians $`H_\lambda `$, providing potentials that grow with distance between color charges. This is quite different a situation from other formulations of the theory, where second order calculations are not expected to tell us anything about confinement. Therefore, the small $`x`$ features of QCD in the light-front Hamiltonian approach deserve extensive studies. Here, we merely report some initial results for third order gluon vertex counterterm, indicating $`x`$-dependent features.
The whole analysis of the previous Section can be repeated step by step and one can derive the interaction term for effective gluons, in the narrow Hamiltonian $`H_\lambda `$ for QCD. The condition that the effective vertex is independent of regularization parameter $`\mathrm{\Delta }`$ gives us the diverging triple-gluon vertex counterterm in the initial QCD Hamiltonian. The new features appear in the coefficient of $`\mathrm{log}\mathrm{\Delta }`$. Namely, the diverging part of the vertex, to be compensated by the ultraviolet counterterm, has the form
$$\frac{g^3}{4\pi ^2}\left[\frac{11}{12}N_c\frac{1}{6}n_f+N_cf(x_1,x_2)\right]\mathrm{log}\mathrm{\Delta }Y_{123},$$
$`(5.2)`$
where the function $`f(x_1,x_2)`$ is symmetric in its arguments. This function originates from three successive actions of the triple gluon interaction, the gluon mass correction providing only a constant contribution, and it depends on the regularization factor $`r_\delta `$ in the initial Hamiltonian. In fact, $`f(x_1,x_2)=[\mathrm{log}x_1+_0^1dx(2/x+1/(1x))(r_{\delta 3}r_{\delta 2})+(12)]/2`$, where $`r_{\delta 3}=r_{\delta 3}(x,x_1)`$ is a product of three factors $`r_\delta `$ for the three successive triple-gluon interactions, and $`r_{\delta 2}=r_\delta ^2(x)`$, is a product of two factors $`r_\delta `$ from the gluon mass counterterm. For $`r_\delta (x)=\theta (x\delta )\theta (1x\delta )`$ in the bare vertex, for gluons carrying $`x`$ and $`1x`$ of the single gluon momentum $`k_3^+`$ in Eq. (1.1), $`f(x_1,x_2)=\mathrm{log}[min(x_1,x_2)]`$, which is negative and hence reduces the rate at which the initial coupling depends on $`\mathrm{\Delta }`$. The latter feature comes about as follows. Since the counterterm must contain the term opposite in sign to (5.2), the coupling constant in the regularized $`H_{\mathrm{}}(q_{\mathrm{}})`$ in QCD is changed to \[cf. Eq. (4.15)\]
$$g_\mathrm{\Delta }=g\frac{g^3}{4\pi ^2}\left[\frac{11}{12}N_c\frac{1}{6}n_f+N_cf(x_0,1x_0)\right]\mathrm{log}\mathrm{\Delta }.$$
$`(5.3)`$
The coefficient of $`\mathrm{log}\mathrm{\Delta }`$ is independent of the parameter $`\delta `$, but it depends on the small $`x`$ regularization in a finite way. Having derived the function $`f(x_1,x_2)`$, for certain choices of $`x_1=x_0`$ in the definition of the Hamiltonian coupling constant, one could obtain triviality instead of asymptotic freedom. This is a peculiar result for the regularization given in Eq. (4.3) and the sharp cutoff on $`x`$ at $`\delta `$. In the limit when both $`r_{\delta 3}`$ and $`r_{\delta 2}`$ are replaced by 1, one would obtain $`f(x_1,x_2)=\mathrm{log}\sqrt{x_1x_2}`$, which is always positive, and would accelerate the asymptotic freedom rate of change of $`g_\mathrm{\Delta }`$ with $`\mathrm{\Delta }`$. One can seek choices of regulating functions $`r_\mathrm{\Delta }`$ and $`r_\delta `$ that eliminate the unusual logarithm (the other terms are standard) but it is not known if a finite function of $`x_1`$ is not necessary in place of $`f(x_1,x_2)`$ in light-front Hamiltonians anyway, to restore symmetries for physical quantities. In addition, as in the scalar theory, the ultraviolet finite part of the counterterm involves an unknown function of $`x_1`$. One could say that the structure of the model from Ref. is closer to QCD than to scalar theory in 6 dimensions. Evaluation of the effective coupling constant $`g_\lambda `$ in QCD, in analogy to Eq. (4.20) in scalar theory, may shed some light on how to disentangle genuine ultraviolet from small $`x`$ singularities in Hamiltonians.
It is clear from the above example that effective light-front Hamiltonians of QCD require careful studies employing various types of regulators before we will know the optimal ways of calculating window Hamiltonians. The interplay of transverse and longitudinal momentum variables may lead to surprising results. However, the calculations are certainly doable and the resulting matrices will tell us about details of QCD dynamics in the Fock space of effective quarks and gluons. The similarity renormalization group procedure for Hamiltonians is able to reveal new features of effective particle dynamics which standard Lagrangian approaches do not reveal.
6. TRANSITION TO NEW DEGREES OF FREEDOM
We may hope to make a transition from the effective QCD degrees of freedom to nuclear physics hadronic interactions, such as pion-nucleon coupling, after we achieve understanding of the narrow $`H_\lambda `$ in QCD. To see the basis for this hope, let us come back again to the matrix picture.
Once we have the narrow Hamiltonian matrix, we can divide it into the boxes as it is illustrated in the figure below, neglecting the small triangles in the band outside of them, which are initially left out. We can find eigenstates of the boxes (they correspond to different invariant mass states) and calculate the band diagonal matrix matrix elements in the basis built from those eigenstates. The states corresponding to the middle energy scale of each box will not interact very strongly with neighboring (in energy, or mass) states, but matrix elements sensitive to the left-out triangles will lead to strong interactions. We can imagine that the lowest box corresponds to nucleons interacting through potential forces, the next box corresponds to nucleons plus one meson, the second box to nucleons and two mesons, etc. This is how one could make a connection between the QCD dynamics and nuclear physics through similarity renormalization group for Hamiltonians (cf. ).
Changing degrees of freedom
However, this is not the only possibility one can try to explore for the change of basis. One can consider new basis states built of quarks and gluons, possibly open gluon string bits with quarks at the ends or closed rings of gluons, and evaluate matrix elements of the effective Hamiltonians between such objects. Since one has a perturbative expression for the operators $`𝒰_\lambda `$, see Eqs. (3.1) and (3.2), one can attempt evaluation of matrix elements between states that are constructed in a variety of ways, using quarks and gluons corresponding to different scales $`\lambda `$. One could even ask if there is a way to calculate a connection between the quark and gluon matrices of intermediate widths and reggeized gluon interactions, once one restricts the space of states to those that dominate in multi-Regge kinematics.
7. CONCLUSION
Asymptotically free theories can be analyzed using Hamiltonian approach. The analysis can be based on the similarity renormalization group procedure for effective particles. The evaluation of running couplings in the effective Hamiltonians can be carried out without introduction of wave function renormalization constants and without invoking any properties of the vacuum state (in the light-front form of Hamiltonian dynamics). In third order calculations, one obtains familiar asymptotic results in scalar $`\varphi ^3`$ theory, plus an $`x`$-dependent finite counterterm. In QCD, the standard asymptotic freedom form of triple-gluon vertex counterterm is supplemented by an ultraviolet diverging and $`x`$-dependent counterterm, and by an ultraviolet finite $`x`$-dependent counterterm. Effects predicted for QCD by the power counting in $`k^{}`$ and $`k^+`$ are confirmed but the analysis is changed by transition to boost invariant variables $`\kappa ^{}`$ and $`x`$, and detailed calculations may produce results that are not expected to emerge from Feynman diagrams. Mixing between the small $`x`$ and large $`\kappa ^{}`$ cutoffs indicates a need for a new precise definition of the ultraviolet domain in the Hamiltonian approach. Nevertheless, one obtains well-defined expressions for effective Hamiltonian interactions without necessity to calculate scattering matrix elements for quarks and gluons as if they were observable particles.
The effective particle calculus preserves cluster properties and allows for evaluation of effective Hamiltonians without limitation to any particular set of matrix elements. In other words, we can derive integral expressions for matrix elements of effective Hamiltonians in the whole Fock space spanned by basis states of effective particles. The renormalization group equations are integrated analytically using gaussian similarity form factors and one fully controls off-shell behavior of effective vertices that correspond to the initial theory. The effective dynamics is invariant with respect to boosts and allows simultaneous analysis of the rest frame and infinite momentum frame structure of bound states.
The effective particle Fock space expansion can converge thanks to the similarity form factors in the interaction vertices. The form factors dampen interactions changing invariant masses by more than $`\lambda `$ and thus can tame the spread of eigenstate wave functions for low lying eigenvalues into regions of high relative momenta of constituents. This feature may lead to exponential convergence of the eigenstate expansion in the effective particle basis. Such convergence is not expected in the case of bare particles. The fine structure of effective particles would then unfold in the transformation $`𝒰_{\lambda _1}𝒰_{\lambda _2}^{}`$ relating effective degrees of freedom at two different scales, one corresponding to the binding scale and the other to the high momentum transfer probe in question.
The near-diagonal Hamiltonian matrices allow for transition to new degrees of freedom by turning to basis states that are eigenstates of small block matrices on the diagonal. In principle, these new degrees of freedom could correspond to mesons and baryons built from constituent quarks and gluons, in case of QCD. One can also consider other changes of basis states and seek most efficient degrees of freedom, such as strings of gluons with quarks at the ends, for solving the nonperturbative eigenvalue problems for narrow effective Hamiltonians.
Acknowledgments
The author thanks Lev Lipatov for discussions about $`x`$-dependent effects and Ken Wilson for comments and discussion about the manuscript. Support and hospitality from Andreas Schreiber, Tony Thomas, Tony Williams, and Frederic Bonnet during the author’s stay at CSSM, are gratefully acknowledged. |
warning/0001/hep-th0001180.html | ar5iv | text | # Quantum Coherent String States in 𝐴𝑑𝑆₃ and 𝑆𝐿(2,𝑅) WZWN Model
## 1 Introduction
The question about the spectrum of bosonic string theory in 3-dimensional Anti de Sitter space, $`AdS_3SL(2,R)SU(1,1)`$, attracted a lot of interest about 10 years ago \[1-10\], as a first example of exact string quantization on a manifold with curved space and curved time. It was immediately realized that the problem of unitarity was much more complicated than in flat Minkowski spacetime , in the sense that the Virasoro constraints for $`AdS_3`$ themselves did not eliminate all negative norm states . However, in a series of papers, going back to (see also \[6-10, 13, 14\]), it has been argued that unitarity can be ensured (at least for free strings) by imposing certain restrictions on the allowed representations, exemplified by the now well-known spin-level restriction for the discrete representations $`|j|<k/2`$, where $`j`$ is the spin and $`k`$ is the level of the $`SL(2,R)`$ WZWN model. For a somewhat different approach towards unitarity, see .
More recently, the interest in $`AdS_3`$ (as well as higher dimensional $`AdS`$ spaces) has increased in connection with the conjecture relating supergravity and superstring theory on $`AdS`$ space with a conformal field theory on the boundary. In such constructions, $`AdS_3`$ often appears on the 10-dimensional supergravity/superstring side in a cartesian product with some other compact spaces, for instance as $`AdS_3\times S^3\times T^4`$. Thus, again it has become extremely important to understand the precise spectrum of string theory in $`AdS_3`$.
Even if the problem of unitarity appearently could be solved by the spin-level condition (although this question is certainly not completely settled yet), several problems remained. One of the most important being that the spin-level restriction together with the mass-shell condition imposes a restriction on the grade (to avoid confusion with the level $`k`$, we use the word ”grade” for what is usually called level). This restriction on the grade means that for fixed level $`k`$, a string living in $`AdS_3`$ can only be excited to its very lowest grades. In other words, we are faced with the problem that it seems impossible to have very massive strings in $`AdS_3`$.
On the other hand, the dynamics of classical strings and their semi-classical quantization in $`AdS_3`$ \[18-22\] does not seem to indicate any particular problems for very long and very massive strings, although the question of unitarity cannot really be addressed exactly in such studies. It is therefore highly interesting and important to understand how such long massive strings can arise in an exact quantization scheme for strings in $`AdS_3`$, without being in conflict with unitarity.
In this paper we suggest that long massive strings can be described as coherent string states based on one of the standard discrete representations of $`SL(2,R)`$. For simplicity and clarity, we shall construct quantum coherent string states corresponding, in the classical limit, to the circular strings discussed in , but our construction can be used for other string configurations too.
As for (most) other families of string states in $`AdS_3`$, coherent string states generally do not have positive norm, even if they fulfil the Virasoro conditions. The condition of finite positive norm for the coherent states gives rise to certain restrictions on the spin $`j`$, which in turn restricts the mass of the states. We show that the finite positive norm condition for our coherent string states leads to a mass-spectrum consisting of two parts: A continuous spectrum of low mass states (partly tachyonic) where $`j`$ fulfils the standard spin-level restriction, as well as an infinite tower of discrete high mass states for which the mass-formula is given by eq.(5.20) and asymptotically is $`m^2\alpha ^{}N^2`$ ($`N`$ integer). This result agrees, to leading order, with what was found using semi-classical quantization .
When completing this paper, a recent preprint by Maldacena and Ooguri appeared , considering a similar problem. Our construction, however, is completely different from theirs. First of all, their massive strings are based on descendents of primary states for a new set of $`SL(2,R)`$ representations obtained from the standard ones by a ”spectral flow” operation , whereas our massive strings are based on coherent states of descendents of primary states from the standard $`SL(2,R)`$ representations. Secondly, their construction relies heavily on the existence of some internal compact manifold $``$, which is assumed to give a large contribution to the total world-sheet energy-momentum tensor, whereas our construction works directy for $`AdS_3`$ without need of any additional internal compact manifold (although of course we could easily include an internal compact manifold as well). In fact, with the internal compact manifold, the construction of gives only a finite number of very massive states, and without the internal manifold it gives at most a few very massive states (or even none, depending on some other parameters). Our construction gives in any case an infinite tower of more and more massive states in the $`AdS_3`$ and $`SL(2,R)`$ background. This is again in agreement with the previous semi-classical quantization results giving an infinite number of string states in the $`AdS_3`$ and $`SL(2,R)`$ background. Yet another difference between the two approaches has to do with the world-sheet energy $`L_0`$. In the construction of , $`L_0`$ is not bounded from below for the representations obtained by the spectral flow, while in our construction $`L_0`$ is bounded from below since we are using the standard representations. However, we are not working with eigenstates of $`L_0`$ thus the standard mass-shell condition $`(L_01)|\psi >=0`$ is replaced by $`<\psi |(L_01)|\psi >=0`$. In any case, the mass-shell condition eventually selects those states with ”$`L_0=1`$”.
It must be noticed however, that in despite of the differences between the construction of and our construction, it turns out that the final results for the mass-spectrum (at least when an internal compact manifold with a large contribution to the world-sheet energy-momentum tensor is assumed) are more or less identical; namely, a low mass continuous spectrum and a high mass discrete spectrum, where the energy (or mass) to leading order grows linearly with an integer.
Interestingly enough, our quantum coherent states are a string generalization of the ordinary coherent states of quantum mechanics. All the usual properties of ordinary coherent states are obtained in the $`k\mathrm{}`$ limit. For instance, the low mass continuous spectrum of string states become the ordinary coherent states, eigenstates of the annihilation operator, for any value of the spin $`j1/2`$, while the high mass discrete spectrum of string states completely disappears, pushed towards infinite mass. These are precisely the properties which in quantum mechanics characterise coherent states as quasi-classical, being the states for which quantum uncertainty is minimal.
Our paper is organized as follows. In Section 2, we review the classical $`SL(2,R)`$ WZWN model, mainly to set our conventions and normalizations. We also give a simple derivation of the reduction of the classical equations of motion to the Liouville equation . In Section 3, we reconsider the classical oscillating circular strings in terms of $`SL(2,R)`$ currents. In Section 4, we present the standard results of the quantization of conformal field theories on a group manifold ; we only give the results which we will use later. In Section 5, we then turn to the construction of the quantum coherent string states. We derive the expression for the norm of such states and show that the condition of finite positive norm leads to a mass-spectrum as explained above. We also show that our coherent string states lead to non-vanishing expectation values only for the components of the currents corresponding to the classical oscillating circular strings. Finally in Section 6, we have some concluding remarks.
## 2 SL(2,R) WZWN Model. The Classical Picture.
Our starting point is the sigma-model action including the WZWN term at level $`k`$
$$S=\frac{k}{8\pi }_M𝑑\tau 𝑑\sigma \eta ^{\alpha \beta }\text{Tr}[g^1_\alpha gg^1_\beta g]\frac{k}{12\pi }_B\text{Tr}[g^1dgg^1dgg^1dg]$$
(2.1)
Here $`M`$ is the boundary of the manifold $`B`$, and $`g`$ is a group-element of the group under consideration (later taken to be $`SL(2,R)`$). The classical string equations of motion are
$$_{}(g^1_+g)=0$$
(2.2)
where we introduced world-sheet light-cone coordinates $`\sigma ^\pm =\tau \pm \sigma `$. The world-sheet energy-momentum tensor is
$$T_{\pm \pm }=\frac{2}{k}\text{Tr}(J_\pm J_\pm )$$
(2.3)
where the conserved currents, $`_\pm J_{}=0`$, are given by
$$J_+=ikg^1(_+g),J_{}=ik(_{}g)g^1$$
(2.4)
and the string constraints are
$$\text{Tr}[(g^1_\pm g)(g^1_\pm g)]=0$$
(2.5)
Equation (2.2) is trivially solved by
$$g(\sigma ^+,\sigma ^{})=g_R(\sigma ^{})g_L(\sigma ^+)$$
(2.6)
It follows that
$$J_+=ikg_L^1(_+g_L),J_{}=ik(_{}g_R)g_R^1$$
(2.7)
and the constraints, eq.(2.5), separate
$$\text{Tr}[(g_L^1_+g_L)^2]=\text{Tr}[(g_R^1_{}g_R)^2]=0$$
(2.8)
In the case of $`SL(2,R)`$, the group elements are given by
$$g_L(\sigma ^+)=\left(\begin{array}{cc}\stackrel{~}{a}(\sigma ^+)& \stackrel{~}{u}(\sigma ^+)\\ \stackrel{~}{v}(\sigma ^+)& \stackrel{~}{b}(\sigma ^+)\end{array}\right),g_R(\sigma ^{})=\left(\begin{array}{cc}a(\sigma ^{})& u(\sigma ^{})\\ v(\sigma ^{})& b(\sigma ^{})\end{array}\right)$$
(2.9)
subject to the normalization conditions
$$\stackrel{~}{a}(\sigma ^+)\stackrel{~}{b}(\sigma ^+)+\stackrel{~}{u}(\sigma ^+)\stackrel{~}{v}(\sigma ^+)=a(\sigma ^{})b(\sigma ^{})+u(\sigma ^{})v(\sigma ^{})=1$$
(2.10)
Then the constraints, eqs.(2.8), are simply (from now on we do not write explicitly the arguments $`(\sigma ^\pm )`$ of the functions)
$$\stackrel{~}{a}_+\stackrel{~}{b}_++\stackrel{~}{u}_+\stackrel{~}{v}_+=a_{}b_{}+u_{}v_{}=0$$
(2.11)
where we introduced the notation $`a_{}=_{}a`$, $`\stackrel{~}{a}_+=_+\stackrel{~}{a}`$, etc.
As for the currents, it is convenient to make a Pauli decomposition
$$J_\pm =\eta _{ab}J_\pm ^at^b$$
(2.12)
in terms of the Pauli matrices,
$$t^1=\frac{i}{2}\sigma ^1,t^2=\frac{i}{2}\sigma ^3,t^3=\frac{1}{2}\sigma ^2$$
(2.13)
such that
$$\text{Tr}(t^at^b)=\frac{1}{2}\eta ^{ab},[t^a,t^b]=iϵ^{abc}t_c$$
(2.14)
(Our conventions are: $`\eta ^{ab}=\text{diag}(1,1,1)`$ and $`ϵ^{123}=+1`$).
It is also standard to introduce
$$J_{}^\pm =J_{}^1\pm iJ_{}^2,J_+^\pm =J_+^1\pm iJ_+^2$$
(2.15)
It is now straightforward to write down explicit expressions for the currents in terms of the group elements, eqs.(2.9),
$`J_{}^\pm `$ $`=`$ $`k\left([au_{}ua_{}+vb_{}bv_{}]\pm 2i[ab_{}+uv_{}]\right)`$
$`J_{}^3`$ $`=`$ $`k[vb_{}bv_{}+ua_{}au_{}]`$
$`J_+^\pm `$ $`=`$ $`k\left([\stackrel{~}{b}\stackrel{~}{u}_+\stackrel{~}{u}\stackrel{~}{b}_++\stackrel{~}{v}\stackrel{~}{a}_+\stackrel{~}{a}\stackrel{~}{v}_+]\pm 2i[\stackrel{~}{v}\stackrel{~}{u}_++\stackrel{~}{a}\stackrel{~}{b}_+]\right)`$
$`J_+^3`$ $`=`$ $`k[\stackrel{~}{v}\stackrel{~}{a}_+\stackrel{~}{a}\stackrel{~}{v}_++\stackrel{~}{u}\stackrel{~}{b}_+\stackrel{~}{b}\stackrel{~}{u}_+]`$ (2.16)
Notice also that
$$T_{\pm \pm }=\frac{1}{k}(J_\pm ^+J_\pm ^{}J_\pm ^3J_\pm ^3)$$
(2.17)
such that the conditions $`T_{\pm \pm }=0`$ again lead to eqs.(2.11), as they should.
We close this section with a few comments about the invariant string size and the reduction of the classical equations of motion to the Liouville equation (for a review of different methods, see Ref .).
The line-element on the group manifold is given by
$$dS^2=\frac{1}{H^2}\text{Tr}[(g^1dg)^2]$$
(2.18)
where $`H^1`$ is the length-scale, which up to a numerical factor is related to $`k`$ and $`\alpha ^{}`$ by
$$k=\frac{1}{H^2\alpha ^{}}$$
(2.19)
where $`\alpha ^{}`$ is related to the string tension $`T`$ in the usual way, $`T=(2\pi \alpha ^{})^1`$.
The line-element on the group manifold induces the following proper line-element on the string world-sheet
$$ds^2=\frac{e^\alpha }{2H^2}d\sigma ^+d\sigma ^{}$$
(2.20)
Here $`\alpha =\alpha (\sigma ^+,\sigma ^{})`$ is the fundamental quadratic form, which determines the invariant string size, and is defined by
$$e^\alpha \text{Tr}[(g^1_+g)(g^1_{}g)]$$
(2.21)
In the case of $`SL(2,R)`$, one finds
$`e^\alpha `$ $`=`$ $`[av_{}va_{}][\stackrel{~}{a}\stackrel{~}{u}_+\stackrel{~}{u}\stackrel{~}{a}_+]+[ub_{}bu_{}][\stackrel{~}{v}\stackrel{~}{b}_+\stackrel{~}{b}\stackrel{~}{v}_+]`$ (2.22)
$`+`$ $`2[ab_{}+vu_{}][\stackrel{~}{b}\stackrel{~}{a}_++\stackrel{~}{v}\stackrel{~}{u}_+]`$
By differentiating this identity twice and by using some simple algebra, we get the equation
$$\alpha _+=2f(\sigma ^{})\stackrel{~}{g}(\sigma ^+)e^\alpha $$
(2.23)
where the functions $`f=f(\sigma ^{})`$ and $`\stackrel{~}{g}=\stackrel{~}{g}(\sigma ^+)`$ are given by
$`f`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{u_{}}{a_{}}}[av_{}va_{}]{\displaystyle \frac{v_{}}{b_{}}}[ub_{}bu_{}]\right)`$ (2.24)
$`\stackrel{~}{g}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{\stackrel{~}{u}_+}{\stackrel{~}{b}_+}}[\stackrel{~}{v}\stackrel{~}{b}_{++}\stackrel{~}{b}\stackrel{~}{v}_{++}]{\displaystyle \frac{\stackrel{~}{v}_+}{\stackrel{~}{a}_+}}[\stackrel{~}{a}\stackrel{~}{u}_{++}\stackrel{~}{u}\stackrel{~}{a}_{++}]\right)`$ (2.25)
The product $`f(\sigma ^{})\stackrel{~}{g}(\sigma ^+)`$ in eq.(2.23) can be absorbed by a conformal transformation, thus we conclude that the most general equation fulfilled by the fundamental quadratic form $`\alpha `$ is
$$\alpha _++Ke^\alpha =0,$$
(2.26)
where:
$$K=\{\begin{array}{c}+1,f(\sigma ^{})\stackrel{~}{g}(\sigma ^+)<0\hfill \\ 1,f(\sigma ^{})\stackrel{~}{g}(\sigma ^+)>0\hfill \\ 0,f(\sigma ^{})\stackrel{~}{g}(\sigma ^+)=0\hfill \end{array}$$
(2.27)
Equation (2.26) is either the Liouville equation ($`K=\pm 1`$), or the free wave equation $`(K=0)`$. This result was obtained in a different way and discussed in detail in Ref. .
## 3 Circular Strings
Circular strings on the $`SL(2,R)`$ group manifold were considered in detail in Ref. . In this section we translate the results into the formalism of $`SL(2,R)`$ currents.
Circular strings are most easily discussed using a different parametrization of $`SL(2,R)`$, corresponding to the static coordinates for Anti de Sitter spacetime. We first write the $`SL(2,R)`$ group-element in the form
$$g=g_Rg_L=\left(\begin{array}{cc}A& U\\ V& B\end{array}\right)=\left(\begin{array}{cc}a\stackrel{~}{a}u\stackrel{~}{v}& a\stackrel{~}{u}+u\stackrel{~}{b}\\ v\stackrel{~}{a}b\stackrel{~}{v}& b\stackrel{~}{b}v\stackrel{~}{u}\end{array}\right)$$
(3.1)
and then introduce coordinates $`(t,r,\varphi )`$ by
$`A`$ $`=`$ $`\sqrt{1+H^2r^2}\mathrm{cos}(Ht)+Hr\mathrm{cos}(\varphi )`$
$`B`$ $`=`$ $`\sqrt{1+H^2r^2}\mathrm{cos}(Ht)Hr\mathrm{cos}(\varphi )`$
$`U`$ $`=`$ $`\sqrt{1+H^2r^2}\mathrm{sin}(Ht)Hr\mathrm{sin}(\varphi )`$
$`V`$ $`=`$ $`\sqrt{1+H^2r^2}\mathrm{sin}(Ht)+Hr\mathrm{sin}(\varphi )`$ (3.2)
In these coordinates, the line-element (2.18) on the group manifold becomes
$$dS^2=(1+H^2r^2)dt^2+\frac{dr^2}{1+H^2r^2}+r^2d\varphi ^2$$
(3.3)
i.e., the standard parametrization of $`2+1`$ dimensional Anti de Sitter spacetime using static coordinates . As usual we unwind the temporal coordinate $`t`$, corresponding to considering the covering group of $`SL(2,R)`$. Moreover, the anti-symmetric tensor which can be read off from eq.(2.1), is given by
$$B_{t\varphi }=B_{\varphi t}=\frac{1}{2}Hr^2$$
(3.4)
with all other components vanishing.
In the $`(t,r,\varphi )`$ coordinates, the oscillating circular strings are given by
$`\varphi `$ $`=`$ $`\sigma `$
$`Ht`$ $`=`$ $`\mathrm{arctan}\left({\displaystyle \frac{1+EH}{\sqrt{1+2EH}}}\mathrm{tan}(\sqrt{1+2EH}\tau )\right)\tau `$
$`r`$ $`=`$ $`{\displaystyle \frac{E}{\sqrt{1+2EH}}}\mathrm{sin}(\sqrt{1+2EH}\tau )`$ (3.5)
where $`E`$ is an integration constant.
It is now straightforward to work backwards and read off the explicit expressions for the leftmoving and rightmoving group-elements, eqs.(2.9). The expressions are however not very enlightening, so we give them in Appendix A. It is more interesting to consider directly the leftmoving and rightmoving currents, eqs.(2.16). After some simple algebra, they are found to be
$`J_{}^\pm `$ $`=`$ $`\pm iEHke^{\pm i\sigma ^{}}`$
$`J_{}^3`$ $`=`$ $`EHk`$
$`J_+^\pm `$ $`=`$ $`iEHke^{i\sigma ^+}`$
$`J_+^3`$ $`=`$ $`EHk`$ (3.6)
Thus, the circular strings contain only modes corresponding to $`n=0`$ and $`n=\pm 1`$. This was of course to be expected for a circular string, c.f. circular strings in Minkowski spacetime, but it is in fact highly implicit in the parametrization (3.5). From the conformal field theory point of view, the parametrization (3.6) is thus more natural.
## 4 Quantization
In this section we give a short review of quantization of conformal field theories, corresponding to bosonic strings on group manifolds (see for instance Refs. ). This is mainly to fix our conventions and normalizations. We only give the results which we shall use in Section 5.
The currents $`J_\pm ^a`$, as introduced in eq.(2.12), can be expanded in Fourier series
$$J_{}^a=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}J_n^ae^{in\sigma ^{}};\left(J_n^a\right)^{}=J_n^a$$
(4.1)
and similarly for $`J_+^a`$ in terms of $`\sigma ^+`$. In the following we shall consider only the rightmovers $`()`$; the construction for the leftmovers $`(+)`$ is of course similar. For simplicity we shall therefore also skip the minus indices of $`J_{}^a`$, $`T_{}`$ etc. The $`SL(2,R)`$ Kac-Moody algebra is
$$[J_m^a,J_n^b]=iϵ_c^{ab}J_{m+n}^c+\frac{k}{2}m\eta ^{ab}\delta _{n+m}$$
(4.2)
In terms of the currents, eq.(2.15), the algebra becomes
$`[J_m^+,J_n^{}]`$ $`=2J_{m+n}^3+km\delta _{m+n}`$ (4.3)
$`[J_m^3,J_n^\pm ]`$ $`=\pm J_{m+n}^\pm `$
$`[J_m^3,J_n^3]`$ $`={\displaystyle \frac{k}{2}}m\delta _{m+n}`$
At the quantum level, the world-sheet energy-momentum tensor takes the Sugawara form
$$T=\frac{1}{k2}\eta _{ab}:J^aJ^b:=\frac{1}{k2}:(J^+J^{}J^3J^3):$$
(4.4)
Its Fourier modes
$$T=\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}L_ne^{in\sigma ^{}}$$
(4.5)
are given by
$$L_n=\frac{1}{k2}\underset{l=\mathrm{}}{\overset{\mathrm{}}{}}:\left(\frac{1}{2}(J_{nl}^+J_l^{}+J_{nl}^{}J_l^+)J_{nl}^3J_l^3\right):$$
(4.6)
They fulfill the Virasoro algebra
$$[L_m,L_n]=(mn)L_{m+n}+\frac{c}{12}m(m^21)\delta _{m+n}$$
(4.7)
where the central charge is given by
$$c=\frac{3k}{k2}$$
(4.8)
Demanding $`c=26`$, corresponding to conformal invariance, gives $`k=52/23`$. Notice also the commutators
$$[L_n,J_m^\pm ]=mJ_{n+m}^\pm ,[L_n,J_m^3]=mJ_{n+m}^3$$
(4.9)
which will be usefull in the following.
The Kac-Moody algebra contains the subalgebra of zero modes $`J_0^a`$, for which the quadratic Casimir is
$$Q=\eta _{ab}J_0^aJ_0^b=\frac{1}{2}\left(J_0^+J_0^{}+J_0^{}J_0^+\right)J_0^3J_0^3$$
(4.10)
The primary states, which are quantum states $`|jm>`$ at grade zero (”base-states” or ”ground-states”), are characterised by
$$Q|jm>=j(j+1)|jm>,J_0^3|jm>=m|jm>$$
(4.11)
Moreover, they fulfill
$$J_0^\pm |jm>=\sqrt{m(m\pm 1)j(j+1)}|jm\pm 1>$$
(4.12)
as well as
$$J_l^a|jm>=0;l>0$$
(4.13)
The primary states must belong to one of the unitary representations of $`SL(2,R)`$ (or its covering group) . We shall return to this point in the next section.
From the primary states, one can construct the excited states as descendents by applying $`J_l^a`$ operators ($`l`$ is a positive integer)
$$|\psi >=J_{l_1}^{a_1}J_{l_2}^{a_2}\mathrm{}..J_{l_r}^{a_r}|jm>$$
(4.14)
and so on. The physical state conditions (the mass-shell condition and the Virasoro primary condition) are then
$$(L_01)|\psi >=0$$
(4.15)
$$L_l|\psi >=0;l>0$$
(4.16)
For a physical state of the form (4.14) at grade $`n=l_i`$, the mass-shell condition gives
$$n\frac{j(j+1)}{k2}=1$$
(4.17)
Identifying (minus) the quadratic Casimir, eq.(4.10), of the base-state with the mass-squared , or more precisely we normalize the mass as
$$m^2\alpha ^{}\frac{j(j+1)}{k2}$$
(4.18)
we then see that the mass-squared grows linearly with the grade
$$m^2\alpha ^{}=n1$$
(4.19)
This is just like for strings in flat Minkowski spacetime.
## 5 Coherent String States
The idea is now to construct exact quantum states with properties similar to the classical circular strings considered in Section 3. More precisely, our aim is to construct states $`|\psi >`$ such that the only components of the currents giving non-vanishing expectation values, $`<\psi |J^a|\psi >0`$, are those components corresponding to the non-vanishing classical currents, eq.(3.6). In other words, considering for simplicity only the rightmovers, then only the components $`J_1^+`$, $`J_{+1}^{}`$ and $`J_0^3`$ should have non-vanishing expectation values; all other components $`J_n^+(n1)`$, $`J_n^{}(n+1)`$ and $`J_n^3(n0)`$ must have zero expectation values.
Several problems immediately appear: If we consider states of the form (4.14), it would be possible to obtain a non-vanishing expectation value for $`J_0^3`$, but it would certainly be impossible to get non-vanishing expectation values of $`J_1^+`$ and $`J_{+1}^{}`$. Moreover, as mentioned at the end of Section 4, states of the form (4.14) give rise to a mass-spectrum where the mass-squared grows linearly with the grade as in flat space. However, semi-classical quantization of the circular strings has been shown to lead to a mass-spectrum where $`m^2\alpha ^{}N^2`$ ($`N`$ positive integer), at least for the high mass states. Fortunately, it turns out that both problems can be solved by considering coherent string states on the $`SL(2,R)`$ group manifold (a similar construction for the $`SU(2)`$ group manifold was considered in ).
As a starting point, we consider states of the form
$$\left(J_1^+\right)^n|jj>;n0$$
(5.1)
where $`|jj>`$ belongs to the highest weight discrete series $`D_j^{}`$ , with states $`|jm>`$
$$j1/2,m=j,j1,\mathrm{}$$
(5.2)
Since we shall consider the covering group of $`SL(2,R)`$, there are no further restrictions on $`j`$, i.e., it need not be integer or half-integer . In particular, from eq.(4.12) it follows that
$$J_0^+|jj>=0,J_0^{}|jj>=\sqrt{2j}|jj1>$$
(5.3)
The states eq.(5.1) fulfill the Virasoro primary condition
$$L_l\left(J_1^+\right)|jj>=0;l>0$$
(5.4)
and the mass-shell condition (4.15) leads to
$$n=1+\frac{j(j+1)}{k2}j=\frac{1}{2}\sqrt{(k2)(n1)+1/4}$$
(5.5)
However, these states generally do not have positive norm. Indeed
$$<jj|\left(J_{+1}^{}\right)^m\left(J_1^+\right)^n|jj>=n!\delta _{nm}\underset{i=1}{\overset{n}{}}\left(k2+i\sqrt{4(k2)(n1)+1}\right)$$
(5.6)
and the right hand side is generally not positive. For example, it is negative for $`n=m=2`$ and $`n=m=3`$ (using that $`k=52/23`$). This is of course just a simple example illustrating the well known unitarity problem for strings on $`SL(2,R)`$ \[1-10, 13\] (for recent reviews, see ).
We consider instead coherent states built from states of the form (5.1)
$$e^{\mu J_1^+}|jj>$$
(5.7)
where $`\mu `$ is an arbitrary complex number. These states certainly also fulfill the Virasoro primary condition but, being coherent states, they obviously are eigenstates of neither the number operator nor of the $`L_0`$ operator. We shall therefore impose a ”weak” mass-shell condition
$$<jj|e^{\mu ^{}J_{+1}^{}}\left(L_01\right)e^{\mu J_1^+}|jj>=0$$
(5.8)
Before evaluating the left hand side of eq.(5.8), it is necessary to consider the normalization of the states (5.7).
In ordinary quantum mechanics with creation and annihilation operators $`a^{}`$ and $`a`$, respectively, and a vacuum state $`|0>`$
$$[a,a^{}]=1,a|0>=0$$
(5.9)
the excited states are constructed as
$$\left(a^{}\right)^n|0>=\sqrt{n!}|n>,<n|m>=\delta _{nm}$$
(5.10)
In that case, a coherent state can always be normalized. In fact, the coherent state
$$|\mu >e^{\frac{1}{2}\mu ^{}\mu }e^{\mu a^{}}|0>$$
(5.11)
has unit norm, for arbitrary complex number $`\mu `$. Notice also that the coherent state is an eigenstate of the annihilation operator
$$a|\mu >=\mu |\mu >$$
(5.12)
which can be taken as the definition of coherent states in ordinary quantum mechanics. (For more discussion, see for instance ).
In our case, the situation is somewhat different since we have a Kac-Moody algebra (4.2) with a non-Abelian term in the current algebra, and in particular since the group manifold $`SL(2,R)`$ is non-compact and has a time-like direction (contrary to the case of $`SU(2)`$ ). It implies that the coherent state (5.7) is not an eigenstate of the annihilation operator $`J_{+1}^{}`$
$$J_{+1}^{}e^{\mu J_1^+}|jj>=\mu \left(2j+k+\mu J_1^+\right)e^{\mu J_1^+}|jj>$$
(5.13)
Moreover, the coherent state (5.7) can not be normalized for arbitrary complex number $`\mu `$. One finds
$$<jj|e^{\mu ^{}J_{+1}^{}}e^{\mu J_1^+}|jj>=1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{(\mu ^{}\mu )^n}{n!}\underset{l=1}{\overset{n}{}}(2j+k1+l)$$
(5.14)
The product on the right hand side goes as $`n!`$. Thus the infinite sum is convergent only if $`\mu ^{}\mu <1`$, or if the infinite sum terminates after a finite number of terms (this happens if $`2j+k1+l=0`$, for some $`l`$). More precisely, the right hand side of eq.(5.14) is a finite positive number in the following two cases
(I): $`\mu ^{}\mu <1`$ and $`j`$ arbitrary ($`j1/2`$).
In this case the normalized state is
$$|\mu I>=(1\mu ^{}\mu )^{j+k/2}e^{\mu J_1^+}|jj>$$
(5.15)
(II): $`\mu ^{}\mu >1`$ and $`j=Nk/2(N=0,1,2,\mathrm{}).`$
In this case the normalized state is
$$|\mu II>=(\mu ^{}\mu 1)^Ne^{\mu J_1^+}|Nk/2,Nk/2>$$
(5.16)
Let us now return to the mass-shell condition, eq.(5.8), which gives rise to some additional constraints on $`\mu `$ and $`j`$. In the two cases, respectively, one finds
(I):
$$\mu ^{}\mu =\frac{1+\frac{j(j+1)}{k2}}{2j+k+1+\frac{j(j+1)}{k2}}<1;\frac{k}{2}<j\frac{1}{2}$$
(5.17)
(II):
$$\mu ^{}\mu =\frac{1+\frac{j(j+1)}{k2}}{2j+k+1+\frac{j(j+1)}{k2}}>1;j=N\frac{k}{2}(N=1,2,\mathrm{})$$
(5.18)
Thus, the spectrum consists of two parts: (I) A continuous spectrum where $`j`$ fulfills the standard spin-level condition \[2-4, 6-10, 13, 14\] $`k/2<j1/2`$, and (II) a discrete spectrum where $`j`$ fulfills $`j=Nk/2`$ ($`N`$ positive integer). If we were considering ordinary descendent states of the form (4.14), we would generally not be allowed to have $`|j|>k/2`$ because of unitarity, thus the discrete part (II) would not be allowed. For the coherent states under consideration here, there is however no problem since the quantization condition $`j=Nk/2`$ ($`N`$ positive integer) precisely ensures that they are all positive norm states.
Introducing the mass with the normalization as in eq.(4.18), we find that the continuous part of the spectrum (I) corresponds to
$$m^2\alpha ^{}[\frac{1}{4(k2)},\frac{k}{4}[=[\frac{23}{24},\frac{13}{23}[$$
(5.19)
where the last equality was obtained using $`k=52/23`$, corresponding to conformal invariance. That is, the continuous part of the spectrum (I) consists of low mass states and is partly tachyonic.
On the other hand, the discrete spectrum (II) gives
$$m^2\alpha ^{}=\frac{N^2}{k2}\left(1+\frac{k1}{N}+\frac{k(k2)}{4N^2}\right)$$
(5.20)
i.e., for the discrete part of the spectrum we find that $`m^2\alpha ^{}N^2`$ ($`N`$ positive integer). Asymptotically, this is precisely what was found using semi-classical quantization . It should be stressed, however, that $`N`$ is not the eigenvalue of the number operator; as already mentioned, since we are working with coherent states, we do not have any eigenstates of the number operator. Thus $`N`$ is simply a positive integer here.
Notice also that $`k`$, $`\alpha ^{}`$ and the length-scale $`H`$ in the quantum theory are related as in eq.(2.19), but with $`k`$ replaced by $`k2`$. With the present normalizations we therefore have exactly the same leading order behaviour, including the numerical coefficient, for the mass-squared as obtained using semi-classical quantization in Ref. .
Finally, the relation with the classical circular strings is established by noting that only the expectation values of $`J_1^+`$, $`J_{+1}^{}`$ and $`J_0^3`$ are non-vanishing ($`i=I,II`$)
$`<\mu i|J_l^+|\mu i>`$ $`=`$ $`\{\begin{array}{cc}(2j+k)\frac{\mu ^{}}{1\mu ^{}\mu },l=1& \\ 0,l1& \end{array}`$ (5.23)
$`<\mu i|J_l^{}|\mu i>`$ $`=`$ $`\{\begin{array}{cc}(2j+k)\frac{\mu }{1\mu ^{}\mu },l=+1& \\ 0,l+1& \end{array}`$ (5.26)
$`<\mu i|J_l^3|\mu i>`$ $`=`$ $`\{\begin{array}{cc}j+(2j+k)\frac{\mu ^{}\mu }{1\mu ^{}\mu },l=0& \\ 0,l0& \end{array}`$ (5.29)
valid for both $`|\mu I>`$ and $`|\mu II>`$ for the respective values of $`\mu `$ and $`j`$, as given in eqs.(5.17)-(5.18). To obtain these results, as well as most other results in this section, we used the commutators listed in Appendix B.
## 6 Conclusion
We have shown that very massive string states in the $`SL(2,R)`$ WZWN model (corresponding to $`AdS_3`$) can be described as coherent states based on the standard discrete representation $`D_j^{}`$. The spectrum of such states was shown to consist of two parts: A continuous low mass (partly tachyonic) part and a discrete high mass part. For the continuous part, the spin $`j`$ must fulfill the standard spin-level restriction $`k/2<j1/2`$, while for the discrete part we get the quantization condition $`j=Nk/2`$ ($`N`$ positive integer). Although the latter seems to be in contradiction with the spin-level restriction, the quantization condition precisely ensures that all our coherent states have finite positive norm. Thus, no ghost-states are included in the spectrum. The mass spectrum of the discrete part of the spectrum, eq.(5.20), shows the asymptotic behaviour $`m^2\alpha ^{}N^2`$. This is in precise agreement with our previous results obtained using semi-classical quantization . The same asymptotic behaviour was also obtained in the recent preprint , although the construction there is completely different from ours, as discussed in more detail in the introduction.
In this paper we used, for simplicity and clarity, the example of an oscillating circular string. That is, our quantum coherent string states were constructed to lead to non-vanishing expectation values for very specific components of the currents, eq.(5.21). It is however easy to generalize our construction to other string configurations too.
We saw in Section 5 that the coherent states, eq.(5.7), do not have all the same properties of standard quantum mechanical coherent states . For instance, they are not eigenstates of the annihilation operator. All the standard properties are however obtained in the following manner: First we must renormalize the currents $`J_n^a\sqrt{k}J_n^a`$, as follows from the algebra, eq.(4.3) (although this does not work for the zero-modes ). Secondly, for the coherent states we must renormalize the complex number $`\mu `$ by $`\mu \mu /\sqrt{k}`$, as follows from eq.(5.7). Finally, we let $`k\mathrm{}`$. Then, it follows that the non-Abelian piece drops out in eq.(5.13) and we get an eigenstate of the renormalized annihilation operator. More generally, by the same prescription we recover all the usual well-known properties of standard quantum mechanical coherent states . For instance, the continuous spectrum, represented by the states (5.15), will now be valid for any $`j1/2`$, and the states become
$$|\mu I>e^{\frac{1}{2}\mu ^{}\mu }e^{\mu J_1^+}|jj>;k\mathrm{}$$
(6.1)
c.f. eq.(5.11). On the other hand, the discrete part of the spectrum, represented by the states (5.16), disappears since all the states are pushed to infinite mass. These are precisely the properties which characterise the usual coherent states of quantum mechanics as quasi-classical states, for which the quantum uncertainty is minimal.
Acknowledgements
A.L.L. would like to thank the Ambassade de France à Copenhague, Service Culturel et Scientifique for financial support in Paris, during the preparation of this paper.
## 7 Appendix A
In this appendix we give the explicit expressions for the group-elements (2.9) corresponding to the circular strings (3.5). For simplicity, we only give the results for the rightmovers.
$`a(\sigma ^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+2EH}}}\mathrm{sin}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)\left[(1+EH)\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)+EH\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\right]`$ (7.1)
$`+`$ $`\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\mathrm{cos}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)`$
$`b(\sigma ^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+2EH}}}\mathrm{sin}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)\left[(1+EH)\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)EH\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\right]`$ (7.2)
$`+`$ $`\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\mathrm{cos}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)`$
$`u(\sigma ^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+2EH}}}\mathrm{sin}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)\left[(1+EH)\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)+EH\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\right]`$ (7.3)
$``$ $`\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\mathrm{cos}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)`$
$`v(\sigma ^{})`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{1+2EH}}}\mathrm{sin}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)\left[(1+EH)\mathrm{cos}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)EH\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\right]`$ (7.4)
$``$ $`\mathrm{sin}\left({\displaystyle \frac{\sigma ^{}}{2}}\right)\mathrm{cos}\left(\sqrt{1+2EH}{\displaystyle \frac{\sigma ^{}}{2}}\right)`$
Now using eq.(2.16), it is straightforward to obtain (3.6) for the rightmovers. The derivation for the leftmovers is similar.
## 8 Appendix B
In this appendix we list some useful commutators used in Section 5.
$$[J_m^3,\left(J_1^+\right)^n]=nJ_{m1}^+\left(J_1^+\right)^{n1}$$
(8.1)
$$[J_m^+,\left(J_{+1}^{}\right)^n]=n\left(2J_{m+1}^3km\delta _{m+1}\right)\left(J_{+1}^{}\right)^{n1}n(n1)J_{m+2}^{}\left(J_{+1}^{}\right)^{n2}$$
(8.2)
$$[J_m^{},\left(J_1^+\right)^n]=n\left(J_1^+\right)^{n1}\left(2J_{m1}^3+km\delta _{m1}\right)+n(n1)\left(J_1^+\right)^{n2}J_{m2}^+$$
(8.3)
These identities are most easily proved by induction using eqs.(4.3). It follows that
$$[J_m^3,e^{\mu J_1^+}]=\mu J_{m1}^+e^{\mu J_1^+}$$
(8.4)
$$[J_m^+,e^{\mu ^{}J_{+1}^{}}]=\mu ^{}\left(2J_{m+1}^3km\delta _{m+1}+\mu ^{}J_{m+2}^{}\right)e^{\mu ^{}J_{+1}^{}}$$
(8.5)
$$[J_m^{},e^{\mu J_1^+}]=\mu e^{\mu J_1^+}\left(2J_{m1}^3+km\delta _{m1}+\mu J_{m2}^+\right)$$
(8.6) |
warning/0001/astro-ph0001362.html | ar5iv | text | # Mirror Dark Matter and Galaxy Core densities
## I Introduction
Dark matter constitutes the bulk of the matter in the universe and a proper understanding of the nature of the new particle that plays this role has profound implications not only for cosmology but also for particle physics beyond the standard model. It is therefore not surprising that one of the major areas of research in both particle physics and cosmology continues to be the physics of dark matter. Apart from the simple requirement that the right particle physics candidate must have properties that yield the requisite relic density and mass to dominate the mass content of the universe, it should be required to provide a satisfactory resolution of three puzzles of dark matter physics: (i) why is it that the contribution of baryons to the mass density ($`\mathrm{\Omega }`$) of the universe is almost of the same order as the contribution of the dark matter to it ? (ii) how does one understand the dark objects with mass $`0.5M_{}`$ observed in the MACHO experiment, which are supposed to constitute up to 50% of the mass of the halo of the Milky way galaxy and presumably be connected to the dark constituent that contributes to $`\mathrm{\Omega }`$ (in a manner that satisfies the “environmental impact” conditions of Freese et al? and, finally (iii) what explains the density profile of dark matter in galactic halos– in particular, the apparent evidence in favor of the fact that the core density of galactic halos remain constant as the radius goes to zero.
There are many particle physics candidates for the dark constituent of the universe. Generally speaking, the prime consideration that leads to such candidates is that they yield the right order of magnitude for the relic density and mass necessary to get the desired $`\mathrm{\Omega }_{DM}0.21`$. This is, of course, the minimal criterion for any such candidate and requires that the annihilation cross section of the particles must be in a very specific range correlated with their mass. The most widely discussed candidates are the lightest supersymmetric particle (LSP) and the Peccei-Quinn particle, the axion. The first is expected to have a mass in the range of 100 GeV whereas the mass of the second would be in the range of $`10^6`$ eV. Compare these values with the proton mass of one GeV). To understand within these models why $`\mathrm{\Omega }_B\mathrm{\Omega }_{DM}`$, one needs to work in a special range for the parameters of the theory. In either of these pictures, the MACHO observations must have a separate explanation. Thus it may not be unfair to say that these two most popular candidates do not resolve the first two of the three dark matter puzzles. In recent years it has been emphasized that the LSP and the axion may also have difficulty in explaining the observed core density behaviour of dwarf spheroidal galaxies which are known to be dark matter dominated. The point here is that both the axions and neutralinos, being collisionless and nonrelativistic, accumulate at the core of galactic halos, leading to a core density $`\rho (R)`$ which goes like $`R^2`$ rather than a constant which seems to fit data better. We will refer to this as the core density puzzle. This last puzzle has motivated Spergel and Steinhardt to revive and reinvigorate an old idea that dark matter may be strongly self interacting, which for the right range of the parameters of the particle may lead to a halo core which is much less dense and hence in better agreement with observations. To be more specific, it was noted in that if the dark matter particle is self interacting and has mean free path of collision of about a kpc to a Mpc, then the core on this scale cannot “keep on accumulating” dark matter particle, since these will now scatter and “diffuse out”. For typical dark matter particle densities of order of one particle per cm<sup>3</sup>, this requires a cross section for scattering of $`\sigma 10^{21}10^{24}`$ cm<sup>2</sup>. Furthermore, in order to prevent dissipation which would lead to cooling and collapse to the core, one has a lower limit on the mass of the exchanged particle that must exceed typical “virial” energy of particles ($``$ keV). An alternative possibility is that the core is optically thick to exchanged particles. If these considerations stand the test of time, a theoretical challenge would be to look for alternative dark matter candidates (different from the popular ones described above) and the associated scenarios for physics beyond the standard model. A class of models known as mirror universe models have recently been discussed. These are motivated theoretically by string theory and experimentally by neutrino physics. These models predict the existence of a mirror sector of the universe with matter and force content identical to the familiar sector (prior to symmetry breaking). Symmetry breaking may either keep the mirror symmetry exact or it may break it. This leads to two classes of mirror models: the symmetric mirror model, where all masses and forces in the two sectors remain the same after symmetry breaking and the asymmetric mirror model where the masses in the mirror sector are larger than those in the familiar sector. The mirror particles interact with the mirror photon and not the familiar photon so that they remain dark to our observations. Since the the lightest particles of the mirror sector (other than the neutrinos), the mirror proton and the mirror electron (like in the familiar sector) are stable and will have abundances similar to the familiar protons and electrons, the proton being heavier could certainly qualify as a dark matter candidate. It has indeed been pointed out that, consistent with the cosmological constraints of the mirror universe theory, the mirror baryons have the desired relic density to play the role of dark matter of the universe. The additional neutrinos of the mirror sector are the sterile neutrinos which appear to be needed in order to have a simultaneous understanding of the three different neutrino oscillations i.e. solar, atmospheric and the LSND observations. In fact, one view of neutrino oscillation explanations of these phenomena fixes the ratio of familiar particle mass to the mirror particle mass thereby narrowing down the freedom of the mirror sector parameters. If indeed sterile neutrinos turn out to be required, mirror universe theory is one of the few models where they appear naturally with mass in the desired range. If we denote the mass ratio $`m_p^{}/m_p=\zeta `$, then a value of $`\zeta 10`$ is required to explain the neutrino puzzles. What is more interesting is that for the same range of parameters that are required to solve the neutrino puzzles, (i.e. $`\zeta 10`$) mirror matter can also provide an explanation of the microlensing observations\- in particular why the observed MACHOs have a mass very near the solar mass and are still dark.
It is the goal of this paper to show that the same mirror universe framework can also explain the core density puzzle of galactic halos. The basic idea is that mirror sector $`H^{}H^{}`$ (i.e. mirror hydrogen) scattering, with its large geometric cross section, is a natural candidate for strongly interacting dark matter of Ref.. Thus the mirror matter model has the desirable properties that it can naturally explain all three dark matter puzzles. It is worth noting that the asymmetric mirror model was not originally designed for this purpose but rather to explain the neutrino puzzles and indeed it is gratifying that slight modification of the framework (increasing the QCD scale) that solves the neutrino puzzles also solves the dark matter puzzles.
## II Mirror matter as dark matter
Let us start with a brief overview of the mirror matter models . The basic idea of the model is extremely simple: duplicate the standard model or any extension of it in the gauge summetric Lagrangian (and allow for the possibility that symmetry breaking may be different in the two sectors). There is an exact mirror symmetry connecting the Lagrangians (prior to symmetry breaking) describing physics in each sector. Clearly the $`W^{}s,\gamma ^{}s`$ etc of each sector are different from those in the other as are the quarks and leptons. When the symmetry breaking scale is different in the two sectors, we will call this the asymmetric mirror model. The QCD scale being an independent scale in the theory could be arbitrary. We will allow both the weak scale as well as the QCD scale of the mirror sector to be different from that of the familiar sector and assume the same common ratio $`\zeta `$ for both scales i.e. $`<H^{}>/<H>=\mathrm{\Lambda }^{}/\mathrm{\Lambda }\zeta `$. It is assumed that the two sectors in the universe are connected by only gravitational interactions. It was shown in that gravity induced nonrenormalizable operators generate mixings between the familiar and the mirror neutrinos, which is one of the ingredients in the resolution of neutrino puzzles. It is of course clear that both sectors of the universe are co-located. Together they evolve according to the rules of the usual big bang model except that the cosmic soups in the two sectors may have different temperature. In fact the constraints of big bang nucleosynthesis require that the post inflation reheat temperature in the mirror sector $`T_R^{}`$ be slightly lower than that in the familiar sector $`T_R`$ (define $`\beta T_R^{}/T_R`$) so that the contribution of the light mirror particles such as $`\nu ^{},\gamma ^{}`$ etc. to nucleosynthesis is not too important. This has been called asymmetric inflation and can be implemented in different ways. Present discussions of BBN can be used to conclude that roughly $`\beta ^41/6`$ is equivalent to $`\delta N_\nu 1`$. Before proceeding to any detailed discussion, let us first note the impact of the asymmetry on physical parameters and processes. First it implies that $`m_i\zeta m_i`$ with $`i=n,p,e,W,Z`$. This has important implications which have been summarized before. For instance, the binding energy of mirror hydrogen is $`\zeta `$ times larger so that recombination in the mirror sector takes place much earlier than in the visible sector. With $`\beta T_R^{}/T_R`$ as above, mirror reombination temperature is $`\zeta /\beta T_r`$ where $`T_r`$ is the recombination temperature in the familiar sector. The mirror sector recombination takes place before familiar sector recombination; this means that density inhomogeneities in the mirror sector begin to grow earlier and familiar matter can fall into it later. One can also compute the contribution of mirror baryons to the mass density of the universe as follows:
$`{\displaystyle \frac{\mathrm{\Omega }_B^{}}{\mathrm{\Omega }_B}}\beta ^3\zeta `$ (1)
Here we have assumed that baryon to photon ratio in the familiar and the mirror sectors are the same as would be expected since the dynamics are same in both sectors due to mirror symmetry. Eq. (1) implies that both the baryonic and the mirror baryon contribution to $`\mathrm{\Omega }`$ are roughly of the same order, as observed. This provides a resolution of the first conceptual puzzle. Furthermore if we take $`\mathrm{\Omega }_B0.05`$, then $`\mathrm{\Omega }_B^{}0.2`$ would require that $`\beta =(4/\zeta )^{1/3}`$. From this one can calculate the effective $`\delta N_\nu `$ using the following formula:
$`\delta N_\nu =3\beta ^4+{\displaystyle \frac{4}{7}}\beta ^4({\displaystyle \frac{11}{4}})^{4/3}`$ (2)
where the last factor $`(11/4)^{4/3}`$ is due to the reheating of the mirror photon gas subsequent to mirror $`e^+^{}e^{^{}}`$ annihilation. For $`\zeta =20`$, this implies $`\delta N_\nu 0.6`$ and it scales with $`\zeta `$ as $`\zeta ^{4/3}`$. Thus in principle the idea that mirror baryons are dark matter could be tested by more accurate measurements of primordial He<sup>4</sup>, Deuterium and Li<sup>7</sup> abundances which determine $`\delta N_\nu `$. Clearly to satisfy the inflationary constraint of $`\mathrm{\Omega }_{TOT}=1`$, we need $`\mathrm{\Omega }_\mathrm{\Lambda }0.7`$. These kinds of numbers for cold dark matter density emerge from current type I supernovae observations. It is interesting to note that if one were to require that $`\mathrm{\Omega }_{CDM}=1`$, the mirror model would require that $`\zeta `$ be much larger (more than 100) which would then create difficulties in understanding both the neutrino data and the microlensing anomalies. Thus mirror baryons satisfy the most requirement to be the cold dark matter of the universe.
## III Mirror matter collision and core density of galaxies
To discuss further implications of the mirror cold dark matter for structure formation and the nature of the dark halos, we need to know various cross sections. Using the asymmetry factor $`\zeta `$, it is easy to see that weak cross sections varies as $`\zeta ^4`$ i.e. $`\sigma _W^{}=\sigma _W/\zeta ^4`$, the Thomson cross section $`\sigma _T^{}=\sigma _T/\zeta ^2`$. Nuclear cross sections will also be different in the mirror sector due to different values of the QCD scale in the two sectors. For $`\mathrm{\Lambda }^{}(1015)\mathrm{\Lambda }`$, we would expect $`\sigma _{N^{}N^{}}^{}=(\sigma _{NN}/\zeta ^2)\times \left(\frac{g_{\pi ^{}N^{}N^{}}}{g_{\pi NN}}\right)^4`$ (for fixed values of energy). With these simple rules, assuming the pion nucleon couplings in both sectors to be identical, we find that $`\sigma _{N^{}N^{}}`$ to be of order $`10^{30}`$ cm<sup>2</sup> or so. This is clearly too small to make a difference in the core density problem. Let us now focus on the atomic forces. We would expect the mirror dark matter to be mostly in the form of atomic hydrogen ($`H^{}`$). In the familiar sector, the hydrogen atomic scattering is of order $`\pi a_0^2`$, (where $`a_0`$ is the Bohr radius) and is of order $`10^{16}`$ cm<sup>2</sup>. Since the Bohr radius scales like $`\zeta ^1`$ when we consider the mirror sector, we would expect the collisional mean free path due to atomic scattering to be of order $`10^{17}\zeta ^3`$ cm. For $`\zeta =20`$, the mean free path is about $`0.3`$ kpc. This is slightly smaller than the lowest allowed value given in Ref. . However, as we argue below, the dangers for lower mean free paths envisioned in Ref. , do not apply to our case due to the dynamics of the mirror universe. As noted in Ref., an additional constraint on the self interacting dark matter arises from the fact that the forces of self interaction must not be such as to allow significant loss of energy from the core of halos since otherwise, core will lose energy and collapse giving us back the problem we wanted to cure in the first place. Essentially similar reasoning led to the lower bound of 1 kpc on the mean free path in the Ref.. In the mirror dark matter model, even though we have mirror photons, there is no dissipation problem. The point is that when there is a collision, it will excite the atom. The atom radiates when it falls back down, but that radiation can be absorbed. Sometimes you ionize, then you get a plasma which also absorbs the radiation. The result would seem to be something like the sun in which it takes a long time for the radiation to get out. In other words, the mirror matter core is optically thick. In fact, one can estimate the number of collisions a core particle makes since we know the mfp is in the range of 1 - 0.1 kpc. Taking mfp/virial velocity, we get a time of 10-100Myr so it makes $`10^310^4`$ collisions during the age of the universe. We can also give a crude estimate of the fractional energy loss per particle by using $`\sigma T^4`$ for the energy radiated per unit area per unit time, dividing by the number of particles and multiplying by $`10^{17}`$ s - the age of the universe, and multiplying by thearea of the “core” of radius 1 kpc. This fractional energy loss per particle is negligible enabling us to have a lower mean free path than the Ref.. For the case of mirror matter, the core is protected against this collapse by the long time it takes to radiate away energy which is in the form of (mirror) photons that can’t get out. The essential point is the similarity of the mirror sector with the familiar sector, where we know that the photons emitted from the core of a star are very few in number due to minimum ratio of surface to volume and therefore do not lead to collapse of the galaxy cores.
## IV Structure Formation
In this section we will try to make plausible that, in spite of the $`\zeta ^2`$ decrease in cross sections, the facts that (a) structure formation begins earlier in the mirror sector (because recombination occurs before matter-radiation equality) and consequently (b) mirror temperatures are higher, for the same processes, than familiar temperatures, permit formation of galactic and smaller structures. In doing this, we will make use of our previous work in and , as well as that of Tegmark et al . Much of the work of can be carried over to the present work, after suitable modification to take into account the fact that, in the current model, the proton mass scales as $`\zeta `$. Here, we will assume that primordial perturbations are ”curvature” or ”adiabatic” perturbations. This means that the scale of the largest structures are set by mirror Silk damping . $`\gamma ^{}`$ diffusion wipes out inhomogeneities until the $`\gamma ^{}`$ mean free path,
$$\lambda ^{}=[\sigma _T\zeta ^2n_e^{^{}}]^1$$
(3)
where $`\zeta ^2\sigma _T`$ is the mirror matter Thomson cross section and $`n_e^{}`$ is its electron number density, becomes greater than one third the horizon distance ($`ct`$). Silk damping turns off because the $`\lambda ^{}`$ increases as $`z^3`$ while $`ct`$ only increases as $`z^2`$. First, we compute, from Silk damping, the masses of the largest structures in this picture. Since structure formation starts with the mirror sector, our assumption is that familiar sector particles will later fall into these. For numerical values below, we will take, $`h=0.7`$ and $`\mathrm{\Omega }_{\stackrel{~}{B}}=0.2`$. We pick $`t(z_1/z)^2s`$ with $`z_1=4\times 10^9`$ and $`n_{\stackrel{~}{e}}=\mathrm{\Omega }_{tildeB}\rho _cz^3/(\zeta m_p)`$ with $`\rho _c=1.9h^2\times 10^{29}g/cm^3`$. Silk damping stops at around $`z_{sd}8\zeta ^3`$ which gives
$`\lambda _{sd}2.5\times 10^{27}/\zeta ^6cm`$ (4)
$`M_{sd}10^{54}/\zeta ^9gm`$ (5)
Note that, for $`\zeta 10`$, this is about the mass (and size) of a large galaxy. This coincidence could be an important factor in understanding galaxy sizes should this model correspond to reality. As in , we parametrize the separation of $`M_{sd}`$ from the expanding universe as taking place at
$$z_{stop}=z_Mz_{sd}$$
(6)
with
$$R_G=\lambda _{sd}/z_M$$
(7)
After violent relaxation we have for the proton temperature
$$T_p=GM_G\zeta m_p/R_G10^4z_M/\zeta ^2ergs$$
(8)
with $`\rho `$, outside the central plateau, given by
$$\rho (R)=A/R^210^{26}z_M/(\zeta ^3R^2)gm/cm^3$$
(9)
We now turn to the question of whether this isothermal sphere is likely to fragment and form mirror stars. For this we compute the amount of mirror molecular hydrogen since it is its collisional excitation (and subsequent radiation) that is believed to be the chief mechanism that provides cooling for formation of stars. If the rate for this mechanism is faster than the rate for free fall into the mass of the structure at issue, we can expect local regions to cool fast enough to result in fragmentation of that structure. We do here a very rough estimate of mirror galaxy fragmentation into mirror globular clusters, using the results of , but leave to a more detailed work further fragmentation into the $`0.5M_{}`$ structures predicted in Reference give a useful approximation to their numerical results for the fraction of molecular hydrogen, $`f_2`$ ($`f_0`$ denotes its primordial value):
$$f_2(t)=f_0+(k_m/k_1)ln[1+x_0nk_1t)$$
(10)
where, as a first try, $`k_m`$ can be taken as just the rate for $`H+e^{}H^{}+\gamma `$ (which they conveniently give as about $`2\times 10^{18}T^{.88}cm^3s^1`$), while $`k_1`$ is the rate for $`H^++e^{}H+\gamma `$ ($`2\times 10^{10}T^{0.64}cm^3s^1`$). Equation (8) is the result of $`H_2`$ production from the catalytic reactions $`H+e^{}H^{}+\gamma `$ followed by $`H+H^{}H_2+e^{}`$ competing against the recombination reaction that destroys the catalyst, free electrons, (approximately) as $`1/t`$ (assuming constant density). Our goal here is to show from Equation (8) that it is plausible that $`f_2`$, the fraction of molecular hydrogen, rises from its primordial value of $`10^6`$ to the region above $`10^4`$ where cooling tends to be competitive with free fall. First, we note that, if $`k=<v\sigma >AT^\gamma cm^3/s`$, for familiar e and p, we expect that, for mirror e and p, scaling with $`\zeta `$ to go as $`\zeta ^{(2+\gamma )}AT^\gamma `$, since $`\sigma `$ must go as $`\zeta ^2`$ and all factors of $`T`$ must be divided by some combination of $`m_e`$ and $`M_p`$, both of which go as $`\zeta `$ (making this model much easier to compute from than that of ). We now estimate fragmentation. From Equation (6) we see that the galactic temperature should begin at about 10 eV at a time when the cosmic temperature is about 1 eV and the cosmic gamma number density is about $`10^9/cm^3`$. The rate for “compton cooling” is very fast at this high density (unlike at later times for the familiar case) and there should be rapid cooling to about 1 eV. We can now compute the Jeans length for fragments as a function of distance R from galaxy central. We use
$$\rho _J=(T/Gm)^3/M^2$$
(11)
If we set the Jeans mass, $`M`$, to $`4\pi r^3\rho _J/3`$, we can solve for $`r`$ obtaining (if we are careful to convert $`T`$ in Equation (6) from ergs)
$$r=R[10^7\zeta ^2/z_M]^{1/2}10^{3.5}R$$
(12)
Now inserting into Equation (8) gives the coefficient of the log term on the order of $`10^2`$ and the argument varying from $`10^{13}`$ to $`10^{17}`$ as $`R`$ varies from 1 to 100 kiloparsecs while the free fall time ($`(G\rho )^{1/2}`$) varies from $`10^{14}`$ to $`10^{16}`$. This would appear to indicate the likelihood of fragmentation of the original Silk damping structure into smaller units, (outside the optically thick core) and the eventual formation of the $`0.5M_{odot}`$ black holes that explain the microlensing events of .
## V Conclusion
The asymmetric mirror model was originally proposed to solve the neutrino puzzles. Subsequently, it was advocated as providing an alternative dark matter candidate. Then it was shown to have the advantage of resolving the microlensing anomaly in a “non-polluting” manner. In this paper we have argued that the model could additionally provide an explanation of a fourth problem. It appears to be a realization of the mechanism of Spergel and steinhardt for understanding the core density profile of galaxies by means of builtin self interactions of mirror matter. The work of R. N. M. is supported by the National Science Foundation grant under no. PHY-9802551 and the work of V. L. T. is supported by the DOE under grant no. DE-FG03-95ER40908. We like to thank P. Steinhardt for some discussions. |
warning/0001/gr-qc0001071.html | ar5iv | text | # Untitled Document
How to test vector nature of gravity
Igor E. Bulyzhenkov
Institute of Spectroscopy, RAS, Troitsk, 142092 Moscow reg., Russia
The covariant scheme is proposed to couple gravity and electrodynamics in pseudo-Riemannian four-spaces with electromagnetic connections. Novel dynamics of the charged particle and electromagnetic dilation-compression of its proper time can be tested in non-relativistic experiments. The vector equations acknowledge unified photon waves without metric modulations of flat laboratory space.
PACS numbers: 04.50.+h, 12.10.-g
The common expectation for the forthcoming search of gravitation waves is that gravity has a tensor nature and gravitational or metric waves ought to be quite different from vector electromagnetic waves. The unorthodox paradigm of curved three-space was successfully employed in the last century to explain the precise gravitational observations \[2-3\], but divorced electrodynamical and gravitational forces: electric charges do not disturb 3D geometry.
Nonetheless the similarity of Newton and Coulomb interactions may suggest a natural similarity or identity of gravitons and photons, which are responsible for these interactions. An alternative opportunity to explain gravitational observations is to keep flat three-space but to derive the pseudo-Riemannian four-interval from a nonlinear relation, $`ds^2=[d\tau (ds)]^2\delta _{ij}dx^idx^j`$. Only three-spaces $`x^i`$ with constant curvature, including flat space with $`\gamma _{ij}=\delta _{ij}`$, are compatible with the well-tested conservation of a system three-momentum at all space points.
The present scheme is based on material states of different charged objects in their proper four-spaces $`x__K^\mu `$ with different pseudo-Riemannian metrics. But evolution of material objects may be compared and observed within 3D intersections $`x__K^i=x^i`$ of the proper four-spaces, because all sub-spaces $`x__K^i`$ keep universal Euclidean geometry and may form laboratory three-space. Flat laboratory space $`x^i`$ and a universal time interval $`dt=\pm |dx^o|`$ are commom for all extended particle-field objects, while any curved four-space $`x__N^\mu `$ may be associated only with one selected object N. The accepted gravitation of the point charges in common curved four-space differs from the novel scheme, called vector electrogravity, where non-relativistic electrodynamic relations may be used to test gravitation of the extended charges and their point sources.
In order to prove joint roots for gravitational and electrodynamic fields one has to derive all variation equations in a joint vector form and test the derived dynamical equations in practice. Let all fields from external sources be based on the unified forming-up four-potentials $`a_{{}_{K}{}^{}\mu }`$ and contribute jointly into the proper tetrad of any selected object N,
$$e_{{}_{N}{}^{}\mu }^\alpha =\delta _\mu ^\alpha +\delta ^{\alpha o}\sqrt{1\delta _{ij}v__N^iv__N^j}U_\mu ^_N,$$
(1)
where $`U_\mu ^_N(x)`$ = $`__K^{{}_{K}{}^{}_{N}^{}}(Gm__K+m__N^1q__Nq__K)a_{{}_{K}{}^{}\mu }`$. One may verify that the proper pseudo-Riemannian tensor, $`g_{\mu \nu }^_N\eta _{\alpha \beta }e_{{}_{N}{}^{}\mu }^\alpha e_{{}_{N}{}^{}\nu }^\beta `$, has universal symmetry $`\gamma _{ij}^_Ng_{oi}g_{oj}g_{oo}^1g_{ij}`$ = $`\delta _{ij}`$, that corresponds to flat 3D sub-space under arbitrary electromagnetic and gravitational fields.
The same external electromagnetic, $`A_\mu ^_N__K^{{}_{K}{}^{}_{N}^{}}q__Ka_{{}_{K}{}^{}\mu }`$, and gravitational, $`B_\mu ^_NG__K^{{}_{K}{}^{}_{N}^{}}m__Ka_{{}_{K}{}^{}\mu }`$, fields determine the proper canonical four-momentum
$$P_{{}_{N}{}^{}\nu }m__Ng_{\mu \nu }^_Ndx__N^\mu /ds__N=m__N\delta _\nu ^\alpha V_\alpha +m__NU_\nu ^_N,$$
(2)
where $`\delta _\mu ^\alpha V_\alpha `$ = $`\{\beta ^1,\beta ^1v_i\}`$, $`\beta =\beta __N`$ = $`\sqrt{1\delta _{ij}v^iv^j}`$, $`c=1`$, $`v^i=v__N^i=dx__N^i/d\tau __N`$, $`ds__N^2=d\tau __N^2\delta _{ij}dx__N^idx__N^j`$.
The proper time rate $`d\tau __N=\beta ^1ds__N=\sqrt{g_{oo}^_N}(dx__N^og_i^_Ndx__N^i)`$ = $`(1+\beta U_o^_N)dx__N^o`$ \+ $`\beta U_i^_Ndx__N^i=dx__N^o+\beta U_\mu ^_Ndx__N^\mu =dx__N^o+\beta ^2U_\mu ^_NP__N^\mu m__N^1d\tau __N`$ depends on all external gravitational and electromagnetic fields,
$$\left(\frac{d\tau __N}{dt}\right)^2=\left(\frac{1+\beta (B_o^_N+m__N^1q__NA_o^_N)}{1\beta (B_i^_N+m__N^1q__NA_o^_N)v__N^i}\right)^2.$$
(3)
The gravitational dilation of the proper time rate $`d\tau __N`$ with respect to the laboratory time interval $`dt`$ coincides for weak fields in (3) with the similar result of general relativity . There are also observations \[4-6\] of electromagnetic time dilation-compression for charges. Vector electrogravity explains this time relativity by coupling gravity and electrodynamics in the ”old” pseudo-Riemannian four-space, but with electromechanical connections, when
$`g_{oo}=(1+\beta U_o)^2,g_{oi}=(1+\beta U_o)\beta U_i,g_{ij}=\beta ^2U_iU_j+\eta _{ij},`$
$`g^i=g^{oi}=\gamma ^{ij}g_j=g_i=g_{oi}g_{oo}^1=U_i(\beta ^1+U_o)^1`$
$`g^{oo}=g_{oo}^1g_ig^i=(1\beta ^2U_iU_j\delta ^{ij})(1+\beta U_o)^2,\gamma _{ij}=\gamma ^{ij}=g^{ij}=\delta _{ij}`$
$`P_\mu =m\{\beta ^1+U_o;\beta ^1v_i+U_i\}=m(\delta _\mu ^\alpha V_\alpha +U_\mu )=g_{\mu \nu }P^\nu `$
$`P^\mu =\{m(\beta ^1+U^o);P^i\}=m\{\beta ^1(U_o+U_iv^i)(1+\beta U_o)^1;\beta ^1v^i\}`$
$`P_\mu P^\mu =g_{oo}(P^og_iP^i)^2\delta _{ij}P^iP^j=P_o^2g_{oo}^1m^2\beta ^2v^2=m^2.`$ (4)
Proper four-space $`x__N^\mu `$, metric tensor $`g_{\mu \nu }^_N`$, time rate $`d\tau __N`$ and four-interval $`ds__N\pm \sqrt{g_{\mu \nu }^_Ndx^\mu dx^\nu }`$ may be introduced only for one selected object N and they do not coincide with similar proper functions of other objects. An ensemble of different charged and neutral objects may be described in common three-space exclusively due to universal geometry, $`\gamma _{ij}^_K=\delta _{ij}`$, of all their three-intervals $`dl__K`$. There is no universal geometry for all four-intervals, and no one rate $`d\tau __K`$ can be used as a universal time interval for an ensemble of interacting elements. One may however employ the common time interval, $`dt\pm \sqrt{\gamma _{oo}^_Kdx__K^odx__K^o}=\pm \sqrt{\delta _{oo}dx__K^odx__K^o}`$ = $`\pm |dx^o|`$, which is appropriate for all matter, $`dt=|dx^o|`$, and antimatter, $`dt=|dx^o|`$, due to the universal metric tensor $`\gamma _{oo}^_K=\delta _{oo}`$ of flat one-dimensional proper intervals $`|dx__K^o|`$. Laboratory evolution of matter is three-dimensional because only the universal flat intervals, $`dt__K`$ and $`dl__K`$, rather than the unique proper four-intervals $`ds__K`$, have a common sense for the total ensemble.
According to the tetrad (1) this proper four-space can take Euclidean metric in a local inertial reference system, where all external fields are absent or balanced, $`U_\mu ^_N(x__N^\nu )=0`$. The equivalency principle leads to the following relations,
$$\frac{DP_{{}_{N}{}^{}\mu }}{dt}=\frac{dx__N^\nu }{dt}\frac{P_{{}_{N}{}^{}\mu }}{x__N^\nu },$$
(5)
for the charged object N in its proper four-space $`x__N^\nu `$. The novel geodesic relations and dynamics of charges under the canonical conservation $`P_{{}_{N}{}^{}\mu }P__N^\mu =m__N^2`$ in arbitrary electromagnetic fields can be tested in experiments.
A new relation between $`ds__N^2g_{\mu \nu }^_Ndx__N^\mu dx__N^\nu `$, $`dl__N^2\delta _{ij}dx__N^idx__N^j`$, and $`dt__N^2\delta _{oo}dx__N^odx__N^o`$ may be derived, due to (1), from the equality $`ds__N^2(dx__N^o+\alpha __Nds__N)^2`$ \+ $`\delta _{ij}dx__N^idx__N^j`$, where $`\alpha __N\beta U_\mu ^_NP_{{}_{N}{}^{}\mu }m__N^1`$ $`(U_o^_N+U_i^_Nv^i)/(1+\beta U_o^_N)`$. The pseudo-Riemannian four-interval for charged matter reads
$$ds^2(\alpha __N)\left(\frac{\alpha __Ndx^o\pm \sqrt{(dx^o)^2dl^2(1\alpha __N^2)}}{(1\alpha __N^2)}\right)^2\frac{dt^2}{(1\alpha __N)^2}\frac{dl^2}{(1\alpha __N)},$$
(6)
when $`(1\alpha __N^2)dl^2/dt^21`$. Notice that there is no Schwarzschild’s divergence in (6) because $`\alpha __N<0`$ for pure gravitational potentials.
The four-interval $`s__N(\alpha __N)`$ depends on external gravitational and electromagnetic fields that can be verified in practice for both neutral and charged non-relativistic objects. Solutions of (6) with $`(\alpha __N)=GM/r1`$ can explain, for example, the measured planet perihelion precession under flat three-space.
One may verify from (4) or (1), that the metric tensor or its tetrad takes four field degrees of freedom due to the external four-potential $`U_\mu ^_N`$. Thus, the Hilbert variation with respect to ten ”independent” components of $`g_{\mu \nu }^_N`$ is not a well-defined procedure. There is no physical notion with ten degrees of freedom behind the metric tensor $`g_{\mu \nu }^_N`$, which is not an independent proper variable because $`g__N^{\mu \nu }P_{{}_{N}{}^{}\mu }P_{{}_{N}{}^{}\nu }=scalar`$.
Due to the variation of the proper particle-field action $`S__N`$,
$$\delta S__N=𝑑x^4\sqrt{g}T__N^{\mu \nu }\delta g_{\mu \nu }(P_{{}_{N}{}^{}\lambda })=𝑑x^4\sqrt{g}T__N^{\mu \nu }\frac{g_{\mu \nu }^_N}{e_{{}_{N}{}^{}\rho }^\alpha }\frac{e_{{}_{N}{}^{}\rho }^\alpha }{P_{{}_{N}{}^{}\lambda }}\delta P_{{}_{N}{}^{}\lambda },$$
(7)
gravitation has to be rewritten in terms of four-vector contraction of the Hilbert energy tensor $`T__N^{\mu \nu }`$. A basic dynamic equation, which determines the forming-up four-potential $`a_{{}_{N}{}^{}\mu }`$, takes a vector Maxwell-type form, $`T__N^{\mu \nu }P_{{}_{N}{}^{}\mu }`$ = 0, when $`\delta S__N`$ = 0 in (7). Wave solutions of this four-vector equation correspond to photons, which are responsible for both gravitation and electromagnetic interactions.
By taking the proper action $`S__N`$ = $`\sqrt{g}d^4x[P_{{}_{N}{}^{}\mu }i__N^\mu +(f__N^{\mu \nu }W_{\mu \nu }^_N/16\pi )]`$, with $`i__N^\mu (dx__N^\mu /dp__N)(g)^{1/2}\widehat{\delta }^4(s__N)𝑑p__N`$, $`f_{\mu \nu }^_N_\mu a_{{}_{N}{}^{}\nu }_\nu a_{{}_{N}{}^{}\mu }`$, and $`W_{\mu \nu }^_N_\mu P_{{}_{N}{}^{}\nu }_\nu P_{{}_{N}{}^{}\mu }`$, one finds the proper energy tensor
$$T__N^{\mu \nu }=\frac{P__N^\mu I__N^\nu (x)+P__N^\nu I__N^\mu (x)}{2}+\frac{W_{{}_{N}{}^{}\rho \lambda }(x)}{16\pi }[g__N^{\mu \nu }f__N^{\rho \lambda }2g__N^{\mu \rho }f__N^{\nu \lambda }2g__N^{\nu \rho }f__N^{\mu \lambda }],$$
(8)
where the Maxwell-type four-vector $`I__N^\mu i__N^\mu (4\pi )^1_\nu f__N^{\mu \nu }`$ may be associated with both kind of the extended homogeneous charges, $`m__N`$ and $`q__N`$, located on the proper light-cone $`s__N(x__N,\xi __N)=0`$. Variations of the scalar action $`S__N`$ with respect to three proper variables, $`x__N^\mu `$, $`P__N^\mu `$, and $`a_{{}_{N}{}^{}\mu }`$, lead to a system of vector equations, $`_\mu T__N^{\mu \nu }=0`$, $`T__N^{\mu \nu }P_{{}_{N}{}^{}\mu }=0`$, and $`_\mu W__N^{\mu \nu }=0`$, respectively, which determines dynamics of one selected object N. The tensor Einstein-type relation, $`T__N^{\mu \nu }=0`$, for all components of the symmetric energy tensor (8) appears in vector electrogravity only for potential (superfluid, gauge-invariant) states of the object without energy exchange and radiation, when $`W_{\mu \nu }^_N=0`$ and $`I__N^\mu =0`$.
Measurements of the electron beam bending by static Coulomb potential can be performed with high accuracy. These measurements could test the dynamic equations (4)-(6) and the basic relation $`P_{{}_{N}{}^{}\mu }P__N^\mu =m__N^2`$ for canonical four-momentum. Electromagnetic dilation-compression of time (3) provides one more opportunity to test the developed double unification (particle with field, gravity with electromagnetism) in the laboratory.
Vector electrogravity operates only with flat three-space for both gravity and electrodynamics. The vector approach to gravitation suggests the unified nature of gravitational and electromagnetic radiation, associated with photon waves. These vector waves can modulate only the four-space metric tensor $`g_{\mu \nu }^_N`$, but they are irrelevant to Euclidean metrics, $`\gamma _{ij}^_N\delta _{ij}`$, of laboratory three-space under all possible experiments.
R. Weiss, Rev. Mod. Phys. 71, S187 (1999).
C.W. Misner, K.S. Thorne, and J.A. Wheeler, Gravitation (San Francisco: Freeman, 1973).
C.M. Will, Theory and experiment in gravitational physics (Cambridge: Cambridge University Press, 1993), revised ed.
E.J. Saxl, Nature 203, 136 (1964).
M.A. Tamers, Nature 339, 588 (1989).
W.A. Barker,US Patent No 5,076,971 (1991). |
warning/0001/hep-ph0001112.html | ar5iv | text | # References
INSTITUT FÜR KERNPHYSIK, UNIVERSITÄT FRANKFURT
D - 60486 Frankfurt, August–Euler–Strasse 6, Germany
IKF–HENPG/01–00
Baryon Number Conservation
and
Statistical Production of Antibaryons
Mark I. Gorenstein<sup>a,</sup><sup>1</sup><sup>1</sup>1Permanent address: Bogolyubov Institute for Theoretical Physics, Kiev, Ukraine<sup>,</sup><sup>2</sup><sup>2</sup>2E–mail: goren@th.physik.uni-frankfurt.de, Marek Gaździcki<sup>b,</sup><sup>3</sup><sup>3</sup>3E–mail: marek@ikf.physik.uni–frankfurt.de and Walter Greiner<sup>a</sup><sup>4</sup><sup>4</sup>4E–mail: greiner@th.physik.uni-frankfurt.de
<sup>a</sup> Institut für Theoretische Physik, Universität Frankfurt, Germany
<sup>b</sup> Institut für Kernphysik, Universität Frankfurt, Germany
## Abstract
The statistical production of antibaryons is considered within the canonical ensemble formulation. We demonstrate that the antibaryon suppression in small systems due to the exact baryon number conservation is rather different in the baryon–free ($`B=0`$) and baryon–rich ($`B2`$) systems. At constant values of temperature and baryon density in the baryon–rich systems the density of the produced antibaryons is only weakly dependent on the size of the system. For realistic hadronization conditions this dependence appears to be close to $`B/(B+1)`$ which is in agreement with the preliminary data of the NA49 Collaboration for the $`\overline{p}/\pi `$ ratio in nucleus–nucleus collisions at the CERN SPS energies. However, a consistent picture of antibaryon production within the statistical hadronization model has not yet been achieved. This is because the condition of constant hadronization temperature in the baryon–free systems leads to a contradiction with the data on the $`\overline{p}/\pi `$ ratio in e<sup>+</sup>+e<sup>-</sup> interactions.
1. Introduction
Among the first models of multiparticle production in high energy interactions were statistical models . In the last decade a significant development of these models and the extension of the area of their applicability took place. The main reason for this is a surprising success of the statistical approach in reproducing new experimental data on hadron multiplicities in nuclear (A+A) and elementary (e<sup>+</sup>+e<sup>-</sup>, p+p, p+$`\overline{\mathrm{p}}`$) collisions . One of the important results of the analysis of hadron yield systematics at high energies (SPS and higher) done within the statistical models is the approximate independence of the temperature parameter $`T=160÷190`$ MeV from the system size and collision energy . This result can be attributed to the statistical character of the hadronization process.
The statistical models are based on the key assumption that all microscopic states of the system allowed by the conservation laws are equally probable. Calculations within these models are straightforward when the mean number of particles of interest is large and consequently it is enough to fulfill the conservation laws in the average sense, i.e., for the macroscopic state. This is achieved by the introduction of parameters: the temperature $`T`$ and chemical potentials $`\mu _i`$, which control correspondingly the average values of the system energy density and its material content (e.g., baryon number, strangeness and electric charge). In this case (grand canonical ensemble, $`g.c.e.`$) the mean particle multiplicities are just proportional to the volume $`V`$ of the system. The particle density and the ratio of the multiplicities of two different particles often used for the comparison with the data are volume independent.
This simple volume dependence is however not valid any more for a small system in which the mean particle multiplicity is low. In this case (canonical ensemble, $`c.e.`$) the material conservation laws should be imposed on each microscopic state of the system. This condition introduces a significant correlation between particles which carry conserved charges. The correlation reduces the effective number of degrees of freedom and consequently leads to the $`c.e.`$ suppression of the ’charged’ particle multiplicity when compared with the result of the calculations done within $`g.c.e.`$ A ’neutral’ particle (a particle which does not carry conserved charges) does not feel the $`c.e.`$ suppression in the small system, its mean multiplicity remains proportional to the volume. Therefore, a different dependence on the system volume is expected within statistical models for ’charged’ and ’neutral’ particles. The magnitude of the $`c.e.`$ suppression of the ’charged’-’anticharged’ pair creation increases with the mass of the lightest hadron needed to compensate the particle charges. Therefore, a strong $`c.e.`$ suppression may be expected for antibaryon production as the mass of the lightest baryon is $`m938`$ MeV which is much larger than the value of the temperature parameter found in the hadronization models. The $`c.e.`$ suppression is still rather essential for strange particle production (see, e.g., ).
The above expectation seems to be violated by the preliminary data on antiproton production presented recently by the NA49 Collaboration . The $`\overline{p}/\pi `$ ratio is found to be approximately the same for p+p interactions and central Pb+Pb collisions at 158 A$``$GeV (antiproton scaling). Thus the hadronization volume increases but the effect of the $`c.e.`$ suppression is not observed. The $`c.e.`$ suppression can be expected in p+p interactions because of the key difference between antiproton and pion: the antiproton carries baryon number in addition to the electric charge carried by both particles. In order to compensate the electric charge of a produced particle it is enough to create an additional charged pion. As the pion mass is much smaller than the nucleon mass, a significantly stronger $`c.e.`$ suppression is expected for antiproton production than for the pion production, consequently it should lead to a strong violation of the experimentally observed scaling. Thus the crucial question is whether the antiproton scaling can be understood within the statistical model of hadron production in which the condition of exact baryon number conservation is imposed. We note that the antiproton multiplicities in high energy collisions were shown to approximately agree with the predictions of statistical models . However different versions of the models with different parameters were used to fit various sets of data. Thus the question whether consistent description of the antiproton data in the statistical model is possible is still opened.
The importance of the exact treatment of the material conservation laws within statistical models of strong interactions was first pointed out by Hagedorn (see also Refs. ). Subsequently a complete treatment has been developed (see, e.g., and references therein) and applied to analyze the hadron yields in elementary collisions . In this letter we derive explicit analytical formulae to study the role of the exact material conservation laws within the statistical model of hadronization and we use them to discuss the antiproton scaling observed experimentally in Pb+Pb collisions at 158 A$``$GeV . We also discuss the data on the $`\overline{p}/\pi `$ ratio in e<sup>+</sup>+e<sup>-</sup> interactions within the statistical hadronization model.
2. Model formulation
Let us consider the system of baryons ’$`b`$’ and antibaryons ’$`a`$’ with total baryon number $`B`$ as the Boltzmann ideal gas in the volume $`V`$, at temperature $`T`$. The c.e. partition function is
$`Z(T,V,B)`$ $`=`$ $`{\displaystyle \underset{N_b^{(1)},N_a^{(1)}}{\overset{\mathrm{}}{}}}\mathrm{}{\displaystyle \underset{N_b^{(j)},N_a^{(j)}=0}{\overset{\mathrm{}}{}}}\mathrm{}\delta _K\left[B{\displaystyle \underset{j}{}}(N_b^{(j)}N_a^{(j)})\right]`$
$`\times `$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{(\lambda _b^{(j)}z_j)^{N_b^{(j)}}}{N_b^{(j)}!}}{\displaystyle \frac{(\lambda _a^{(j)}z_j)^{N_a^{(j)}}}{N_a^{(j)}!}},`$
where the index $`j`$ runs over all (non-strange) baryon states $`N,\mathrm{\Delta },N^{},\mathrm{}`$, and the single baryon (antibaryon) partition function reads
$`z_j=z_j(T,V)`$ $`=`$ $`{\displaystyle \frac{g_jV}{(2\pi )^3}}{\displaystyle d^3k\mathrm{exp}[(k^2+m_j^2)^{1/2}/T]}=`$ (2)
$`=`$ $`{\displaystyle \frac{g_jV}{2\pi ^2}}Tm_j^2K_2(m_j/T)Vf_j(T).`$
The baryon mass and the baryon degeneracy factor are denoted here by $`m_j`$ and $`g_j`$, respectively. Auxiliary parameters $`\lambda _b^{(j)}`$ and $`\lambda _a^{(j)}`$ are introduced in order to calculate the mean number of baryons and antibaryons and they are set to unity in the final formulae. By expressing $`\delta _K`$ as
$$\delta _K(n)=\frac{1}{2\pi }_0^{2\pi }𝑑\varphi e^{in\varphi },$$
Eq. (S0.Ex1) becomes
$`Z(T,V,B)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi e^{iB\varphi }{\displaystyle \underset{j}{}}{\displaystyle \underset{N_b^{(j)}=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{N_a^{(j)}=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\lambda _b^{(j)}z_je^{i\varphi })^{N_b^{(j)}}}{N_b^{(j)}!}}{\displaystyle \frac{(\lambda _a^{(j)}z_je^{i\varphi })^{N_a^{(j)}}}{N_a^{(j)}!}}=`$ (3)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}𝑑\varphi e^{iB\varphi }\mathrm{exp}\left[{\displaystyle \underset{j}{}}z_j(\lambda _b^{(j)}e^{i\varphi }+\lambda _a^{(j)}e^{i\varphi })\right].`$
This form of the $`c.e.`$ partition function allows one to derive the mean numbers of baryons and antibaryons
$$N_b^{(j)}=\left(\frac{\mathrm{log}Z}{\lambda _b^{(j)}}\right)_{\lambda _b=\lambda _a=1}=z_j\frac{Z(T,V,B1)}{Z(T,V,B)},$$
(4)
$$N_a^{(j)}=\left(\frac{\mathrm{log}Z}{\lambda _a^{(j)}}\right)_{\lambda _p=\lambda _a=1}=z_j\frac{Z(T,V,B+1)}{Z(T,V,B)}.$$
(5)
For $`\lambda _b=\lambda _a=1`$ the partition function (3) can be presented as the modified Bessel function
$$Z(T,V,B)=\frac{1}{2\pi }_0^{2\pi }𝑑\varphi e^{iB\varphi }\mathrm{exp}(2z\mathrm{cos}\varphi )=I_B(2z),$$
(6)
where $`z_jz_j`$. This yields final expressions for the mean number of baryons and antibaryons
$$N_b^{(j)}=z_j\frac{I_{B1}(2z)}{I_B(2z)},N_a^{(j)}=z_j\frac{I_{B+1}(2z)}{I_B(2z)}.$$
(7)
As the exact baryon number conservation is imposed on each microscopic state it is evidently fulfilled also by the average values (7):
$$N_bN_a\underset{j}{}N_b^{(j)}\underset{j}{}N_a^{(j)}=B,$$
(8)
as indeed can be easily seen from the identity $`I_{n1}(x)I_{n+1}(x)=2nI_n(x)/x`$ . Eq. (7) is valid for all combinations of $`B`$ and $`z`$ values. For a specific case of $`B=0`$ in the nucleon–antinucleon gas (i.e., no resonances included) our results (7) are reduced to the result of Rafelski and Danos .
The $`c.e.`$ expressions for the mean number of baryons and antibaryons can be further simplified for the two limiting cases: $`z1`$ (small systems) and $`z1`$ (large systems). Using the representation of $`I_n`$ as the infinite series
$$I_n(2z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^{n+2k}}{k!(n+k)!},$$
one obtains for small systems
$$N_b^{(j)}B\frac{z_j}{z}+\frac{z_jz}{B+1}+o(z_jz^3),N_a^{(j)}\frac{z_jz}{B+1}+o(z_jz^3).$$
(9)
The dependence $`N_aV^2/(B+1)`$ is therefore observed from Eq. (9) for the antibaryon yield in small systems. Such a dependence can be intuitively understood from the kinetic picture of the baryon–antibaryon pair creation and annihilation. Let’s consider the time dependence $`N_a(t)`$ (the time averaging of $`N_a(t)`$ should reproduce our statistical average $`N_a`$). At each time moment the value of $`N_a(t)`$ equals 0 or 1 (configurations with $`N_a(t)2`$ can be safely neglected as $`N_a1`$). The kinetics of the evolution of $`N_a(t)`$ is defined by the frequency $`\omega `$ of the baryon–antibaryon pair creation and the life–time $`\mathrm{\Delta }t`$ of the produced baryon–antibaryon pair: $`N_a(t)=\omega \mathrm{\Delta }t.`$ As a baryon–antibaryon pair is locally produced in any point of the system we have $`\omega V`$. Being produced as baryon–antibaryon pair the antibaryon can be then locally annihilated by any baryon existing in the system. Therefore, $`\mathrm{\Delta }t1/n_B`$ ($`n_B`$ is the density of baryons). $`n_B=1/V`$ for $`B=0`$ (as only one baryon exists in the volume $`V`$), and $`n_B=(B+1)/V`$ if $`B`$ baryons are present in the system before the baryon–antibaryon pair creation. These lead to the dependence of the antibaryon number on $`V`$ and $`B`$ as given by Eq. (9).
For large systems ($`z1`$) the c.e. becomes equivalent to the g.c.e. where the partition function and the average number of baryons and antibaryons are calculated as
$$Z(V,T,\mu _B)=\underset{b}{}\mathrm{exp}\left(\frac{\mu _Bb}{T}\right)Z(T,V,b)=\mathrm{exp}\left(ze^{\mu _B/T}+ze^{\mu _B/T}\right),$$
(10)
$$N_b^{(j)}=z_j\mathrm{exp}(\mu _B/T),N_a^{(j)}=z_j\mathrm{exp}(\mu _B/T).$$
(11)
Here $`\mu _B`$ is a baryon chemical potential which using Eq. (8) is defined for $`B=b0`$ as:
$$\mathrm{exp}(\mu _B/T)=\frac{B}{2z}+\sqrt{1+\left(\frac{B}{2z}\right)^2}.$$
(12)
Note that the function $`f_j=f_j(T)`$ introduced in Eq. (2) has the physical meaning of the density of the $`j`$-th baryon and antibaryon in the g.c.e. formulation for $`\mu _B=0`$. Using the uniform asymptotic expansion of the modified Bessel functions at $`n\mathrm{}`$
$$I_n(nx)\frac{1}{\sqrt{2\pi n}}\frac{\mathrm{exp}(n\eta )}{(1+x^2)^{1/4}}\left[1+o\left(\frac{1}{n}\right)\right];\eta \sqrt{1+x^2}+\mathrm{ln}\frac{x}{1+\sqrt{1+x^2}},$$
(13)
results (11) and (12) for the $`g.c.e.`$ are also easy to obtain from the $`c.e`$. (7) in the thermodynamical limit $`V\mathrm{},B\mathrm{}`$ with $`B/V\rho _B=const(V)`$.
In order to remove a ’trivial’ linear dependence of the particle multiplicities on the system volume it is convenient to make a comparison between the particle ratios from the model and the experimental data. In the statistical model the multiplicity of any ’neutral’ meson state $`M`$ is just proportional to the volume ($`N_M=z_M=Vf_M(T)`$) for both small ($`c.e.`$) and large ($`g.c.e.`$) systems. Therefore, the system volume dependence of the ratio $`N_a/N_M`$ at fixed temperature is the same as for the antibaryon density. From Eqs. (7) and (11) the antibaryon densities for the c.e. and g.c.e. are equal to
$`{\displaystyle \frac{N_a^{(j)}}{V}}_{ce}`$ $``$ $`n_a^{(j)}_{ce}=f_j{\displaystyle \frac{I_{B+1}(2z)}{I_B(2z)}}=f_j{\displaystyle \frac{I_{B+1}(xB)}{I_B(xB)}}f_j{\displaystyle \frac{x}{2}}{\displaystyle \frac{B}{B+1}},`$ (14)
$`{\displaystyle \frac{N_a^{(j)}}{V}}_{gce}`$ $``$ $`n_a^{(j)}_{gce}=f_j\mathrm{exp}(\mu _B/T)=f_j{\displaystyle \frac{x}{1+\sqrt{1+x^2}}},`$ (15)
where $`f_jf_j`$, $`x2z/B=2f/\rho _B`$, and the last approximation in Eq. (14) is valid for small systems only. Note that introducing the variable $`x`$ we have transformed the finite size $`V`$-dependence of the $`c.e.`$ density (14) into its dependence on the baryon number $`B`$. Eqs. (14,15) give us the primary thermal density for all individual antibaryon states $`j`$. Each non-strange resonance (anti)baryon state decays finally into (anti)nucleon plus meson(s). Therefore, the total (primary plus resonance decay) antinucleon density equals to the total thermal antibaryon density, $`n_a=_jn_a^{(j)}`$ and is given by Eqs. (14,15) with the substitution of $`f_j`$ by a sum $`f=_jf_j`$.
For the purpose of the following discussion we define a canonical suppression factor
$$F_{cs}\frac{\left(n_a\right)_{ce}}{\left(n_a\right)_{gce}}.$$
(16)
It quantifies the antinucleon suppression due to the exact baryon number conservation. We note also that the suppression factor $`F_{cs}`$ (16) is the same for any individual antibaryon state.
3. Discussion
The results derived in the previous section are used here to discuss antibaryon production in high energy collisions.
In the $`B=0`$ case the baryon and antibaryon densities are equal and Eqs. (14) and (15) yield
$$n_a^{(j)}_{ce}=n_b^{(j)}_{ce}=f_j\frac{I_1(2z)}{I_0(2z)}f_jz,n_a^{(j)}_{gce}=n_b^{(j)}_{gce}=f_j,$$
(17)
where the approximation for the c.e. density is valid for small system only. The canonical suppression factor (16) for $`B=0`$ is equal to
$$F_{cs}^0=\frac{I_1(2z)}{I_0(2z)}fV.$$
(18)
The approximation in Eq. (18) is valid for small system only ($`zfV1`$).
The behavior of the canonical suppression factor $`F_{cs}^0`$ (18) is shown by the solid lines in Fig. 1 for $`T=160`$ MeV, 170 MeV and 180 MeV, assuming that $`f`$ is the sum of $`f_j`$ over all non-strange baryons. The lines start from $`V=`$ 5 fm<sup>3</sup>, which is approximately equal to the estimate of the hadronization volume for $`e^++e^{}`$ interactions at $`\sqrt{s}=29`$ GeV . One observes (see Fig. 1) a strong $`c.e.`$ suppression of the (anti)baryon density. For $`T=160`$ MeV the (anti)baryon density increases by a factor of 10 from its value at $`V=5`$fm<sup>3</sup> to its $`V\mathrm{}`$ g.c.e. limit. For the small systems the (anti)baryon density increases approximately linearly with $`V`$, i.e., the (anti)baryon multiplicity for the small systems is proportional to $`V^2`$. The $`c.e.`$ suppression becomes less pronounced and the volume region with linear increase of the (anti)baryon density is reduced for increasing temperature.
Let us now turn to the antibaryon production in baryon rich system. In the analysis of data on particle multiplicities in p+p, p+A and A+A collisions one usually assumes that all participating nucleons in the collisions (wounded nucleons) take part in the statistical hadronization of the system. It means that in the analysis of the NA49 results on antiprotons from p+p interactions to central Pb+Pb collisions at 158 A$``$GeV we should study statistical systems with $`2B400`$. It was found that the mean multiplicity of pions per wounded nucleon increases (at the SPS collision energies) only by about 20% when going from p+p interactions to central Pb+Pb collisions . The pion to baryon ratio in the statistical model is determined by two parameters: the temperature and baryon density. Thus as the temperature is found to be constant ($`T=175\pm 15`$ MeV) we conclude that the baryon density at hadronization in nuclear collisions at 158 A$``$GeV is also approximately constant.
Therefore, for the comparison with the NA49 results we study the evolution of the antibaryon density with increasing net baryon number $`B`$ at $`T=const`$ and $`\rho _B=const`$. The $`c.e.`$ suppression factor (16) is found at these conditions from Eqs. (14,15)
$$F_{cs}^B=\frac{1+\sqrt{1+x^2}}{x}\frac{I_{B+1}(xB)}{I_B(xB)};x\frac{2f}{\rho _B}.$$
(19)
Its $`B`$-dependence is plotted in Fig. 2 for several different values of the parameter $`x`$. Note that our assumption $`T=const`$ and $`\rho _B=const`$ for statistical hadronization at different values of $`B`$ can be substituted by a weaker one, $`x=const`$. From Fig. 2 one observes that the $`c.e.`$ suppression of antibaryon density becomes stronger at high baryon density (i.e., small $`x`$). For $`x<1`$ the $`c.e.`$ suppression $`F_{cs}^B`$ (19) becomes close to its $`x0`$ limit:
$$F_{cs}^B=\frac{B}{B+1}.$$
(20)
Eq. (20) shows that the strongest $`c.e.`$ suppression of the antibaryon density is for the $`B=2`$ (nucleon–nucleon interactions) case and it leads to the suppression factor of 2/3. This moderate effect of $`c.e.`$ suppression is in strong contrast with the large $`c.e.`$ suppression (i.e., $`F_{cs}^01`$) in the baryon–free system. A mathematical reason of this very different behavior for $`B=0`$ and $`B2`$ (with $`\rho _B=const(V)`$) is due to the fact that in the latter case both the order of the modified Bessel functions and their arguments are dependent on $`B`$ (i.e., on $`V`$) whereas in the $`B=0`$ case only the argument increases with $`V`$.
The presence of non-zero baryon number $`B>0`$ has a twofold effect on antibaryon production. First, it suppresses the production of antibaryons: the additional factors $`\mathrm{exp}(\mu _B/T)=x/(1+\sqrt{1+x^2})<1`$ and $`1/(B+1)<1`$ appear respectively in the ’large’ and ’small’ systems for the antibaryon density in comparison with the $`B=0`$ case. On the other hand, the $`c.e.`$ suppression effect due to the exact baryon number conservation becomes smaller: at fixed $`T`$ and $`V`$ the following inequality is always valid, $`F_{cs}^B>F_{cs}^0`$. For fixed $`B>0`$ the $`c.e.`$ suppression of antibaryons becomes smaller when $`\rho _B`$ decreases and it disappears completely (i.e., $`F_{cs}^B1`$) in the limit $`\rho _B0`$ (and respectively $`V\mathrm{}`$ in order to keep the $`B`$ value fixed). This is because the total number of baryon–antibaryon pairs becomes large due to large $`V`$. Note that in this case the last approximation in Eq. (14) is no more valid. Instead one should use the large argument asymptotic of the modified Bessel functions.
Thus for $`B2`$ systems at constant $`x=2f/\rho _B`$ the $`c.e.`$ suppression factor $`F_{cs}^B`$ (19) ranges between 2/3 and 1 for $`x1`$ and between $`(11/4x)`$ and 1 for $`x1`$.
Previous analyses of hadron production at the CERN SPS indicate large baryon densities at hadronization. The typical values of $`T170`$ MeV and $`\mu _B250`$ MeV found for Pb+Pb collisions lead to the estimate $`x=2f/\rho _B=sh^1(\mu _B/T)0.5`$. As seen from Fig. 2 the $`c.e.`$ suppression factor $`F_{cs}^B`$ (19) for this ’small’ value of $`x`$ is close to its limiting pattern $`B/(B+1)`$ (20). In Fig. 3 the NA49 results on the $`\overline{p}/\pi `$ ratio in p+p and Pb+Pb collisions at 158 A$``$GeV are compared with this limiting pattern. From this comparison we conclude that the model of statistical production of antiprotons at hadronization in baryon–rich system correctly reproduces the observed antiproton scaling.
Let us return again to the case of the baryon–free system. The statistical model calculations for e<sup>+</sup>+e<sup>-</sup> and p+$`\overline{\mathrm{p}}`$ interactions include large $`c.e.`$ suppression effects. As discussed in the introduction we assume that the hadronization temperature reflects a universal property of the hadronization process and therefore should be collision energy independent. In the case of e<sup>+</sup>+e<sup>-</sup> interactions the hadronization volume is small and therefore one expects approximate proportionality to $`V^2`$ of the multiplicity of nucleon–antinucleon pairs but only a linear increase with $`V`$ of the pion multiplicity. Therefore, the $`\overline{p}/\pi `$ ratio ratio calculated within the model increases linearly with increasing pion multiplicity. However experimental data contradict this expectation of the statistical model. The $`\overline{p}/\pi `$ ratio which is plotted in Fig. 4 as a function of pion multiplicity for e<sup>+</sup>+e<sup>-</sup> interactions at different energies, $`\sqrt{s}=14÷91`$ GeV, is approximately constant.
Within the discussed statistical hadronization model one can try to solve the problem by assuming an increase of the temperature $`T`$ with decreasing volume $`V`$. The function $`f(T)`$ strongly increases with $`T`$ which allows to compensate the $`c.e.`$ suppression effect (to keep the (anti)nucleon density, $`n_af^2(T)V`$, constant) for moderate (of about 10 MeV) changes of $`T`$. This indeed is observed in the fit results of the statistical model for the e<sup>+</sup>+e<sup>-</sup> and p+$`\overline{\mathrm{p}}`$ data: the increase of the volume is always accompanied with the decrease of the temperature parameter . Thus one may argue that the hadronization condition $`T=const`$ has to be substituted by a different criterion which should explain the decrease of the temperature with increasing size of the system in e<sup>+</sup>+e<sup>-</sup> and p+$`\overline{\mathrm{p}}`$ interactions. The constant energy per particle was recently discussed as a chemical freeze–out condition . The detailed study of this question is, however, outside of the scope of the present paper. We note only that the statistical production of heavy particles (e.g., $`J/\psi `$ mesons ) is very sensitive to the temperature parameter. Their yields, therefore, can be used to clarify the problem.
4. Summary
The role of baryon number conservation in the calculations of antibaryon multiplicity within statistical model of hadronization was investigated. We derived explicit analytical formulae for the antibaryon multiplicity in baryon–free and baryon–rich small and large systems. This formalism was further used to discuss antiproton scaling observed experimentally in A+A collisions. The statistical model with constant hadronization temperature correctly reproduces the weak dependence of the $`\overline{p}/\pi `$ ratio on the system size in p+p and nuclear collisions at the CERN SPS energy. A description of the ratio of $`J/\psi `$ mesons to pions within the statistical hadronization model requires also a constant temperature parameter in p+p and A+A collisions at the CERN SPS. However, the same model with $`T=const`$ does not give a natural explanation of the approximate independence of the $`\overline{p}/\pi `$ ratio of collision energy in e<sup>+</sup>+e<sup>-</sup> interactions. Therefore, a consistent description of hadron production within the statistical hadronization model has not yet been achieved.
Acknowledgements
We thank F. Becattini, K. Bugaev, J. Cleymans, A. Korol, A. Kostyuk, I. Mishustin, St. Mrówczyński, L. Neise P. Seyboth and H. Stöcker for fruitful discussions. We acknowledge the financial support of BMBF and DFG, Germany. |
warning/0001/physics0001010.html | ar5iv | text | # Spectrum of atomic radiation at sudden perturbation
A general expression for the spectrum of photons emitted by atom at sudden perturbation is obtained. Some concrete examples of application of the obtained result are considered. The conclusion about the coherence of radiation of the atomic electrons under the such influences is made.
PACS number: 32.30.\*
It is known many examples when the excitation or ionization of atoms occurs as result of the action of sudden perturbations. First of all these are atomic excitation or ionization in the nuclear reactions . For example in $`\beta `$-decay of nucleus, when the fast $`\beta `$-electron’s escape is perceived by atomic electrons as a sudden changing of nuclear charge or in neutron impact with nucleus, when the sudden of momentum transfer to the nucleus occurs etc. The sudden approximation can be used for consideration multielectron transition in complex atoms, when transition occurring in internal shells, are perceived by relatively slow electrons of external shells as instantaneous (see ). As a result of action of sudden perturbation can be considered inelastic processes in the collisions of fast multicharged ions with atoms \[6 - 12\] and in the collisions of charged particles with highly-excited atoms . After action of sudden perturbation, the excited atom can relax with radiation of photons belonging to known spectrum of isolated atom. However, if sudden perturbation causes the change of velocities of atomic electrons, atom can radiate during the action of perturbation. Classical analogue of such a problem is the radiation of a free electron under the sudden changing of velocity. In many practically important cases perturbation is not sufficiently small to use a perturbation theory. However the situations when the time of action of perturbation is considerably less than the characteristic atomic time that enables one to solve the problem without restricting the value of perturbation (see for instance \[9,15-17\]).
Thus, it is necessary to state a general problem on the spectrum of photons emitted by atom during the time of action of sudden perturbation, i.e. - on the spectrum of photons emitted simultaneously by all atomic electrons as a result of action of perturbation. In this paper we derive a general expression for the spectrum of photons emitted by the atom under sudden perturbation and apply this result to some concrete processes.
Consider ”collision” type sudden perturbation , when the perturbation $`V(t)V(𝐫_a,t)`$ (where $`𝐫_a`$ are the coordinates of atomic electrons) acts only during the time $`\tau `$, which is much smaller than the characteristic period of unperturbed atom, describing by Hamiltonian $`H_0`$. To be definite we will assume that $`V(t)`$ is not equal zero near $`t=0`$ only. Then in the exact solution of Schrödinger equation ( atomic units are used throughout in this paper)
$$i\frac{\psi }{t}=(H_0+V(t))\psi $$
one can neglect by evolution of $`\psi `$ (during the time $`\tau `$) caused by unperturbed Hamiltonian $`H_0`$. Therefore the transition amplitude of atom from the initial state $`\phi _0`$ to a final state $`\phi _n`$, as a result of actions of sudden perturbation $`V(t)`$, has the form :
$$a_{0n}=\phi _nexp(i\underset{\mathrm{}}{\overset{+\mathrm{}}{}}V(t)𝑑t)\phi _0,$$
(1)
where $`\phi _0`$ and $`\phi _n`$ belong to the full set of orthonormalized eigenfunctions of the unperturbated Hamiltonian $`H_0`$, i.e. $`H_0\phi _n=ϵ_n\phi _n`$.
Thus in the sudden perturbation approximation the evolution of the initial state has the form
$$\psi _0(t)=exp(i\underset{\mathrm{}}{\overset{t}{}}V(t^{})𝑑t^{})\phi _0,$$
(2)
where $`\psi _0(t)`$ satisfies the equation
$$i\frac{\psi _0(t)}{t}=V(t)\psi _0(t),$$
(3)
and $`\psi _0(t)\phi _0`$ under $`t\mathrm{}`$. Let’s introduce full and orthonormal set of functions
$$\mathrm{\Phi }_n(t)=exp(i\underset{t}{\overset{+\mathrm{}}{}}V(t^{})𝑑t^{})\phi _n,$$
(4)
obeying eq. (3), and $`\mathrm{\Phi }_n(t)\phi _n`$ at $`t+\mathrm{}`$. Obviously the amplitude (1) can be rewritten as
$$a_{0n}=\mathrm{\Phi }_n(t)\psi _0(t).$$
Therefore the one photon radiation amplitude can be calculated in the first order of perturbation theory (as a corrections to the states (2) and (4)) over the interaction of atomic electrons with electromagnetic field
$$W=\underset{a,𝐤,\sigma }{}\left(\frac{2\pi }{\omega }\right)^{\frac{1}{2}}𝐮_{𝐤\sigma }(a_{𝐤\sigma }^+e^{i\mathrm{𝐤𝐫}_a}+a_{𝐤\sigma }e^{i\mathrm{𝐤𝐫}_a})\widehat{𝐩}_a,$$
where $`a_{𝐤\sigma }^+`$ and $`a_{𝐤\sigma }`$ are the creation and annihilation operators of the photon with a frequency $`\omega `$, momentum k and polarization $`\sigma ,(\sigma =1,2),𝐮_{𝐤\sigma }`$ are the unit vectors of polarization, $`𝐫_a`$ are the coordinates of atomic electrons ($`a=1,..,Z_a`$), here $`Z_a`$ is the number of atomic electrons, $`\widehat{𝐩}_a`$ are the momentum operators of atomic electrons. Then in the dipole approximation the amplitude of emission of photon with simultaneous transition of atom from the state $`\phi _0`$ to a state $`\phi _n`$ has the form
$$b_{0n}(\omega )=i\left(\frac{2\pi }{\omega }\right)^{\frac{1}{2}}𝐮_{𝐤\sigma }\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑te^{i\omega t}\mathrm{\Phi }_n(t)\underset{a}{}\widehat{𝐩}_a\psi _0(t).$$
Integrating this expression by parts over the time and omitting the terms vanishing (at $`t\pm \mathrm{}`$) in turning off the interaction with electromagnetic field we have
$`b_{0n}(\omega )=i\left({\displaystyle \frac{2\pi }{\omega }}\right)^{\frac{1}{2}}𝐮_{𝐤\sigma }{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}dt{\displaystyle \frac{e^{i\omega t}}{\omega }}\times `$
$`\times \phi _n{\displaystyle \underset{a}{}}{\displaystyle \frac{V(t)}{𝐫_a}}exp\left(i{\displaystyle \underset{\mathrm{}}{\overset{+\mathrm{}}{}}}V(t^{^{}})𝑑t^{^{}}\right)\phi _0.`$ (5)
Summing $`b_{0n}(\omega )^2`$ over polarization and integrating over the photon’s emission angles and summing, after this, over all final states of the atom $`\phi _n`$, we find the total radiation spectrum
$$\frac{dW}{d\omega }=\frac{2}{3\pi }\frac{1}{c^3\omega }\phi _0\underset{a}{}\frac{\stackrel{~}{V}^{}(\omega )}{𝐫_a}\underset{b}{}\frac{\stackrel{~}{V}(\omega )}{𝐫_b}\phi _0,$$
(6)
where c = 137 a.u. is the speed of light,
$$\stackrel{~}{V}(\omega )=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}V(t)e^{i\omega t}𝑑t.$$
(7)
Thus we have obtained the radiation spectrum of atom during the time of sudden perturbation $`V(t)`$.
As an application we consider the radiation spectrum of atom in the sudden transmission of momentum $`𝐩`$ to the atomic electrons when $`V(t)`$ has the (widely used for collision problems) form
$$V(t)=𝐟(t)\underset{a}{}𝐫_a,𝐩=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}𝑑t𝐟(t),$$
(8)
and $`𝐟(t)`$ is the perturbing force which not depends on $`𝐫_a`$ and interacts during a time $`\tau `$ that is considerable less than the characteristic periods of the unperturbed atom. The total radiation spectrum (6) in this case has the form
$$\frac{dW}{d\omega }=\frac{2}{3\pi }\frac{1}{c^3\omega }\stackrel{~}{𝐟}(\omega )^2Z_a^2,$$
(9)
where $`\stackrel{~}{𝐟}(\omega ),`$ is the Fourier transform of the functions $`𝐟(t)`$, defined according to (7), $`Z_a`$ is the number of atomic electrons. In this case the spectrum coincides (after producting to $`\omega `$) with the radiation spectrum of the classical particle with mass equal to electron’s one and with charge $`Z_a`$, moving in the field of homogeneous forces $`𝐟(t)`$. This gives us the information about the value of the spectrum (9). Since $`𝐟(t)`$ $`0`$ just during the time $`\tau `$ , and the spectrum (9) is proportional to $`\stackrel{~}{𝐟}(\omega )^2`$, only the photons belonging to continuum with characteristic frequencies $`\omega 1/\tau `$ can be emitted by atom.
Analogously one can consider the radiation of atom in the ”switching” type sudden perturbation (we use the classification of sudden perturbations introduced in ).
Formula (5) allows one to obtain the spectrum of photons in the transition of atom from the state $`\phi _0`$ to a state $`\phi _n`$ under the influence of perturbation (8):
$$\frac{dw_{0n}}{d\omega }=\frac{2}{3\pi }\frac{1}{c^3\omega }\stackrel{~}{𝐟}(\omega )^2Z_a^2\phi _nexp(i𝐩\underset{a}{}𝐫_a)\phi _0^2.$$
(10)
Here $`dW/d\omega =_ndw_{0n}/d\omega `$, where $`_n`$ means summing over the complete set of atomic states. Formula (10) allows one to express the relative contribution of transitions with excitation to an arbitrary state $`\phi _n`$ to the total spectrum (9)
$$\frac{dw_{0n}/d\omega }{dW/d\omega }=\phi _nexp(i𝐩\underset{a}{}𝐫_a)\phi _0^2$$
via the well known inelastic atomic formfactors $`\phi _nexp(i𝐩\underset{a}{}𝐫_a)\phi _0`$.
In the most simple case of instantaneous transferring to atomic electrons the momentum $`𝐩`$, when in (8) $`𝐟(t)=𝐩\delta (t),`$ where $`\delta (t)`$ is the Dirac $`\delta `$-function, then $`\stackrel{~}{𝐟}(\omega )=𝐩`$ and spectrum (9) coincides, after producting to $`\omega `$, with the radiation spectrum of free classical particle with charge $`Z_a`$, which takes (suddenly) a velocity $`𝐩.`$
As an another example we give the radiation spectrum in the influence of pulse having the Gausian form
$$𝐟(t)=𝐟_0exp(\alpha ^2t^2)cos(\omega _0t),$$
respectively
$$\stackrel{~}{V}(\omega )=\frac{\sqrt{\pi }}{2\alpha }𝐟_\mathrm{𝟎}\underset{a}{}𝐫_a\left\{exp\left[\frac{(\omega \omega _0)^2}{4\alpha ^2}\right]+exp\left[\frac{(\omega +\omega _0)^2}{4\alpha ^2}\right]\right\}.$$
Therefore the radiation spectrum has the form
$$\frac{dW}{d\mathrm{\Omega }}=\frac{f_0^2}{6\mathrm{\Omega }c^3\alpha ^2}\left\{exp\left[(\mathrm{\Omega }+\mathrm{\Omega }_0)^2\right]+exp\left[(\mathrm{\Omega }\mathrm{\Omega }_0)^2\right]\right\}Z_a^2,$$
where for the sake of convenience the frequencies $`\mathrm{\Omega }=\omega /(2\alpha )`$ and $`\mathrm{\Omega }_0=\omega _0/(2\alpha )`$ are introduced.
One should note an important generality of radiation at sudden perturbation, namely, the radiation intensity for the multielectron atoms is proportional to the square of the number of atomic electrons. This fact allows one to conclude on the coherence of radiation of atomic electrons under such type influences.
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warning/0001/quant-ph0001101.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In semi-classical treatments we construct approximate solutions of a quantum mechanical problem from the knowledge of the solutions of the corresponding classical problem. The usual WKB procedure yields excellent results sufficiently far away from turning points. Yet at turning points this approximation leads to singularities. This occurs for all energies below the maximal value of the potential. Furthermore degenerate situations occur at maxima of this potential. Some kind of regularisation (uniformisation) is used near such points. The most common ad hoc solution uses Airy functions near turning points and Pearcy functions near maxima .
The purpose of this paper is to present a semi-classical global approximation without singularities. To achieve this we take advantage of the following facts:
First, the solution of a problem linear in momentum such as $`p+V(q)E=0`$ is trivial under canonical quantisation. Second, the usual Hamiltonian $`H=p^2/2+V(q)`$ can be factorised into two factors of the first type. Third one of these factors can be converted to a simple momentum by a canonical transformation. For quantum mechanics we need Fourier and gauge transformations . Their composition introduces errors of order $`\mathrm{}^2`$ thus leading to an integral representation of an approximate solution. A judicious selection of the integration path depending on the coordinate yields converging integrals everywhere, thus guaranteeing a uniform approximation.
We shall show, that the saddle point approximation of the integral again yields the WKB solution, and indicate how the path of integration has to be laid to avoid singularities at turning points and extremal points of the potential. We will restrict our discussion to the one dimensional case and a Hamiltonian with local potential, where the action for a given energy $`E`$ is defined as
$$S(q)=^q𝑑s[2(EV(s))]^{1/2}$$
(1)
## 2 The factorisation
We start by introducing the function $`\mathrm{\Omega }(q,p,E)=H(q,p)E`$. Then the wave function is an eigenfunction with eigenvalue 0 of the operator version of $`\mathrm{\Omega }`$ under the usual canonical quantisation. If $`\mathrm{\Omega }`$ were linear in $`p`$, we could bring it into the form
$$\mathrm{\Omega }=pg(q,E)$$
(2)
where $`g`$ is the derivative of the action function (which would have only one branch in this case) and the exact wave function would have the simple form
$$\psi (q)=c\mathrm{exp}[iS(q)/\mathrm{}]$$
(3)
For Hamiltonian functions of the standard form $`\mathrm{\Omega }`$ is a polynomial of second degree in p and after multiplication by 2 we can decompose it into the linear factors
$$\mathrm{\Omega }=[pg(q,E)][p+g(q,E)]$$
(4)
where the two functions $`g`$ are the derivatives of the two branches of the action function which differ only in their sign. It is again trivial to construct eigenfunctions with eigenvalue 0 for each factor. But this does not solve the entire problem because of the ordering problem of quantum mechanics. In order to obtain an hermitian operator we must write the factors in some symmetrical order, and there are infinitely many ways to do so. Fortunately they only differ in $`\mathrm{}^2`$ and higher orders in $`\mathrm{}`$. Since we want to construct semi-classical solutions all these possibilities are equally correct for our purpose and we can select the one which is most convenient. We choose
$$\mathrm{\Omega }=[pg(q,E)]^{1/2}[p+g(q,E)][pg(q,E)]^{1/2}$$
(5)
Next we remember, that for canonical transformations which are represented by gauge transformations, this representation is quantum mechanically precise . At this stage it helps to apply the following canonical transformation
$$pp+g(q,E),qq$$
(6)
to bring $`\mathrm{\Omega }`$ into the form
$$\mathrm{\Omega }=p^{1/2}[p+2g(q,E)]p^{1/2}$$
(7)
Now assume that $`\stackrel{~}{\chi }(p)`$ is the Fourier transform of the wave function $`\chi (q)`$, which is eigenfunction with eigenvalue 0 of the operator version of $`[p+2g(q,E)]`$. Then $`\stackrel{~}{\varphi }(p)=p^{1/2}\stackrel{~}{\chi }(p)`$ is eigenfunction with eigenvalue 0 of the operator version of $`\mathrm{\Omega }`$ of Eq. 7 in momentum representation.
## 3 Construction of the wave function
To obtain an expression for the final wave function $`\psi (q)`$ we have to assemble all steps in reverse order. First it is obvious that the function $`\chi (q)`$ has the form
$$\chi (q)=c\mathrm{exp}[i(2S(q,E)/\mathrm{}]$$
(8)
such that its Fourier transform
$$\stackrel{~}{\chi }(p)=cdr\mathrm{exp}[i(rp+2S(r,E)/\mathrm{}]$$
(9)
yields $`\stackrel{~}{\chi }(p)`$ as used above, Here, and in what follows, we put all uninteresting factors into the normalisation constant $`c`$. According to Eq. 8 we obtain $`\stackrel{~}{\varphi }(p)`$ and next $`\varphi (q)`$ by an inverse Fourier transform as
$$\varphi (q)=c𝑑p𝑑rp^{1/2}\mathrm{exp}[i(qp+rp+2S(r,E))/\mathrm{}]$$
(10)
The p integral can be done in closed form giving
$$\varphi (q)=c𝑑rqr^{1/2}\mathrm{exp}[2iS(r,E)/\mathrm{}]$$
(11)
Now we make a substitution of the integration variable introducing $`s`$ as new integration variable according to $`r=qs^2`$ arriving at
$$\varphi (q)=c𝑑s\mathrm{exp}[2iS(qs^2,E)/\mathrm{}]$$
(12)
Finally to arrive at $`\psi (q)`$ we must undo the canonical transformation Eq. 6. This is done by multiplying the wave function by the gauge factor $`\mathrm{exp}[iS(q,E)/\mathrm{}]`$ and we obtain
$$\psi (q)=c\mathrm{exp}[iS(q,E)/\mathrm{}]𝑑s\mathrm{exp}[2iS_2(qs^2,E)/\mathrm{}]$$
(13)
This is the formal global solution we wished to obtain, once we determine the path of integration. Yet for formal manipulations the fact that we have this closed form may be quite important. For example we can readily see that we retrieve the WKB approximation by a power expansion, as long as we are far away from any turning point. Then we formally expand $`S(qs^2,E)`$ in a power series in $`s^2`$ and only keep the first two terms giving $`S(q,E)s^2S(q,E)/q`$. Plugging in into Eq. 21 gives
$`\psi (q)`$ $`=`$ $`\mathrm{exp}[iS(q,E)/\mathrm{}]{\displaystyle 𝑑s\mathrm{exp}[2is^2S(q,E)/q/\mathrm{}]}`$ (14)
$`=`$ $`c(S(q,E)/q)^{1/2}\mathrm{exp}[iS(q,E)/\mathrm{}]`$
which is the usual WKB solution. To arrive at this result we have taken the stationary phase contribution of the point $`s=0`$ to the $`s`$ integral. The exponent can have further stationary points at values $`s_c`$ of $`s`$ such that $`qs_c^2`$ is a turning point. However, in general, $`S`$ varies as $`(ss_c)^{3/2}`$ in the vicinity of such argument values; therefore its second derivative goes as $`(ss_c)^{1/2}`$. Accordingly, the contribution of such points to the saddle point evaluation of the integral has weight zero. Therefore the point $`s=0`$ is the only point giving contributions in saddle point evaluation of the integral.
## 4 The integration path
We have started from a second order differential equation. Therefore we must be able to obtain two linearly independent solutions. One way is to reverse the sign of $`S`$, the other is by appropriate choices of the integration path for the variable $`s`$ in the complex plane. The integrand is the exponential of some function $`f(s)`$, the square bracket in Eq. 20. Along some sectors for the angle $`\alpha `$ of the complex variable $`s`$ the function $`f(s)`$ acquires large negative real parts and the integrand decays exponentially. Let us call these intervals $`I_j`$. For the sectors in between the integrand explodes exponentially. An appropriate choice for the integration path is to come in from infinity in one angle sector $`I_{in}`$, to pass near the origin and to return to infinity in a different sector $`I_{out}`$. Some combinations of the two intervals will produce the same solution, and some the solution identically zero. But there should be two different choices leading to two different solutions. In general the function $`f(s)`$ can have isolated singularities, whose position depends on $`q`$ and $`E`$. Then we may eventually deform the integration path into one, encircling some of these singularities.
By a good choice of the integration path we can also be sure that the solution does not have singularities. To understand this, let us fix a value of $`E`$ and consider $`q`$ in a small neighbourhood of an arbitrary fixed point $`q_0`$. Assume for the moment that the potential $`V(q)`$ is an analytic function. Then the integrand may have some singularities in isolated points in $`s`$, but outside of them it is analytic. We call the singular points $`s_j`$. As we vary $`q`$ the singular points in $`s`$ will also move in general but remain in small neighbourhoods of $`s_j(q_0)`$. When we choose the integration path such that it avoids all these little neighbourhoods, then we obtain a function $`\psi (q)`$ which is analytic in the neighbourhood of the point $`q_0`$. When we vary $`q`$ over large ranges, then we have eventually to shift the integration path accordingly to avoid singularities.
The two most important situations where we need the explicit path of the integral are near the maximum of potentials and near turning points. It can be shown that in the first case, if we use the quadratic approximation for the extremum we retrieve the exact solution, i.e. the Pearcy function. In the second case it does not seem easy to find a path that yields the exact solution, but by choosing a path that fulfills the above conditions, we obtain a solution that has no singularities and numerical inspection shows it to be quite close to the Airy function.
## 5 Conclusions
We have obtained an integral representation for a semi-classical approximation, of the wave-function of a standard Hamiltonian with local potential. The method involves a factorisation, which causes errors in higher orders of $`\mathrm{}`$, which is acceptable for a semi-classical approximation. As we may expect the usual WKB method results from a saddle-point approximation of this integral. Judicious choices of the path of the integral in the complex plane lead to approximate solutions which have no singularities and thus converge everywhere adequately. The method may be readily generalized to Hamiltonians which are of higher order in the momenta. In particular examples it might be useful to modify the method slightly by multiplying and dividing by additional factors. This can in many cases provide the exact solution for all values of the energy. In a future publication we will present such examples in detail.
The advantages of this method consist in the fact that the points where WKB breaks down are not treated piecemeal, but are covered by the same integral representation. Whenever we wish to make an analytic statement about semi-classics this can be a great advantage.
## Acknowledgments
This work is supported by the UNAM DGAPA project IN-102597. |
warning/0001/quant-ph0001076.html | ar5iv | text | # COVARIANCE, CORRELATION AND ENTANGLEMENT
## 1 Introduction
Several measures of entanglement or quantum correlations have been proposed: some are associated with the preparation of the state, others with the process of purification or distillation and yet others with the notion of mutual information or relative entropy . In this paper we wish to suggest another measure, based on covariance, in which the acts of state creation and observation are considered in a dual manner.
In practice it is very natural to describe the condition of the system (its method of preparation or lack of it) in terms of a density matrix which is tied to the subsequent observations on it. This is how the linkage between observer and observed occurs quantum mechanically, and of course the results are expressed in terms of traces over appropriate functions of the density matrix and of the operators being measured . Indeed Mermin has taken that view that the density matrix, and the correlations between observables which thereby ensue, constitute all of physical reality.
In this paper we will also focus on the density matrix. Because binning of observations is a necessity in practice, the dimension of the density matrix is thereby determined: $`N`$ separate bins produce a density matrix $`\rho `$ that is an $`N\times N`$ hermitian matrix, satisfying the usual hermiticity and trace conditions. In this way, we can regard the basis as an $`N`$-level system, rather like a particle of angular momentum $`J=(N1)/2`$. Hence, although we might be studying the probability distribution of an observable which actually possesses a continuous spectrum, we can still regard it as a spin-like system; in practical terms, the bigger the binning number $`N`$, the greater the precision of the information about the continuous variable, but obviously $`N`$ is never infinite. For spin measurements, we need not go to such pains because $`N`$ is fixed for us at the start.
With the focus on density matrices, we will carry out measurements (without mutual interference) on two subsystems, 1 and 2 say, so their corresponding observables, superscripted by (1) and (2), are commuting operators. The system will be separable or “disentangled” if the larger density matrix $`\rho `$ is merely a direct product of density matrices associated with the two subsystems, or $`\rho =\rho ^{(1)}\rho ^{(2)}`$; a particular case arises when the initial or prepared state is the direct product of two subsystem states, $`|\psi =|\varphi ^{(1)}|\chi ^{(2)}`$. When two subsystems are disentangled, the results of measuring any quantity $`A^{(1)}`$ in the first subsystem are not tied to the results of measuring any quantity $`B^{(2)}`$ in the second subsystem; necessarily
$$A^{(1)}B^{(2)}=A^{(1)}B^{(2)},$$
for all choices of $`A`$ and $`B`$. However, if $`\rho \rho ^{(1)}\rho ^{(2)}`$, the configuration is non-factorizable and the covariance,
$$\mathrm{cov}(A^{(1)}B^{(2)})A^{(1)}B^{2)}A^{(1)}B^{(2)},$$
(1)
no longer disappears.
The real issue is how to quantify the entanglement or lack of factorizability of the larger density matrix. Several proposals have been advanced in the literature, but none of them is entirely simple or definitive . However all researchers in this field seem to agree on the following three conditions for an entanglement measure $`E(\rho )`$:
1. $`E(\rho )=0`$ iff $`\rho `$ is separable, ie if the density matrix can be written as $`\rho =_ip_i\rho _i^{(1)}\rho _i^{(2)}`$.
2. Local unitary transformations should leave $`E(\rho )`$ invariant.
3. $`E(\rho )`$ should not increase under local measurement and classical communication procedures, we intuitively know that such procedures cannot add non-locality characteristics to the system being measured.
As an extra requirement, it would be nice if $`E(\rho )`$ gave some indication of the extent of violation of Bell-type inequalities .
In this paper we want to put forward a concrete scheme for quantifying correlations between two subsystems and their possible entanglement. The scheme is based on a generalization of eq. (1) and particular choices of operators $`A^{(1)}`$ and $`B^{(2)}`$, which are readily applicable and rooted in the density matrix notion. In the next section we discuss several matters connected with non-separability of states and their influence on subsequent subsystem measurements. Because we deal with practical observations, the density matrix is truly discrete and we can assume that the elements of the vector space on which it lives have equal weight. As already mentioned, one may regard the dimension $`N=N^{(1)}+N^{(2)}`$ as corresponding to a “spin system”, with each component carrying equal weight, and can adopt the same stance for the subsystem dimensions $`N^{(1,2)}`$. (This restriction can be relaxed if the components have unequal weights, such as atomic energy levels at a finite temperature.)
The next section deals with the generalities of simultaneous measurements and their covariance properties . This is followed by our suggestion for quantifying entanglement of two subsystems within a larger entity, which is shown to be consistent with normal expectations for two spin 1/2 subsystems, when $`N^{(1)}=N^{(2)}=2`$. We also discuss the use of total spin dispersion as another measure of entanglement, with an allied appendix concerning the Majorana-Penrose representation of spin states on the Poincaré sphere. The subsequent sections deal with entanglement measures for larger value of $`N^{(1)}`$ and $`N^{(2)}`$. Finally we discuss general questions pertaining to our suggested measure; these include Rovelli’s notion that information in quantum mechanics is relational , Mermin’s notions of correlations between local observables , and the difference between our modified correlation measure with classical correlations for impure states.
## 2 Correlations and density matrices
Elementary texts on quantum mechanics teach us that the results of all physical measurements and processes can be tied to the evaluation of traces of products of observables with the hermitian density matrix $`\rho `$. Thus statistical formulae like
$$F=\mathrm{Tr}(\rho F);\mathrm{Tr}(\rho ^2)1,$$
etc. are part of the standard repertoire. Of course, the $`\rho `$-eigenvalues lie between 0 and 1; in the latter case we are dealing with a pure state when the density matrix reduces to a projector $`\rho P_\psi |\psi \psi |`$, while the most random situation $`\rho =1/N`$ corresponds to the case of maximum entropy.
The covariance for any two commuting observables $`A,B`$ in a mixed state $`\rho `$ is defined as
$$\mathrm{cov}_\rho (A,B)ABAB=\mathrm{tr}(\rho AB)\mathrm{tr}(\rho A)\mathrm{tr}(\rho B).$$
(2)
Clearly, $`\mathrm{var}_\rho (A)=\mathrm{cov}_\rho (A,A)`$. Less well-known is the fact that pure state dispersions and correlations can be neatly expressed in terms of a single trace. Consider the quantity
$$C_\rho (A,B)\mathrm{tr}([\rho ,A][B,\rho ])/2=\mathrm{tr}(\rho ^2\{A,B\}/2\rho A\rho B),$$
(3)
where $`A`$ and $`B`$ are any two operators. This quantity will be referred to as the alternative covariance <sup>1</sup><sup>1</sup>1Evaluating traces of larger numbers of pure state commutators, one may establish algebraically that for odd numbers of products, the traces do vanish. For instance, $`\mathrm{tr}\left([A,\rho ][B,\rho ][C,\rho ]\right)=0`$, etc..
We now present some elementary results about $`C_\rho `$ which follow simply from this definition:
1. $`C_\rho (A,B)=C_\rho (Aa,Bb)`$, where $`a,b`$ are any two constants.
2. $`C_\rho (aA,bB)=abC_\rho (A,B)`$.
3. $`C_\rho (_iA_i,_jB_j)=_{i,j}C_\rho (A_i.B_j).`$
4. $`C_\rho (A,A)=\mathrm{tr}(\rho ^2A^2\rho A\rho A).`$
5. $`C_{U\rho U^{}}(A,B)=C_\rho (U^{}AU,U^{}BU)`$, where $`U`$ is any unitary transformation. Thus a change of basis for the state is equivalent to an inverse change of basis for the operators.
6. $`C_\rho (A,A^{})C_\rho (B,B^{})|C_\rho (A,B^{})|^2`$. This follows by considering the operator $`T=[\rho ,AcB]`$, with $`c=\mathrm{tr}([\rho ,A][A^{},\rho ])/\mathrm{tr}([\rho ,A][B^{},\rho ],`$ and noting that $`\mathrm{tr}(TT^{})0.`$
7. $`C_\rho (A,A^{})`$ is real.
8. $`|C_\rho (A,B)|^2`$ is symmetrical under interchange, conjugation and change of phase of the two operators.
All of these properties are shared by the usual covariance cov$`{}_{\rho }{}^{}(A,B)`$. Nevertheless, alternative covariance $`C_\rho `$ does not provide an indication of variance and covariance in the usual sense. For instance, if the state is one of maximum entropy, on the one hand we have $`C_\rho (A,B)=0`$ for all $`A,B`$ because $`\rho `$ is proportional to unity; on the other hand, the $`\mathrm{cov}_\rho (A,B)`$ need not vanish.
Some special cases for the operators $`A,B`$ can now be studied.
1. If $`A`$ and $`B`$ commute, $`C_\rho (A,B)=\mathrm{tr}(\rho ^2AB\rho A\rho B).`$
2. If $`A`$ and $`B`$ are both hermitian, $`C_\rho (A,B)`$ becomes real.
3. If $`A`$ and $`B`$ are both unitary, $`C_\rho (A,A^{})=\mathrm{tr}(\rho ^2\rho A\rho A^{})\mathrm{tr}(\rho ^2)1.`$ Likewise for $`B`$. Since $`C_\rho (A,A^{})C_\rho (B,B^{})|C_\rho (A,B^{})|^2`$, it follows that $`|C_\rho (A,B)|^21.`$
More particular cases arise when the system is prepared in a pure state $`|\psi `$, so that $`\rho `$ becomes a projection operator and $`C_\rho (A,B)`$ reduces to
$`C_\rho (A,B)`$ $``$ $`{\displaystyle \frac{1}{2}}\psi |\{A,B\}|\psi \psi |A|\psi \psi |B|\psi ={\displaystyle \frac{1}{2}}\{AA,BB\}`$ (4)
$`=`$ $`\mathrm{cov}_\rho (A,B),\mathrm{when}[A,B]=0.`$
Thus
$$C_\rho (A,A)A^2A^2=\mathrm{var}_\rho (A).$$
(5)
This is in keeping with the familiar variance-covariance inequality:
$$\mathrm{var}_\rho (A)\mathrm{var}_\rho (B)|\mathrm{cov}_\rho (A,B)|^2.$$
(6)
If $`A`$ and $`B`$ are commuting unitary operators and because $`|C_\rho (A,B)|^2C_\rho (A,A)C_\rho (B,B)(\mathrm{tr}(\rho ^2))^21,`$ we see that alternative covariance only attains a value of 1 for pure states.
It is worthwhile comparing the two covariance functions, in relation to two commuting observables, $`A,B`$. Since $`[A,B]=0`$, select an orthonormal basis $`|i`$ wherein the operators are simultaneously diagonalised, so
$$A=\underset{i}{}|ia_ii|,B=\underset{i}{}|ib_ii|.$$
Then
$`\mathrm{cov}_\rho (A,B)`$ $`=`$ $`\mathrm{tr}(\rho AB)\mathrm{tr}(\rho A)\mathrm{tr}(\rho B)`$
$`=`$ $`{\displaystyle \underset{i}{}}a_ib_ii|\rho |i{\displaystyle \underset{i,j}{}}a_ib_ji|\rho |ij|\rho |j`$
$`=`$ $`{\displaystyle \underset{i,j}{}}a_i(b_ib_j)i|\rho |ij|\rho |j,`$
and similarly
$`C_\rho (A,B)`$ $`=`$ $`\mathrm{tr}(\rho ^2AB)\mathrm{tr}(\rho A\rho B)`$
$`=`$ $`{\displaystyle \underset{i,j}{}}a_i(b_ib_j)i|\rho |jj|\rho |i,`$
since $`_jj|\rho |j=\mathrm{tr}(\rho )=1`$. Furthermore note that $`_{i,j}i|\rho |jj|\rho |i=\mathrm{tr}(\rho ^2)1.`$ Upon symmetrising the sums, we obtain the neater expressions,
$$\mathrm{cov}_\rho (A,B)=\underset{i,j}{}(a_ia_j)(b_ib_j)i|\rho |ij|\rho |j/2$$
(7)
$$C_\rho (A,B)=\underset{i,j}{}(a_ib_j)(b_ib_j)i|\rho |jj|\rho |i/2$$
(8)
Whilst the ordinary covariance has a clear meaning—namely, a measure of the correlations between the results of local measurements $`A`$ and $`B`$ that commute—the interpretation of the alternative covariance is less obvious.
We can obtain more insight by choosing $`A=B`$. Since $`\rho `$ is a positive definite hermitian operator, $`i|\rho |ij|\rho |ji|\rho |jj|\rho |i`$, for any two states $`|i,|j`$. Therefore for a general (mixed state) density matrix,
$$\underset{i,j}{}(a_ia_j)^2i|\rho |ij|\rho |j\underset{i,j}{}(a_ia_j)^2i|\rho |jj|\rho |i,$$
or
$$\mathrm{var}_\rho (A)C_\rho (A,A).$$
(9)
In the light of the variance inequality above it is natural to ask whether
$$|\mathrm{cov}_\rho (A,B)||C_\rho (A,B)|$$
is true. In fact a single (but carefully chosen) counterexample suffices to show that it is false: in local bases $`u,d`$ for two local operators $`A1`$ and $`1B`$, take
$$\rho =|(uu+dd)(uu+dd)|/4+|udud|/4+|dudu|/4$$
and select the local operators to be diagonal,
$`A1`$ $`=`$ $`(|uuuu|+|udud|)(|dudu|+|dddd|)`$
$`1B`$ $`=`$ $`(|uuuu||udud|)(|dudu||dddd|).`$
Evaluation of the two types of covariance leads to
$$\mathrm{cov}_\rho (A,B)=0,\mathrm{but}C_\rho (A,B)=1/4.$$
Thus the variance inequality cannot be extended to covariance.
However, an immediate consequence of the inequality, var$`{}_{\rho }{}^{}(A)C_\rho (A,A)`$, is that when var$`{}_{\rho }{}^{}(A)=0`$, $`C_\rho (A,A)=0`$ too for any observable $`A`$. But $`2C_\rho (A,A)=\mathrm{tr}([A,\rho ][A,\rho ]^{})`$; so $`[A,\rho ]=0`$, which means that $`\rho `$ is purely in an eigenstate of $`A`$. This accords with the basic tenets of quantum mechanics of course. The contrapositive of this result is that if $`[A,\rho ]0`$, then $`\mathrm{var}_\rho (A)>0`$.
Another worthwhile comment stems from the observation that if $`X`$ is conjugate to $`A`$ in the sense $`[A,X]=i\mathrm{}`$, then
$$2C_\rho (A,A)=\mathrm{tr}\left([\rho ,A][\rho ,A]^{}\right)=\mathrm{}^2\mathrm{tr}\left(\frac{\rho }{X}\frac{\rho }{X}^{}\right).$$
Thus,
$$\mathrm{var}_\rho (A)=(\mathrm{\Delta }A)^2\mathrm{}^2\mathrm{tr}\left(\left|\frac{\rho }{X}\right|^2\right)/2,$$
with equality only applying to pure states. For example, the energy uncertainty is given by $`(\mathrm{\Delta }H)^2\mathrm{}^2\mathrm{tr}\left(\left|\frac{d\rho }{dt}\right|^2\right)/2`$, while the momentum uncertainty is given by the derivative of the density matrix with respect to position: $`(\mathrm{\Delta }P)^2\mathrm{}^2\mathrm{tr}\left(\left|\frac{\rho }{X}\right|^2\right)/2`$, and so on.
For a general mixed configuration, the two inequalities,
$$\mathrm{var}_\rho (A)\mathrm{var}_\rho (B)C_\rho (A,A)C_\rho (B,B)\left|C_\rho (A,B)\right|^2$$
together with the well-known
$$\mathrm{var}_\rho (A)\mathrm{var}_\rho (B)|\mathrm{cov}_\rho (A,B)|^2,$$
provide a lower bound for the experimentally observed variance products of any two operators, whether or not they commute. For instance,
$$\mathrm{var}_\rho (X).\mathrm{var}_\rho (P)\left|\frac{1}{2}\mathrm{}^2\mathrm{tr}\left(\frac{\rho }{X}\frac{\rho }{P}\right)\right|^2.$$
In the next section we present examples of operators $`A,B`$, for which there exist states such that $`|\mathrm{cov}_\rho (A,B)||C_\rho (A,B)|`$ and also other states for which $`|C_\rho (A,B)||\mathrm{cov}_\rho (A,B)|`$. Hence both inequalities must be considered jointly in an examination of the minimum of the variance product, together with Heisenberg’s well-known lower bound, $`|\mathrm{tr}(\rho [A,B])|/2`$.
## 3 Correlation measures for pure states of two subsystems
This section examines the correlation properties of the entanglement of two subsystems in a tensor product Hilbert space $`^{(1)}^{(2)}`$. By definition, measurements can be carried out without mutual interference on the two subsystems so their corresponding observables, superscripted by (1) and (2), are commuting operators. As mentioned in the introduction, a state of the system is factorisable or disentangled if the larger density matrix $`\rho `$ is merely a direct product of density matrices associated with the two subsystems, or $`\rho =\rho ^{(1)}\rho ^{(2)}`$. For any two local operators $`A^{(1)}=A1,B^{(2)}=1B`$, it is easy to show that the covariance disappears:
$$\mathrm{cov}_{\rho ^{(1)}\rho ^{(2)}}(A^{(1)},B^{(2)})=(AA)(BB)=0.$$
(10)
This includes the case of a pure disentangled state, $`|\varphi =|\varphi ^{(1)}|\varphi ^{(2)}`$.
Having noted that the covariance is non-zero in disentangled states, we now refer to the conditions imposed upon any measure of entanglement. The second condition is that it be invariant under local unitary transformations. With this in mind, define the covariance entanglement for pure states as
$$E_{A^{(1)},B^{(2)}}(\rho )\underset{U=U^{(1)}U^{(2)}}{\mathrm{max}}|\mathrm{cov}_{U\rho U^{}}(A^{(1)},B^{(2)})|.$$
The maximum will clearly be invariant under additional local unitary transformations.
Since the operation of permuting the elements of a Hilbert space is unitary, all elements of the Hilbert spaces are equally important. For this reason, it is natural to select the operators $`A^{(1)}`$ and $`B^{(2)}`$ so as to equally weight the elements of the Hilbert space. The next section describes several ideas for achieving this, starting with the simplest case.
### 3.1 Pure state correlations for $`N^{(1)}=N^{(2)}=2`$
This section investigates a method for quantifying pure state entanglement in the simplest possible non-trivial case, corresponding to two spin 1/2 systems, with Hilbert space $``$, where $`=^2`$ is a local Hilbert space, with orthonormal basis $`|u,|d`$. Consider two local operators which distinguish between elements of the local Hilbert spaces. With the aim of weighting local basis elements equally, define operators in the product basis $`|uu,|ud,|du,|dd`$ by
$$A^{(1)}=\sigma _3^{(1)}=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right),B^{(2)}=\sigma _3^{(2)}=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right),$$
and so
$$A^{(1)}B^{(2)}=\sigma _3^{(1)}\sigma _3^{(2)}=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right).$$
Next consider the pure (normalized but arbitrary) state,
$$|\varphi =\alpha |uu+\beta |ud+\gamma |du+\delta |dd.$$
Working out cov$`{}_{\rho }{}^{}(A,B)`$ in this state, it is straightforward to show that the covariance is maximised provided that $`|\alpha |=|\delta |=1/\sqrt{2},\beta =\gamma =0`$, or $`|\beta |=|\gamma |=1/\sqrt{2},\alpha =\delta =0`$. Thus one may take the four independent Bell states,
$$|10[|ud+|du]/\sqrt{2},|00[|ud|du]/\sqrt{2},$$
$$|1+[|uu+|dd]/\sqrt{2},|1[|uu|dd]/\sqrt{2},$$
as the ones that have the largest covariance. (These pure states are also known to be the most entangled ones.) Of course they are all local unitary transforms of just one of them, say the Bell state, $`|1+|(uu+dd)/\sqrt{2}`$, with a corresponding $`\rho =|uu+dduu+dd|/2`$. If one rotates the operators $`A,B`$ together about the “$`y`$-axis” by the same amount, we can get a good idea of how the covariance varies with rotation angle; maximization is attained when the angle is $`n\pi `$. See Figure 1.
If a pure state is disentangled, then there is an orientation of the local $`A^{(1)}`$ which has zero variance. To see this, rotate the state by local unitary transformations until the reduced density matrix for the local operator in question is diagonalised. Since the initial state was pure and disentangled, then it may be represented by a separable projector $`\rho =\rho ^{(1)}\rho ^{(2)}`$, in which both reduced density matrices are projectors. Thus the main diagonals of $`\rho ^{(1)}`$ and $`\rho ^{(2)}`$ can be reduced to a single 1, with 0s elsewhere. Hence for diagonalised local operators $`A`$ and $`B`$ in this basis,
$$A^2=A^2,B^2=B^2,$$
so the variances vanishes. Figure 2 illustrates the behaviour of the variance in a 2-variable parametrisation. The horizontal axis variable $`x`$ parametrises a set of pure states which range from disentangled to a maximally entangled Bell state, and back to disentangled again, i.e. $`\rho =|\psi (x)\psi (x)|`$ where $`|\psi (x)=\mathrm{cos}(x)|uu+\mathrm{sin}(x)|dd`$. The second variable $`y`$ parametrises $`y`$-axis rotations of the local spin basis (1) associated with $`A`$ alone.
These results may be applied to actual spin measurements. If one knows that a state is pure, but is not certain of the degree of entanglement, local spin measurements can be made in a variety of directions. If the variance and covariance of these measurements vanishes in a pure state, then the state must be disentangled.
### 3.2 Pure state correlations for $`N^{(1)}=N^{(2)}N>2`$
Many different choices of local operators $`A^{(1)},B^{(2)}`$ are possible, and different choices will lead to different behaviour of the covariance. However before considering two sensible choices, let us note that for $`N^{(1)}=N^{(2)}=2`$, one of the maximally entangled states can be taken to be the state of total spin $`|10`$, while minimally entangled states are $`|11,|11`$. Now for these state combinations maximal entanglement happens to equate with maximal dispersion $`(\mathrm{\Delta }J)^2`$ and zero entanglement equates with minimal dispersion $`(\mathrm{\Delta }J)^2`$, where $`J`$ stands for total spin. This suggests that for higher spin, some maximally entangled states might be found by minimising the total angular momentum dispersion and vice versa. This approach towards quantifying entanglement is quite interesting in its own right and is pursued in Appendix A, where we also tie it to the Majorana-Penrose pictorial view of spin. In Appendix B, by contrast, we classify any measure of entanglement via an integrity basis for density matrix invariants.
#### 3.2.1 Pair-discrimination
Consider two-system Hilbert spaces where the local spaces may each have dimension greater than 2. Select local operators $`A^{(1)}=A1,B^{(2)}=1B`$, where $`A,B`$ may be expressed in their respective local bases as unitary transformations of the following diagonal matrix,
$$\left(\begin{array}{cccc}1& 0& 0& \mathrm{}\\ 0& 1& 0& \mathrm{}\\ 0& 0& 0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right).$$
(11)
Next, maximise the covariance over all such unitarily-transformed matrices; the result is perforce invariant under local unitary transformations of the state. Labelling the local bases $`|a_i`$ and $`|b_i`$ respectively, where $`i`$ runs from 1 to $`N`$, operators like the above discriminate between pairs of elements in a local subspace of the full Hilbert space, and treat terms of the form $`|a_ib_k+a_jb_l`$ as the basic element of entanglement<sup>2</sup><sup>2</sup>2Another possibility is to replace all the zeroes along the main diagonal of $`A`$ or $`B`$ with $`1`$ or $`1`$; this makes the operator $`U`$ unitary, which means that its variance is simply $`1|U|^2`$. However the distribution of the $`\pm 1`$ eigenvalues is not self-evident, except for the spin $`\frac{1}{2}`$ spin$`\frac{1}{2}`$ case..
#### 3.2.2 Equal weight unitary operators
Since we wish to handle all the subsystem states democratically, let us define an equal weight local unitary operator as consisting of some unitary transformation of a diagonal matrix comprising the $`N`$-th roots of unity. (Note that these matrices are not hermitian when $`N3`$ and cannot correspond to observables.) Here the local weight unitary matrices in their corresponding diagonalising basis, up to an overall phase, are given by
$$\left(\begin{array}{cccc}\mathrm{exp}(2i\pi /N)& & & \\ & \mathrm{exp}(4i\pi /N)& & \\ & & \mathrm{}& \\ & & & 1\end{array}\right).$$
Only for the spin $`\frac{1}{2}`$ $`\times `$ spin $`\frac{1}{2}`$ case, do these operators correspond exactly to the pairwise local unitary operators used previously. The next nontrivial case is $`N=3`$, or spin 1 $`\times `$ spin 1. In this case one can see that all states which are local unitary transformations of the pure state $`|a_1b_1+a_2b_2+a_3b_3/\sqrt{3}`$ have a maximal covariance of 1. To understand why, rotate the equal-weight unitary operators so that the local states $`a_1,a_2,a_3`$ and $`b_1,b_2,b_3`$ produce eigenvalues which are respectively conjugate pairs. This yields $`A^{(1)}B^{(2)}=1`$. However, $`A^{(1)}=B^{(2)}=0`$, so the maximised covariance is 1.
What of other states, such as $`|a_1b_1+a_2b_2`$ when $`N=3`$? The following theorem gives a necessary and sufficient condition for a state to exhibit $`C_\rho =1`$, with respect to these equal weight operators and is in agreement with all other pure state entanglement measures.
Theorem: The only pure states which attain the maximised covariance of 1 under equal-weight local operators are states which are local unitary transformations of $`|a_1b_1+a_2b_2+\mathrm{}+a_Nb_N/\sqrt{N}`$. Any other states exhibit smaller correlations.
Proof: In a tensor product space $`^{(1)}^{(2)}`$, consider the state
$$|\psi =\underset{i,j}{}c_{ij}|a_ib_j,$$
where $`p_{ij}|c_{ij}|^20`$ and $`_{i,j}p_{ij}=1`$, and the orthonormal bases $`|a_i,|b_j`$ are a complete set of eigenstates for the diagonalised equal-weight unitary operators $`A^{(1)},B^{(2)}`$, with eigenvalues $`a_k=\mathrm{exp}(2i\pi k/N),b_j=\mathrm{exp}(2i\pi j/N)`$ respectively. Recalling the result for such $`A^{(1)},B^{(2)}`$ that
$$(1|A^{(1)}|^2)(1|B^{(2)}|^2)|A^{(1)}B^{(2)}A^{(1)}B^{(2)}|^2,$$
we see that a maximised covariance of 1 is only attainable in a state where $`A=B=0`$, whereupon the covariance reduces to
$$C_\psi (A^{(1)},B^{(2)})=\left|\underset{i,j}{}p_{ij}a_ib_j\right|.$$
In what cases is this expression maximised, subject to the condition that the mean values of the operators remain zero?
We are seeking to maximise $`|p_{ij}a_ib_j|`$ subject to
$$p_{ij}=1,|p_{ij}a_i|=0,|p_{ij}b_j|=0.$$
By the triangle inequality,
$$|p_{ij}a_ib_j|p_{11}|a_1b_1|+|\underset{i,j>1}{}p_{ij}a_ib_j|\mathrm{}\underset{i,j}{}p_{ij}|a_ib_j|,$$
where equality holds at every stage only if all the complex numbers $`a_i,b_i`$ have equal and opposite phase. This means that we are pairing $`a_i,b_j`$ such that non-zero $`p_{ij}`$ (only for $`i=j`$) are associated in a one-to-one manner with $`a_ib_i=1`$ for all such pairs, otherwise the parallelism of the complex numbers will be lost. This will ensure that
$$\underset{i}{}p_{ii}|a_ib_i|=\underset{i}{}p_{ii}=1.$$
At the same time we have to guarantee that the average values of $`A`$ and $`B`$ vanish or $`|p_{ii}a_i|=0,|p_{ii}b_j|=0`$. Since $`_ia_i=_ib_i=0`$, a sufficient condition for this is that for all such pairs, the weightings are equal or $`p_{ii}=1/\sqrt{N}`$; in other words every $`a_i`$ and its corresponding $`b_i=a_i^{}`$ only occurs at most once in the terms with equal non-zero weighting. (Actually one may introduce an arbitrary phase into $`c_{ii}`$ without affecting this conclusion, but we have chosen not to do so.)
Having established that the states which maximise the covariance can take the form $`|a_ia_1^{}+a_2a_2^{}+\mathrm{}`$, we should point out that it is not necessary for all the terms to be paired up. Consider the case $`N=4`$, or spin $`3/2`$ spin $`3/2`$ systems, with bases $`a_1,\mathrm{},a_4`$ and $`b_1,\mathrm{},b_4`$ respectively. It is possible to attain maximised covariances of 1 under the equal-weight measures both for $`|\psi =|a_1b_1+a_2b_2/\sqrt{2}`$ and for $`|\varphi =|a_1b_1+a_2b_2+a_3b_3+a_4b_4/2`$. This is achieved by choosing eigenvalues so that $`A^{(1)}=B^{(2)}=0`$ and yet pick eigenvalues such that the values of $`A^{(1)}B^{(2)}`$ in the states $`|a_1b_1\mathrm{}|a_4b_4`$ are all 1. As we are in a 4-dimensional space, we can achieve this by taking the eigenvalue sets $`\{1,1\}`$ and $`\{1,i,1,i\}`$ for both operators on both states $`|\psi `$ and $`|\varphi `$ respectively.
This observation means that the ‘equal-weight’ $`C`$-based measures for $`N>2`$ are not measures of entanglement, under the standard criteria. Information-based entanglement measures, such as the relative entropy, specify that $`|\psi `$ is less entangled than $`|\varphi `$. It appears that the covariance-based measures of entanglement are most useful when dealing with spin $`\frac{1}{2}`$ spin$`\frac{1}{2}`$ systems, since in this case there is no ambiguity as to the choice of eigenvalues. Similarly for operators of prime dimension, such ambiguities are absent, because there is only one way to arrive at a maximally correlated state: the non-uniqueness only pertains to composite-dimensional local spaces.
The above result demonstrates that for pure states, the reduced density matrices must be diagonalised in order to maximise the covariance of the diagonalised local unitary operators. However, for mixed states it is not at all obvious that the reduced density matrices must be diagonalised in order to maximise the unitary matrix covariance. This issue will will be examined in section 4.
### 3.3 Pure state correlations for $`N^{(1)}N^{(2)}`$
For spaces of differing dimension, such as spin $`1\times `$ spin $`1/2`$, the covariance of the pairwise operators behave just as in the other cases. However, if the equal-weight local operators are used, covariances of 1 are not attainable. This reflects the fact that the bases have different sizes, and so there is no way to pair up elements between the bases in a one-to-one manner so as to produce a set of product eigenvalues with the same phase.
The simplest example which exhibits this effect is a spin $`\frac{1}{2}`$ spin $`1`$ space, with equal-weight operators
$$A^{(1)}=\mathrm{diag}(1,1)1,B^{(2)}=1\mathrm{diag}(1,\mathrm{exp}(2i\pi /3),\mathrm{exp}(4i\pi /3))$$
(12)
As with the proof that the states of maximum covariance are $`|a_1b_1+\mathrm{}+a_nb_n`$, a covariance of 1 is only attainable if $`A^{(1)}B^{(2)}=1`$; we also need the complex numbers $`p_{ij}a_ib_j`$ to have the same phase (where the state has Schmidt decomposition $`|\varphi =_{ij}\sqrt{p_{ij}}|a_ib_j`$ in the basis of eigenvalues of the operators $`A^{(1)},B^{(2)}`$). Following similar arguments to those used in the proof, it is not possible to choose $`p_{ij}`$ so that the direction condition is satisfied; this is because no repetitions of $`a_i`$ or $`b_j`$ values can occur in the set of non-zero $`p_{ij}`$, since the resulting $`p_{ij}a_ib_j`$ would not have the same phase. But it is not possible to partially pair up the given set of eigenvalues so that the directions of the products are the same, by straightforward enumeration of the cases. Thus states in this basis cannot attain covariances of $`1`$.
Clearly many other local operators can be defined which provide a variety of different correlation measures for two-system states but none stands out.
## 4 Correlation measures for mixed states
Making a distinction between quantum and classical correlations has proved a thorny problem in the study of quantum entanglement. The nature of the problem may be seen when comparing the two states, one pure and one mixed, which possess the same covariance for $`A^{(1)}=B^{(2)}=\sigma _3`$:
$$\frac{1}{2}|uu+dduu+dd|\mathrm{and}\frac{1}{2}|uuuu|+\frac{1}{2}|dddd|.$$
The first state is a Bell state and is maximally entangled, whilst the second state is a mixture of disentangled projectors, and is normally regarded as being disentangled. As both states exhibit correlations, it is natural to ask whether the alternative covariance introduced in section 2 provides a way of distinguishing between classical and quantum correlations.
### 4.1 Distinction between $`C_\rho `$ and cov<sub>ρ</sub>
Take any two commuting local measurements, $`A^{(1)}=A1,B^{(2)}=1B`$, and define the function
$$E_{AB}(\rho )\mathrm{max}_U|C_{U\rho U^{}}(A^{(1)},B^{(2)})|^2,$$
(13)
where the maximum is now taken over all local unitary transformations of the general mixed density matrix $`\rho `$. In previous sections we investigated the behaviour of this function for pure states (when $`C_\rho `$ reduces to the covariance), and found that it appeared to have many of the properties desirable in a pure state entanglement measure.
The situation where the density matrix corresponds to an impure configuration is more intriguing. Shown below is a comparison of the behaviour of the two maximised covariance entanglement measures, in several example configurations, which illustrate the distinction between $`|\mathrm{cov}_\rho (A^{(1)},B^{(2)})|^2`$ and $`|C_\rho (A^{(1)},B^{(2)})|^2`$.
| Density matrix | $`|\mathrm{cov}_\rho (A^{(1)},B^{(2)})|^2`$ | $`|C_\rho (A^{(1)},B^{(2)})|^2`$ |
| --- | --- | --- |
| $`\rho _1=[|uuuu|+|dddd|]/2`$ | 1 | 0 |
| $`\rho _2=[\frac{1}{2}|uuuu|+\frac{1}{4}|uu+dduu+dd|]`$ | 3/4 | 1/4 |
| $`\rho _3=|uuuu|`$ | 0 | 0 |
| $`\rho _4=|uu+dduu+dd|/2`$ | 1 | 1 |
The illustrative state $`\rho _4`$ is pure but entangled, $`\rho _3`$ is factorizable and therefore disentangled, while the matrix $`\rho _1`$ is not factorizable but can be expressed as a sum of separable projectors; therefore $`\rho _1`$ should represent a disentangled configuration, according to standard expectations. By inspecting the table we see that cov$`{}_{\rho }{}^{}(A^{(1)},B^{(2)})`$, being nonzero, is not a good entanglement measure $`E(\rho )`$, but the alternative $`C_\rho (A^{(1)},B^{(2)})`$ is better in that it does vanish.
The alternative covariance $`C`$ in mixed states is bounded above by $`\mathrm{tr}(\rho ^2)`$, so the state must be pure to obtain a covariance of 1 under local unitary transformations. Another point worth remembering is that for configurations of maximum entropy where $`\rho 1`$ and commutes with all operators, $`C`$ is automatically zero.
Figures 3 and 4 depict the squares of the equal weight classical covariance and alternative covariance of a parametrised mixture of different Bell states. They show that both covariance measures attain a maximum of 1 only for the pure states, and that the behaviour of these functions depends upon the kind of Bell-states involved in the mixtures.
For any choice of $`A^{(1)},B^{(2)}`$, a maximised $`C_\rho (A^{(1)},B^{(2)})`$ of zero occurs for many mixed configuration $`\rho `$ which are disentangled, according to the standard definition of a mixed state, such as the states of maximum entropy with $`\rho 1`$. This behaviour of the alternative covariance for ‘partially entangled’ configurations is what one would naively expect for a mixed configuration entanglement measure $`E(\rho )`$. However, there is a particular class of mixed states which are disentangled according to the commonly accepted definition of a mixed state, whilst exhibiting non-zero alternative covariance properties.
### 4.2 Conditions on entanglement measures
This section assesses the alternative covariance as a measure of entanglement, according to the principles outlined by Vedral and Plenio.
Firstly, consider the cases where a state exhibits zero quantum correlation. A mixed configuration is usually defined to be separable if it can be written as a sum of separable projectors; thus
$$\rho =\underset{j}{}p_j\rho _j^{(1)}\rho _j^{(2)};\rho _i^{(n)}\rho _j^{(n)}=\delta _{ij}\rho _i^{(n)}.$$
Then for any such $`\rho `$,
$`C_\rho (A^{(1)},B^{(2)})`$ $`=`$ $`\mathrm{tr}(\rho ^2A^{(1)}B^{(2)}\rho A^{(1)}\rho B^{(2)})`$
$`=`$ $`\mathrm{tr}\left(\left({\displaystyle p_i\rho _i^{(1)}A\rho _i^{(2)}}\right)\left({\displaystyle p_j\rho _j^{(1)}B\rho _j^{(2)}}\right)\right)`$
$`\mathrm{tr}\left(\left({\displaystyle p_i\rho _i^{(1)}A\rho _i^{(2)}}\right)\left({\displaystyle p_j\rho _j^{(1)}\rho _j^{(2)}B}\right)\right)`$
$`=`$ $`{\displaystyle }{\displaystyle }p_ip_j\mathrm{tr}(\rho _i^{(1)}A\rho _j^{(1)}\rho _i^{(2)}B\rho _j^{(2)}`$
$`\rho _i^{(1)}A\rho _j^{(1)}\rho _i^{(2)}\rho _j^{(2)}B)`$
$`=`$ $`{\displaystyle p_ip_j\mathrm{tr}_{(1)}(\rho _i^{(1)}A\rho _j^{(1)})\mathrm{tr}_{(2)}(\rho _i^{(2)}B\rho _j^{(2)}\rho _i^{(2)}\rho _j^{(2)}B)}`$
$`\mathrm{or}`$ $`{\displaystyle p_ip_j\mathrm{tr}_{(1)}(\rho _i^{(1)}A\rho _j^{(1)}\rho _i^{(1)}\rho _j^{(1)}A)\mathrm{tr}_{(2)}(\rho _i^{(2)}B\rho _j^{(2)})},`$
which is guaranteed to vanish if the $`\rho _i^{(n)}`$ are projectors.
Now consider cases where the density matrix is a mixture like
$$\frac{1}{2}|uuuu|+\frac{1}{8}|(u+d)(u+d)(u+d)(u+d)|.$$
Here, it can be shown that these mixtures produce non-zero $`C_\rho (A,B)`$, because their expansion into projector traces may be used to extract off-diagonal elements of observables $`A,B`$ in suitable local bases. Take for instance the two local operators on separate spin $`\frac{1}{2}`$ bases, $`A=B=\sigma _2`$, which is a unitary transformations of the equal-weight matrix, $`\sigma _3`$ used previously. The terms in the expansion of $`C`$ above are non-vanishing only if non-parallel, non-orthogonal projectors are used. Thus we evaluate the two non-vanishing terms, with local projectors $`|u+du+d|/2,|uu|`$ giving
$$u+d|uu|A|u+d(u|B|u+du+d|uu+d|B|uu|u+d)$$
$$=(i)(ii)=2,$$
$$u|u+du+d|A|u(u+d|B|uu|u+du|B|u+du+d|u)$$
$$=(i)(i+i)=2.$$
These terms do not cancel, and the other two terms in the expansion are zero, so the alternative covariance $`C_\rho `$ is non-zero even though we are dealing with a disentangled state, according to the usual terminology.
From this we deduce that the measure $`E`$ based on alternative covariance is not a measure of entanglement (according to the conditions provided by Vedral and Plenio), when mixed states are encountered, since it violates the first condition for such measures.
### 4.3 Local purification procedures
The third condition on entanglement measures proposed by Vedral and Plenio is that the entanglement of a state should not increase under the three types of purification processes (LGM, CC, and PS).
Let us consider all possible LGM, CC and PS measurement operations which act on a disentangled state
$$\rho _0=|uuuu|$$
and yield states which exhibit non-zero quantum covariance $`C_\rho (A^{(1)},B^{(2)})`$, such as
$$\rho _1=\frac{1}{2}|uuuu|+\frac{1}{8}|(u+d)(u+d)(u+d)(u+d)|$$
The change from $`\rho _0`$ to $`\rho _1`$ may be performed by the classically correlated set of local measurements $`V_1\mathrm{}V_8`$, where
$`V_1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|uuuu|={\displaystyle \frac{1}{\sqrt{2}}}|uu||uu|`$
$`V_2`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{2}}}|(u+d)(u+d)uu|={\displaystyle \frac{1}{2\sqrt{2}}}|(u+d)u||(u+d)u|`$
together with $`V_3,\mathrm{},V_8`$ which have $`ud|,du|,dd|`$ in place of $`uu|`$. Thus $`V_i\rho _0V_i^{}=\rho _1`$, and $`V_i^{}V_i1`$, so this represents a complete measurement.
Because this set of operations $`V_i`$ is local, we would expect any measure of entanglement to not increase under these operations. However, the state which results from this procedure does in fact exhibit non-zero $`C_\rho `$, as was demonstrated earlier. Therefore we have found a complete local general measurement for which the entanglement increases based upon alternative covariance. Thus the unitarily maximised $`C_\rho `$ does not satisfy the conditions for a measure of mixed state entanglement proposed by Vedral and Plenio.
## 5 Conclusion
We have shown that for pure two 2-state systems, the maximised covariance agrees with other measures of entanglement in specifying which states are disentangled and which are maximally entangled. For subspaces of larger dimension, the situation is less clear-cut since there are ambiguities in the process of selecting eigenvalues for the local operators, but it is possible to obtain information about the degree of higher-order correlations of two subsystems in this way.
The variance and covariance could be used to indicate the best measurements to make to detect entanglement of bipartite systems, by locating the unitary transformation which produces maximum covariance. However, the problem of mixed state separability measures is not resolved by the correlation functions $`C`$ and the covariance, as we have used them in this paper. Nevertheless the link between alternative variance and the conjugate variable derivatives for position, momentum, energy and time is intriguing, and might be extended to other quantised models involving conjugate variables. The inequality $`\mathrm{var}(X)C(X,X)`$, with equality for pure states, should be applicable to situations where observations are made upon impure states.
It is interesting to observe that for pure density matrices $`\rho `$, both the average value and the variance of an operator $`F`$ are symmetrical under the interchange of $`\rho `$ and $`F`$, when expressed in the form
$$F=\mathrm{tr}(\rho F),\mathrm{var}(F)=C_\rho (F,F)=\mathrm{tr}(\rho ^2F^2\rho F\rho F).$$
They motivate the suggestion that a quantum state $`|\psi `$ is best represented not as a ket, but as a projection operator $`P_\psi =|\psi \psi |`$. Any measurement process on the state is relational in the sense of Rovelli : the notion of state vector reduction is replaced by the notion that every new measurement process requires a new Hilbert space to be defined, with new operators corresponding to the new state and any successive measurement made upon the state. The basic idea is that a measurement treats the state information in a comparative sense, with meaning only in relation to the operator corresponding to the observable quantity being measured.
The process of state reduction to the eigenstate of some observable under a measurement is an apparent one resulting not from the action of the operator corresponding to the measurement, but from some other physical process which occurs in the measurement, such as a filtering process of some kind (like the case of a Stern-Gerlach experiment), causing later measurements on the physical system to be represented in a completely new basis by new state and measurement operators. The adoption of the projector as a unit of relational information does not change any predictions of quantum mechanics in terms of average values<sup>3</sup><sup>3</sup>3The paradox of Schroedinger’s Cat is resolved by representing the cat’s state and the measurement on an equal footing as projection operators in the (relational) basis for that measurement. Traces are taken to predict average real-number values, but at no stage does one say ‘the cat is in a state’, since the state projection operator $`P_\psi `$ only has relational significance in this interpretation..
## Appendix A The Majorana-Penrose Representation of Symmetrised States
It is well known that a spin state $`|j,m`$ can be obtained by forming a fully symmetric direct product of $`n=2j`$ spin 1/2 states. Denoting the spin 1/2 states by $`|u=|1/2,1/2,|d=|1/2,1/2`$, the spin $`j`$ states are given by:
$`|j,j`$ $`=`$ $`|\stackrel{2j}{\stackrel{}{uu\mathrm{}u}}`$
$`|j,j1`$ $`=`$ $`[|du\mathrm{}u+ud\mathrm{}u+\mathrm{}+uu\mathrm{}d]/\sqrt{2j}`$
$`\mathrm{}`$ $`\mathrm{}`$
$`|j,m`$ $`=`$ $`{\displaystyle \underset{\mathrm{sym}}{}}|(u)^{j+m}(d)^{jm}/\sqrt{{}_{}{}^{2j}C_{j+m}^{}}`$ (14)
$`\mathrm{}`$ $`\mathrm{}`$
$`|j,j`$ $`=`$ $`|dd\mathrm{}d.`$
In the Majorana-Penrose representation these states are mapped onto the unit sphere, with stereographic projection $`z`$ taken from the South Pole, onto the complex plane by making the association:
$$|j,m\sqrt{{}_{}{}^{2j}C_{j+m}^{}}z^{j+m}.$$
(15)
This association may be rewritten in terms of the spin-$`\frac{1}{2}`$ basis in an elegant way which emphasises that the $`z^i`$ powers are proportional to a symmetrised sum where $`i`$ spin-$`\frac{1}{2}`$ terms of spin $`|u`$, and $`(ni)`$ spin $`|d`$ terms are selected, namely
$${}_{}{}^{2j}C_{j+m}^{}z^{j+m}|u\mathrm{}ud\mathrm{}d+\mathrm{}+|d\mathrm{}du\mathrm{}u$$
(16)
where the sum is taken over all ways of selecting $`i`$ spin-$`\frac{1}{2}`$ elements from the $`2j+1`$ elements in the spin-$`j`$ space.
Thus from a general spin state $`|\psi =_m\psi _{jm}|jm`$ we can define a polynomial
$$p(z)=a_{2j}z^{2j}+a_{2j1}z^{2j1}+\mathrm{}+a_01,$$
where $`a_{j+m}\sqrt{{}_{}{}^{2j}C_{j+m}^{}}\psi _{jm}z^{j+m}`$. The roots of this polynomial<sup>4</sup><sup>4</sup>4The action of exponential functions of angular momentum operators on Majorana-Penrose polynomials are amusing. We quote them without proof: (i) $`\mathrm{exp}(\xi J_z).p(z)=\mathrm{exp}(\xi j)p(z\mathrm{e}^\xi );`$ (ii) $`\mathrm{exp}(\xi J_{}).p(z)=p(z+\xi );`$ (iii) $`\mathrm{exp}(\xi J_+).p(z)=z^{2j}p(z/(\xi z+1))`$. give $`2j`$ points on the stereographic plane and, correspondingly, a general spin $`j`$ state, being some superposition over $`m`$ values, will be described by $`2j`$ distinct points on the Poincaré sphere, obtained by stereographic projection of the roots placed on the horizontal $`x`$-$`y`$ plane. It is worth remarking that another spin state $`|\psi ^{}`$, represented by another polynomial,
$$p^{}(z)=a_{2j}^{}z^{2j}+a_{2j1}^{}z^{2j1}+\mathrm{}+a_0^{}1,$$
yields a scalar product,
$$\psi ^{}|\psi =\underset{m}{}a_{j+m}a_{j+m}^{}/^{2j}C_{j+m}.$$
In the case where the degree $`m`$ of the polynomial is less than $`2j`$, $`2jm`$ additional points at the South pole (the projective point) are added to the projection, corresponding to ‘roots at infinity’. Figure 5 depicts an example of stereographic projection for the spin-$`1`$ state represented by the quadratic polynomial
$$p(z)=(zi)(z+i)=z^21[|11+|11]/\sqrt{2}.$$
(17)
By contrast, for spin 1, the maximum and minimum spin states $`|1,1`$ and $`|1,1`$ correpond to two repeated points at the North and South Pole, respectively; on the other hand the pure intermediate spin state $`|1,0=|ud+du/\sqrt{2}`$ corresponds to one point at the North pole and the other at the South pole. The Bell states, $`|1,\pm |uu\pm dd/\sqrt{2}`$ are also antipodal points, but are located around the equator, specifically at coordinates (1,0,0)& (-1,0,0) and (0,1,0) & (0,-1,0). Under all accepted measures of entanglement, the spin zero state $`|0,0=|uddu/\sqrt{2}`$, the Bell states $`|1,\pm `$ and $`|1,0`$ (in fact, all local unitary transformations of these states) are maximally entangled. This suggests that taking the two points as “far apart as one another” on the Poincaré sphere is one possible way of producing maximal entanglement.
Another point of interest is that for these spin 1 states, the rotationally invariant dispersion measure,
$$(\mathrm{\Delta }J)^2\stackrel{}{J}.\stackrel{}{J}\stackrel{}{J}.\stackrel{}{J},$$
attains the maximum value of 2, because $`\stackrel{}{J}`$ vanishes; so it also suggests that the quantity $`(\mathrm{\Delta }J)^2/\stackrel{}{J}^2`$ might serve as a way of characterising the entanglement of the individual 1/2 spins that make up the $`j`$ state. In fact the dispersion in $`J`$ is equivalent to dispersion in any of the components $`J^{(i)}`$ and to the covariance of any two components. This is because the total angular momentum is $`J=J^{(1)}+J^{(2)}+\mathrm{}+J^{(n)}`$ and all states on the Poincaré sphere are symmetrised; therefore the values taken by all the local $`J^{(i)}`$ are the same. Thus when acting on these symmetrised states, $`J=nJ^{(i)}`$, for all $`i=1,2,\mathrm{},n`$, and the variance of any local $`J^{(i)}`$ is just as good an entanglement measure. We can also take any two subspaces $`i,j`$ and evaluate the covariance of $`J^{(i)}`$ and $`J^{(j)}`$; the space being symmetrised, any $`i,j`$ (including $`i=j`$) may be taken! Thus the symmetrised states of maximised $`(\mathrm{\Delta }J)^2`$ have maximal covariance of the local $`J^{(i)}`$, as well. Finally, since we are dealing with spin 1/2 states, $`J^{(i)}`$ have only have two distinct eigenvalues, and hence these local operators are directly proportional to the equal weight local operators of dimension two defined in section 3.
We may use these considerations to find the symmetrised states of maximum dispersion for arbitrarily high $`J`$ \- these states will have the maximum covariance of any two local $`J^{(i)}`$. From the result in section 4, the local covariance attains the maximum value of 1 whenever the reduced density matrices for the subspace on which those local operators jointly act are of maximum entropy. For a spin $`\frac{1}{2}`$ space, the reduced density matrices must be of the form $`\rho _i=\frac{1}{2}[|u_iu_i|+|d_id_i|]`$.
For general $`j`$-values, observe that disentangled states on the Poincaré sphere are fully factorisable since the states are symmetrised. Clearly a fully factorisable state will have zero local operator variance for any operator with the local state as an eigenstate. Their Poincaré sphere representation simply consists of $`n=2j`$ repeated roots, since the general symmetrised, factorised state is a product of $`n`$ kets
$$(a|u+b|d))\mathrm{}(a|u+b|d))\mathrm{}(az+b)^n$$
because of the $`z`$-symmetrisation properties for spin $`n`$ systems.
Next it is useful to ask which states maximise $`(\mathrm{\Delta }J)^2`$ if this is to serve as a possible indication of entanglement. When $`j=1`$, as we have seen in section 3, the states are represented on the Poincaré sphere by two diametrically opposed points, one example being the state $`|uu+dd/\sqrt{2}`$. When $`j=3/2`$ the states having maximum $`(\mathrm{\Delta }J)^2`$ are those states which maximise the covariance of any two local operators. As was seen such states must be expressible as symmetrised unitary transformations of $`(|uu+|dd)/\sqrt{2}`$ in the local basis for the two operators in question. Thus the overall states must be expressible, by symmetry, as global rotations of $`|uuu+ddd/\sqrt{2}`$. These are the ‘triangular’ states, namely those represented by 3 points on the Poincaré sphere arranged in an equilateral triangle around any great circle—a global unitary transformation (which preserves symmetrisation) simply rotates this configuration around the sphere. Choosing the circle to lie equatorially, a polynomial producing such roots is $`p(z)=z^3+1`$, which corresponds to the state $`|\psi =[|3/2,3/2+|3/2,3/2]/\sqrt{2}=|uuu+ddd/\sqrt{2}.`$
When $`j=2`$ the symmetrised states of maximum dispersion are the states with a tetrahedral representation on the sphere. Choosing one of the apices of the tetrahedron at the North Pole, we arrive at the polynomial
$$p(z)=z(z+\sqrt{2})(z\sqrt{2}\mathrm{e}^{i\pi /3})(z\sqrt{2}\mathrm{e}^{i\pi /3})=z^4+2\sqrt{2}z,$$
which corresponds to the maximally entangled state
$$|\psi =[|2,2+\sqrt{2}|2,1]/3=|uuuu+uddd+dudd+ddud+dddu/\sqrt{5}.$$
Note that here, as in the previous case, the maximal $`|\psi `$ lead to vanishing mean values $`\psi |\stackrel{}{J}|\psi .`$
For $`j=5/2`$ there are two classes of states with maximum dispersion, with slightly different geometries. The first of these is pyramidal with one apex at the North pole and the other four apices at equal latitude (or any global rotation of this state), while the second has one point at the North pole another at the South Pole and the remaining three points distributed equilaterally on the equator. Hence the first configuration corresponds to the polynomial $`p(z)=z^5+z^4/\sqrt{3}`$ or the state $`|\psi =[|5/2,5/2+\sqrt{5/3}|5/2,3/2]/\sqrt{8/3}`$, whereas the second configuration leads to $`p(z)=z^4+z`$, or the state $`[|5/2,5/2+|5/2,3/2]/\sqrt{2}.`$ Both choices have maximum $`(\mathrm{\Delta }J)^2`$ and so the total spin variance is unable to discriminate between them. However, local equal-weight measures are able to discriminate between these states, so it seems that in more complicated cases entanglement is most naturally described by keeping the covariance of local operators in mind.
In order to use the (rotationally invariant) dispersion as a measure of pure state entanglement, one needs to consider that the symmetrised nature of the state is a reflection of the choice of basis. Thus one might define the variance-entanglement for a general state as the variance of the symmetrised form of the state, under an appropriate local unitary transformation. However, it is not always possible to symmetrise an arbitrary state with local unitary transformations in spaces of spin 1 or higher. Thus it seems that the dispersion is not a perfect entanglement measure, although it does give an indication of the degree of entanglement of these particular states, because of its connection with the covariance through symmetrisation.
## Appendix B An Integrity Basis for Density Matrix Invariants
Consider $`\rho `$ for a composite $`N^{(1)}N^{(2)}`$-dimensional system as an object transforming under $`U(N^{(1)})\times U(N^{(2)})`$ like $`\{\overline{1}\}\{1\}\times \{\overline{1}\}\{1\}`$, where $`\{\overline{1}\}\{1\}`$ denotes the reducible $`N^2`$ representation of $`U(N)`$; in tensor notation we can write $`\rho `$ in the form $`\rho _{ai}^{bj}`$, where early Latin letters refer to the first unitary group and the later letters stand for the second group. Here we want to count the number of $`U(N^{(1)})\times U(N^{(2)})`$ singlets $`S_n`$ in the symmetrised product $`\rho ^n`$. Thus in the notation where representations are labelled by partitions ,
$$\rho ^n(\{\overline{1}\}\{1\}\times \{\overline{1}\}\{1\})\{n\}\underset{\kappa \lambda \mu \nu \{n\}}{}\{\overline{\kappa }\}\{\lambda \}\times \{\overline{\mu }\}\{\nu \}.$$
But $`\{\overline{\kappa }\}\{\lambda \}`$ can only contain a singlet if $`\kappa \lambda `$. Moreover $`\kappa =\lambda `$ and $`\mu =\nu `$ have to be respectively $`N^{(1)}`$ and $`N^{(2)}`$ part partitions, say “$`\kappa _{N^{(1)}}n`$”, etc., otherwise $`\{\kappa \}`$ vanishes in $`U(N^{(1)})`$. Therefore
$$\rho ^n|_{\{0\}\times \{0\}}\underset{\kappa _{N^{(1)}}n,\lambda _{N^{(2)}}n}{}\{\overline{\kappa }\}\{\kappa \}\times \{\overline{\lambda }\}\{\lambda \}.$$
But it is known that $`\alpha \alpha \{n\}`$ for any $`\alpha n`$ and moreover the order is immaterial, because the Clebsch series is symmetric. Thus we have only to count the appropriate partitions,
$$S_n=|\rho ^n|_{\{0\}\times \{0\}}=|\{\kappa ,\lambda :\kappa _{N^{(1)}}n,\lambda _{N^{(2)}}n\}|.$$
This is not easy to work out in the general case, but is relatively simple for the case $`N^{(1)}=N^{(2)}=2`$. When $`n`$ is even or odd, the partitions are:
$$n=2k:\{2k,0\},\{2k1,1\},\mathrm{}\{k,k\}$$
$$n=2k+1:\{2k+1,0\},\{2k,1\},\mathrm{},\{k+1,k\}.$$
Now the generating function for invariants of order $`n`$ in $`\rho `$ is written $`F(q):=_{n=0}^{\mathrm{}}S_nq^n`$. Including the even and odd cases,
$$F(q)=\underset{k=0}{\overset{\mathrm{}}{}}[(k+1)^2q^{2k}+(k+1)^2q^{2k+1}]=\frac{1+q^2}{(1q^2)^2(1q)}.$$
The denominator of $`F(q)`$ is crucial for its interpretation: we can recognize that for the $`2\times 2`$ case invariants are freely generated by two quadratic factors (namely $`(1q^2)^2`$) and one linear factor (viz. $`(1q)`$), but there is an extra quadratic factor in the numerator which may only be used once. We may associate these factors with
$$\mathrm{Linear}:\mathrm{tr}\rho =\mathrm{tr}_{(1)}\mathrm{tr}_{(2)}\rho =1\mathrm{anyway},$$
$$\mathrm{Quadratic}:\chi _1\mathrm{tr}_{(1)}(\mathrm{tr}_{(2)}\rho )^2=\rho _{ai}^{bi}\rho _{bj}^{aj},$$
$$\&\chi _2\mathrm{tr}_{(2)}(\mathrm{tr}_{(1)}\rho )^2=\rho _{ai}^{aj}\rho _{bj}^{bi},$$
$$\mathrm{Extra}\mathrm{quadratic}:ϵ_{aa^{}}ϵ^{bb^{}}ϵ_{ii^{}}ϵ^{jj^{}}\rho _{bj}^{ai}\rho _{b^{}j^{}}^{a^{}i^{}}.$$
The last of these invariants is obviously related to $`\chi _1,\chi _2`$ and tr($`\rho ^2`$); however we do not get a new invariant from its square because it is then expressible entirely as products of $`\chi _1,\chi _2`$, tr $`\rho ^2`$ and (tr $`\rho )^21`$. Thus effectively tr($`\rho ^2`$) is only allowed once.
We conclude that any local unitary invariant entanglement measure in the 2$`\times `$2 case must take the form:
$$E(\rho ):=F(\chi _1,\chi _2)+\mathrm{tr}(\rho ^2).G(\chi _1,\chi _2),$$
where $`F,G`$ are functions which depend on the way $`E(\rho )`$ is defined. We think that this result must be useful for classifying entanglement.
## Acknowledgement
We thank V. Vedral for helpful discussions by electronic mail. |
warning/0001/hep-ph0001332.html | ar5iv | text | # 1 Values of the parameter A₂ as function of Q², obtained using eq. () and quark distribution in the pion suggested in [].
RUB-TP2-01/00
Partonic structure of $`\pi `$ and $`\rho `$ mesons from data on hard exclusive production of two pions off nucleon
B. Clerbaux<sup>a,∗</sup> and M.V. Polyakov<sup>b,c</sup>
<sup>a</sup> Inter-University Institute for High Energies (U.L.B.), Brussels, Belgium
$``$ Since January 2000 at C.E.R.N., Geneva, Switzerland
e-mail: Barbara.Clerbaux@cern.ch
<sup>b</sup>Petersburg Nuclear Physics Institute, 188350, Gatchina, Russia
<sup>c</sup> Institut für Theoretische Physik II, Ruhr-Universität Bochum, D-44780 Bochum, Germany
e-mail: maximp@tp2.ruhr-uni-bochum.de
## Abstract
We fitted the $`\pi \pi `$ mass distribution in the range $`0.5M_{\pi \pi }1.1`$ GeV measured in hard exclusive positron-proton reactions at HERA by the form dictated by QCD at leading twist level. Extracted parameters are related to valence quark distribution in the pion, and to the pion and $`\rho `$ meson distribution amplitudes. We obtain, for the first time, a measurement of the second Gegenbauer coefficient of the $`\rho `$ meson distribution amplitude: $`a_2^{(\rho )}=0.10\pm 0.20`$ for a photon virtuality of $`Q^2=21.2`$ GeV<sup>2</sup>.
1. Introduction
Owing to QCD factorisation theorem for hard exclusive reactions the dependence of the amplitude of the reaction <sup>1</sup><sup>1</sup>1With $`Q^22(pq)\mathrm{\Lambda }_{\mathrm{QCD}}^2`$ and $`M_T^2,M_T^{}^2,(p^{}p)^2,M_{\pi \pi }^2Q^2.`$
$$\gamma _L^{}(q)+T(p)\pi ^+\pi ^{}+T^{}(p^{})$$
(1)
with longitudinally polarised virtual photon on the di-pion mass $`M_{\pi \pi }`$ factors out in an universal (independent of the target) factor. At leading order in $`\alpha _s(Q^2)`$ this factor has the form :
$$𝒜(M_{\pi \pi })_0^1\frac{dz}{z}\mathrm{\Phi }^I(z,\zeta ,M_{\pi \pi };Q^2).$$
(2)
Here $`\mathrm{\Phi }^I(z,\zeta ,M_{\pi \pi };Q^2)`$ is the two-pion light cone distribution amplitude ($`2\pi `$DA) , which depends on $`z`$–longitudinal momentum carried by the quark, $`\zeta `$ characterising the distribution of longitudinal momentum between the two pions<sup>2</sup><sup>2</sup>2For detailed definition of kinematical variables $`z`$ and $`\zeta `$ see refs. ., and the invariant mass of produced pions $`M_{\pi \pi }`$, the superscript $`I`$ standing for isospin of produced pions ($`I=0,1`$). The dependence on the virtuality of the incident photon $`Q^2`$ is governed by the ERBL evolution equation .
For the process (1) at small $`x_{Bj}=Q^2/2(pq)`$ the production of two pions in the isoscalar channel is strongly suppressed relative to the isovector channel , because the former is mediated by $`C`$-parity odd exchange. At asymptotically large $`Q^2`$ QCD predicts the following simple form for the isovector $`2\pi `$DA :
$$\mathrm{\Phi }_{asym}^{I=1}(z,\zeta ,M_{\pi \pi })=6z(1z)(2\zeta 1)F_\pi (M_{\pi \pi }),$$
(3)
where $`F_\pi (M_{\pi \pi })`$ is the pion electro-magnetic (e.m.) form factor in time-like region, measured with high precision in low energy experiments . From eqs. (2,3) we conclude that at asymptotically large $`Q^2`$ QCD predicts unambiguously the shape of the di-pion mass distribution. Asymptotically the dependence of the amplitude on $`M_{\pi \pi }`$ has the form:
$$𝒜_{asym}\beta (M_{\pi \pi })F_\pi (M_{\pi \pi })\mathrm{cos}\theta ,$$
(4)
where $`\beta (M_{\pi \pi })=\sqrt{1\frac{4m_\pi ^2}{M_{\pi \pi }^2}}`$ is the velocity of pions in their centre of mass system (cms) and $`\theta `$ is the angle between the directions of the positive pion and the momentum of produced $`\pi ^+\pi ^{}`$ system in the $`\pi ^+\pi ^{}`$ cms. This angle is related to $`\zeta `$ in the following way:
$`\mathrm{cos}\theta ={\displaystyle \frac{2\zeta 1}{\beta (M_{\pi \pi })}}.`$ (5)
The corresponding di-pion mass distribution has asymptotically the form
$$\frac{dN(M_{\pi \pi })}{dM_{\pi \pi }}M_{\pi \pi }\beta (M_{\pi \pi })^3|F_\pi (M_{\pi \pi })|^2.$$
(6)
The asymptotic shape for any di-meson production (mesons $`M_1`$, $`M_2`$) was derived in , where it was related to the cross section of $`e^+e^{}M_1,M_2`$ at low $`\sqrt{s}`$.
At non-asymptotic $`Q^2`$ values, the $`2\pi `$DA deviates from its asymptotic form (3). This deviation can be described by a few parameters which can be related to quark distributions (skewed and usual) in the pion and to distribution amplitudes of mesons ($`\pi ,\rho `$, etc… ), for details see .
2. Deviations from the asymptotic form
The first non-trivial deviation from the asymptotic form of the $`2\pi `$DA occurs in $`P`$ and $`F`$ waves . Generically their effect on $`M_{\pi \pi }`$ dependence of the hard amplitude can be written as:
$`𝒜(M_{\pi \pi })`$ $``$ $`\beta (M_{\pi \pi })e^{i\delta _1}|F_\pi (M_{\pi \pi })|\left(1+D_1(M_{\pi \pi })\right)P_1(\mathrm{cos}\theta )`$ (7)
$`+`$ $`\beta (M_{\pi \pi })^3e^{i\delta _3}D_2(M_{\pi \pi })P_3(\mathrm{cos}\theta ),`$
where $`P_l(\mathrm{cos}\theta )`$ are Legendre polynomials and $`\delta _1(M_{\pi \pi })`$ and $`\delta _3(M_{\pi \pi })`$ are the $`P`$-wave and the $`F`$-wave $`\pi \pi `$ scattering phase shifts, which are well known from low-energy experiments. The functions $`D_{1,2}(M_{\pi \pi })`$ describe the deviation of the amplitude’s $`M_{\pi \pi }`$ dependence from the asymptotic form. These functions are real and can be parametrised as:
$`D_1(M_{\pi \pi },Q^2)`$ $`=`$ $`A_1(Q^2)e^{\overline{b}_1M_{\pi \pi }^2}{\displaystyle \frac{6m_\pi ^2}{M_{\pi \pi }^2}}A_2(Q^2)e^{\overline{b}_2M_{\pi \pi }^2}`$
$`D_2(M_{\pi \pi },Q^2)`$ $`=`$ $`A_2(Q^2)e^{b_3M_{\pi \pi }^2}.`$
The dependence of $`A_{1,2}(Q^2)`$ on $`Q^2`$ is governed by the QCD evolution and in leading order is given by:
$`A_{1,2}(Q^2)=A_{1,2}(\mu _0)\left({\displaystyle \frac{\alpha _s(Q^2)}{\alpha _s(\mu _0)}}\right)^{50/(996n_f)}.`$ (8)
With increasing $`Q^2`$, the parameters $`A_{1,2}(Q^2)`$ go logarithmically to zero and one reproduces the asymptotic formula (4). The parameters $`A_{1,2}(Q^2)`$ are directly related to partonic structure of $`\pi `$ and $`\rho `$ mesons, see section 4 and ref .
The parameter $`b_3`$ is $`Q^2`$-independent but is $`M_{\pi \pi }`$ dependent. The latter dependence is fixed by $`\pi \pi `$ scattering phase shifts (see for derivation ):
$`b_3(M_{\pi \pi })=\overline{b}_3+\text{Re}{\displaystyle \frac{M_{\pi \pi }^2}{\pi }}{\displaystyle _{4m_\pi ^2}^{\mathrm{}}}𝑑s{\displaystyle \frac{\delta _3(s)}{s^2(sM_{\pi \pi }^2i0)}}.`$ (9)
In above equations $`\overline{b}_i`$ are subtraction constants of corresponding dispersion relations for functions $`D_{1,2}(M_{\pi \pi })`$, see details in .
Using expression (7) we can derive the form of the $`M_{\pi \pi }`$ distribution:
$`{\displaystyle \frac{dN(M_{\pi \pi })}{dM_{\pi \pi }}}=N[`$ $`\beta (M_{\pi \pi })^3M_{\pi \pi }|F_\pi (M_{\pi \pi })|^2\left(1+D_1(M_{\pi \pi })\right)^2+`$
$`{\displaystyle \frac{3}{7}}M_{\pi \pi }\beta (M_{\pi \pi })^7D_2(M_{\pi \pi })^2+`$
$`\mathrm{higher}\mathrm{waves}l5],`$
where the higher partial waves can be safely neglected.
3. Angular distributions of produced pions
Another way to obtain sensitivity to the non-asymptotic parameters $`A_{1,2}(Q^2)`$ is to study the two pion angular distributions. From the expression of the amplitude (7) one can derive the $`M_{\pi \pi }`$ dependence of intensity densities $`Y_l^m(\theta ,\varphi )`$ defined as:
$`{\displaystyle \frac{d}{dM_{\pi \pi }}}Y_l^m=M_{\pi \pi }\beta (M_{\pi \pi }){\displaystyle _1^1}d(\mathrm{cos}\theta ){\displaystyle _0^{2\pi }}𝑑\phi Y_l^m(\theta ,\phi )|𝒜(M_{\pi \pi },\theta ,\phi )|^2,`$ (11)
where $`Y_l^m(\theta ,\phi )`$ are spherical harmonics, $`\phi `$ is the azimuthal angle between the pion decay plane and plane formed by momentum of the virtual photon and total momentum of produced $`\pi ^+\pi ^{}`$ system.
For the first nontrivial intensity density $`Y_2^0`$ we have:
$`{\displaystyle \frac{d}{dM_{\pi \pi }}}Y_2^0`$ $``$ $`M_{\pi \pi }[\beta (M_{\pi \pi })^3|F_\pi (M_{\pi \pi })|^2(1+D_1(M_{\pi \pi }))^2`$
$`+`$ $`{\displaystyle \frac{9}{7}}\beta (M_{\pi \pi })^5|F_\pi (M_{\pi \pi })|\left(1+D_1(M_{\pi \pi })\right)D_2(M_{\pi \pi })\mathrm{cos}\left[\delta _1(M_{\pi \pi })\delta _3(M_{\pi \pi })\right]`$
$`+`$ $`{\displaystyle \frac{2}{7}}\beta (M_{\pi \pi })^7D_2(M_{\pi \pi })^2].`$
The combination $`Y_0^0\sqrt{5}/2Y_2^0`$ is especially sensitive to the deviations from asymptotic form because it is exactly zero asymptotically:
$`{\displaystyle \frac{d}{dM_{\pi \pi }}}Y_0^0{\displaystyle \frac{\sqrt{5}}{2}}Y_2^0`$ $``$ $`M_{\pi \pi }[\beta (M_{\pi \pi })^5|F_\pi (M_{\pi \pi })|(1+D_1(M_{\pi \pi }))`$ (12)
$`\times `$ $`D_2(M_{\pi \pi })\mathrm{cos}\left[\delta _1(M_{\pi \pi })\delta _3(M_{\pi \pi })\right]`$
$``$ $`{\displaystyle \frac{1}{9}}\beta (M_{\pi \pi })^7D_2(M_{\pi \pi })^2].`$
Let us note, however that the expressions for intensity densities obtained above are based on leading twist expression for the amplitude (7), in particular, the contribution of transversely polarised photon is neglected (it contributes at higher twist level). Since in the combination of intensity densities (12) the leading twist contribution tends to cancel, the effect of higher twists (for which the cancellation is not expected) can be numerically large.
The study of angular distribution allows to make separation<sup>3</sup><sup>3</sup>3Under assumption of $`s`$channel helicity conservation (SCHC), which holds with good accuracy experimentally . of productions by longitudinally (leading twist) and transversely (higher twist) polarised photons. Indeed, we can form combinations of intensity densities which get contributions only from $`\gamma _L^{}`$, so we minimise the contributions of higher twists. One of such combinations <sup>4</sup><sup>4</sup>4Strictly speaking this combination gets also contribution from transverse polarisation of the photon, but it occurs only in $`F`$-wave and therefore is expected to be small. is, for example,
$`{\displaystyle \frac{d}{dM_{\pi \pi }}}Y_0^0+\sqrt{5}Y_2^0`$ $``$ $`M_{\pi \pi }[\beta (M_{\pi \pi })^3|F_\pi (M_{\pi \pi })|^2(1+D_1(M_{\pi \pi }))^2`$
$`+`$ $`{\displaystyle \frac{6}{7}}\beta (M_{\pi \pi })^5|F_\pi (M_{\pi \pi })|\left(1+D_1(M_{\pi \pi })\right)`$
$`\times `$ $`D_2(M_{\pi \pi })\mathrm{cos}\left[\delta _1(M_{\pi \pi })\delta _3(M_{\pi \pi })\right]`$
$`+`$ $`{\displaystyle \frac{1}{3}}\beta (M_{\pi \pi })^7D_2(M_{\pi \pi })^2].`$
The details of the derivation and the analysis of the data will be presented elsewhere.
The advantage of the formalism presented in this section is that it maximally uses the information on pion e.m. form factor $`F_\pi (M_{\pi \pi })`$ and pion phase shifts $`\delta _l(M_{\pi \pi })`$ at low $`M_{\pi \pi }`$ which are very well known from low energy experiments. For example, from the known phase shifts $`\delta _1(M_{\pi \pi })`$ and $`\delta _3(M_{\pi \pi })`$ we conclude that the term proportional to $`\mathrm{cos}[\delta _1(M_{\pi \pi })\delta _3(M_{\pi \pi })]`$ changes the sign around $`M_{\pi \pi }=0.8`$ GeV, so to increase the sensitivity to this term it would be interesting to consider a “$`M_{\pi \pi }`$ asymmetry” for various observables:
$`\left({\displaystyle _{2m_\pi }^{0.8}}𝑑M_{\pi \pi }{\displaystyle _{0.8}^{M_{\mathrm{max}}}}𝑑M_{\pi \pi }\right){\displaystyle \frac{d}{dM_{\pi \pi }}}(\mathrm{an}\mathrm{observable}).`$ (14)
4. Expected values for the parameters
By crossing relations the parameter $`A_2`$ is related to the third moment of the valence quark distribution in the pion :
$`A_2(Q^2)={\displaystyle \frac{7}{6}}M_3(Q^2)={\displaystyle \frac{7}{6}}{\displaystyle _0^1}𝑑xx^2(u_{\pi ^+}\overline{u}_{\pi ^+}).`$ (15)
If the parametrisation suggested in is used for the quark distributions in the pion, we obtain the values of the parameter $`A_2(Q^2)`$ listed in Table 1.
The value of $`A_1`$ is constrained by the soft pion theorem :
$$A_1(Q^2)=a_2^{(\pi )}(Q^2)A_2(Q^2),$$
where $`a_2^{(\pi )}`$ is the second Gegenbauer coefficient of the pion distribution amplitude. Unfortunately this parameter is not very well measured. In ref. the value of $`a_2^{(\pi )}=0.19\pm 0.13`$ at $`Q^2=2.4`$ GeV<sup>2</sup> is quoted.
Additionally the parameters $`A_1(Q^2)`$ and $`\overline{b}_1`$ can be related to the second Gegenbauer coefficient $`a_2^{(\rho )}(Q^2)`$ of the $`\rho `$ meson distribution amplitude as :
$`a_2^{(\rho )}(Q^2)=A_1(Q^2)e^{\overline{b}_1M_{\pi \pi }^2}.`$ (16)
Up to now there is no direct experimental information about $`a_2^{(\rho )}(Q^2)`$.
The $`M_{\pi \pi }`$ dependence of the parameter $`b_3`$ is estimated to be weak in the range of $`0.5<M_{\pi \pi }<1.1`$ GeV and is then replaced by a constant $`\overline{b}_3`$. The parameters $`\overline{b}_i`$ were estimated in the instanton model of QCD vacuum to be around the following values:
$`\overline{b}_1`$ $`=`$ $`0.75\mathrm{GeV}^2`$
$`\overline{b}_2`$ $`=`$ $`0.75\mathrm{GeV}^2`$ (17)
$`\overline{b}_3`$ $`=`$ $`0.75\mathrm{GeV}^2.`$
5. Results of fits to the HERA data
Recently data on di-pion mass distribution in hard exclusive reaction was measured at HERA by the reaction
$$e+pe+p+\pi ^++\pi ^{}.$$
(18)
This section presents attempts to fit HERA data by the leading twist parametrisations described in the sections above.
In order to study the $`Q^2`$ evolution of the di-pion mass distribution three sets of data, the ZEUS photoproduction data ($`Q^2`$ $``$ 0) , the ZEUS low $`Q^2`$ data ($`0.25<Q^2<0.85`$ GeV<sup>2</sup> and the H1 high $`Q^2`$ data ($`2.5<Q^2<60`$ GeV<sup>2</sup> are used. The mean W value for these samples is around 70 GeV, where $`W`$ is the energy in the photon-proton cms ($`W^2`$ $``$ $`Q^2/x_{Bj}Q^2`$). For the three samples, the di-pion mass distribution <sup>5</sup><sup>5</sup>5 Note that the di-pion mass distributions measured at HERA are not corrected for the transverse photon production. Nevertheless for $`Q^2`$ greater than 2 GeV<sup>2</sup> the longitudinal cross section dominates the transverse cross section. was studied in the ranges $`0.55<M_{\pi \pi }<1.2`$ GeV, $`0.55<M_{\pi \pi }<1.1`$ GeV and $`0.50<M_{\pi \pi }<1.1`$ GeV, respectively. The main background to eq. (18) consists of events in which the proton diffractively dissociates into hadrons, and is estimated to be around 20 % for the ZEUS samples and around 10 % for the H1 sample. This background does not distort the di-pion mass distribution discussed in previous sections if the mass of recoiled baryonic system is much smaller than the hard scale $`Q^2`$. The contamination from the production of $`\omega `$ (decaying into $`\pi ^+\pi ^{}\pi ^0`$) and $`\varphi `$ (decaying into $`K^+K^{}`$ or $`\pi ^+\pi ^{}\pi ^0`$) mesons was estimated to be few percents for the ZEUS samples and 7 % for the H1 sample. This background is mainly situated at low mass $`M_{\pi \pi }<`$ 0.6 GeV.
Fig. 1 presents the di-pion mass distribution (black points) for six $`Q^2`$ values: the first distribution corresponds to the ZEUS photoproduction sample, the two following ones to the ZEUS low $`Q^2`$ sample and the three last ones to the H1 sample. For the H1 data, the distributions are corrected bin per bin for the production of $`\omega `$ and $`\varphi `$ mesons.
Firstly, we attempt to fit the HERA data with the simplest formula of the asymptotic form (6) where the di-pion mass distribution is directly proportional to the square of the pion e.m. form factor. This latter was recently precisely measured in low energy <sup>6</sup><sup>6</sup>6 At leading twist the energy dependence factories out. The shape of the di-pion mass distribution is then energy independent in the present formalism. experiments (see also ), and we use directly these experimental data for $`|F_\pi (M_{\pi \pi })|`$. The result of the fits is superposed to the data in Fig. 1 (black curves). The asymptotic form does not reproduce the distribution in photoproduction and at low $`Q^2`$ as expected<sup>7</sup><sup>7</sup>7The asymptotic form (6) is expected to work only in the hard regime $`Q^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, also let us stress once more that the form of di-pion mass distribution of eq. (S0.Ex5) also makes sense only in the hard regime., but describes well the data at high $`Q^2`$, the $`\chi ^2/ndf`$ of the fits for the six $`Q^2`$ bins of Fig. 1 being respectively 1625/25, 169/10, 69/10, 30.5/23, 16.6/21 and 11.6/18. We see that di-pion mass distributions in soft (low $`Q^2`$) and hard (large $`Q^2`$) regimes are essentially different and are described by different physics. Therefore the parametrisations designed for soft processes are not relevant for hard processes, and the values of the parameters extracted from large $`Q^2`$ data are not related directly to physical observables.
Secondly, we study the deviations from the asymptotic form of the di-pion mass distribution measured at HERA. As the formalism is valid in the hard regime $`Q^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$, only the H1 samples ($`Q^2`$ $`>`$ 2.5 GeV<sup>2</sup>) are considered. The parametrisation (S0.Ex5) is fitted to the data, with three free parameters: the normalisation $`N`$, and the parameters $`A_1`$ and $`A_2`$. For the parameters $`\overline{b}_1,\overline{b}_2`$, and $`\overline{b}_3`$ we use the values given by the calculation in the instanton model of QCD vacuum, see eq. (17). The result of the fits is presented in Fig. 2 as the black curves, and the values obtained for the parameters are listed in Table 2. The three fits have a $`\chi ^2/ndf`$ value smaller than one.
We observe an indication for a decrease of the parameter $`A_1`$ when $`Q^2`$ increases, what is expected from QCD evolution, see eq. (8). In general the values of parameters $`A_1`$ and $`A_2`$ are in agreement with theoretical expectations (see section 4). However the errors of the parameters $`A_1`$ and $`A_2`$ are large, due to limited statistic available. Nevertheless the sensitivity of the data to the parameter $`A_1`$ is encouraging. More precise data from the 1997-99 years data taking at HERA would bring considerably more accurate information on these parameters which are related directly to the valence quark distribution in the pion, and pion and $`\rho `$ meson distribution amplitudes.
6. Discussion of the results and conclusions
The analysis presented here is based on the leading order, twist-2 formalism. The power corrections (higher twist contributions) to the amplitude of the reaction (1) are systematically neglected. The size of higher twist corrections might be rather large , so that they constitute the main theoretical systematic uncertainty in our analysis. On general grounds the higher twist effects should be smaller for such observables as the shape of the di-pion mass distribution considered here, also the study of angular distribution of produced pions as discussed in section 3 can be used to minimise the contributions of higher twists. Another source of theoretical errors is that we fixed values of parameters $`\overline{b}_i`$ by the model values (17). Hopefully more precise data will allow us to measure these parameters.
To be on more safe side in respect to higher twist corrections we use the highest values of $`Q^2=21.2`$ GeV<sup>2</sup> available in the data sample to evaluate observables related to partonic structure of the pion and $`\rho `$ meson. The results for the third moment of valence quark distribution in the pion $`M_3`$ (see eq. (15)), the second Gegenbauer coefficients of the pion and $`\rho `$ meson distribution amplitudes ($`a_2^{(\pi )}`$ and $`a_2^{(\rho )}`$ respectively) at $`Q^2=21.2`$ GeV<sup>2</sup> are presented in Table 3.
Although, due to the limited statistic, the error bars for physical observables are large we see that the obtained values are in agreement with other experiments. This agreement is especially interesting because previously the restrictions for $`a_2^\pi `$ and $`M_3`$ were obtained from completely different measurements: the value of second Gegenbauer coefficient $`a_2^\pi `$ was obtained in ref. from analysis of data on $`\gamma \pi `$ transition form factor at large $`Q^2`$, whereas the third moment $`M_3`$ is restricted by $`\pi N`$ Drell-Yan data <sup>8</sup><sup>8</sup>8Lattice calculations of $`M_3`$ are in agreement with results of analysis of : $`0.09\pm 0.03`$ and $`0.10\pm 0.02`$ at $`Q^2=21.2`$. We see that the formalism to extract the partonic structure of the pion and $`\rho `$ mesons suggested here is complementary to already known methods.
Our analysis allowed us to obtain for the first time an experimental estimate of the second Gegenbauer coefficient of the $`\rho `$ meson distribution amplitude $`a_2^{(\rho )}=0.1\pm 0.2`$ at $`Q^2=21.2`$ GeV<sup>2</sup>. Unfortunately the precision of determination of $`a_2^{(\rho )}`$ is still low to discriminate between different model predictions for this quantity (QCD sum rule: $`0.12\pm 0.06`$ , $`0.05\pm 0.01`$ and instanton model: $`0.07`$ ). With increasing of statistics the accuracy of determination of the parameters related to the partonic structure of the mesons can be considerably improved. Also the studies of angular distributions of produced pions as discussed in section 3 can considerably increase the sensitivity to parameters of partonic structure of pions and their resonances.
As a final remark we note that the formalism developed here can be easily generalised to the case of production of other mesons in hard exclusive reactions (e.g. $`K\overline{K}`$, $`3\pi ,4\pi ,K\overline{K}\pi \pi `$, etc.). The study of meson mass and angular distributions of produced meson can provide us with rich information on partonic structure of these mesons and their resonances. Interesting prediction can be made for the asymptotic shape of meson cluster mass distribution for the reaction:
$`\gamma _L^{}+p(hadrons)+N^{},`$ (19)
where the mass of the hadron cluster $`M_X^2Q^2`$. In this case at asymptotically large $`Q^2`$ we have:
$`{\displaystyle \frac{dN}{dM_X^2}}`$ $``$ $`R(M_X),\mathrm{with}`$ (20)
$`R(\sqrt{s})`$ $`=`$ $`{\displaystyle \frac{\sigma (e^+e^{}hadrons)}{\sigma (e^+e^{}\mu ^+\mu ^{})}}.`$
It would be interesting to check experimentally this asymptotic formula and try to detect deviations from it. Generically, this deviation is described by the vacuum correlator of two light-ray opertors. We showed here that the asymtotic expression (20) works pretty well at low $`M_X`$ (for $`Q^2>7`$ GeV<sup>2</sup>) where the $`2\pi `$ channel dominates over all other hadronic channels, it would be interesting to check how good this formula works at higher $`M_X`$. Also it would be interesting to see whether the formula like (20) is applied also for multimeson partial cross section.
Acknowledgements: We gratefully acknowledge discussions with K. Goeke, B. Lehmann-Dronke, L. Mankiewicz, G. Piller, A. Schäfer, A.G. Shuvaev, W. Weise, and C. Weiss. Special thanks are due to Lyonya Frankfurt and Mark Strikman for suggestion to look into the problem and many valuable contributions. The work of B.C. is supported by the National Foundation for Scientific Reseach (FNRS). M.V.P. is supported by RFBR grant 96-15-96764, by DFG and BMBF. |
warning/0001/gr-qc0001096.html | ar5iv | text | # A Metric Theory of Gravity with Condensed Matter Interpretation
## 1 Introduction
There is a close analogy between condensed matter theory and gravity. It has been recognized that “effective gravity, as a low-frequency phenomenon, arises in many condensed matter systems” . This has been used to study Hawking radiation and the Unruh effect and vacuum energy for condensed matter examples. The general exchange of ideas with high energy physics, which “includes global and local spontaneous symmetry breaking, the renormalization group, effective field theory, solitons, instantons, and fractional charge and statistics” , is also worth to be mentioned.
The theory of gravity we present here suggests that this is not an accident, but the gravitational field is a medium in Newtonian space-time, described by usual condensed matter variables, with an interesting Lagrange formalism. Few general assumptions are sufficient to obtain a Lagrangian very close to GR, which fulfills the Einstein equivalence principle:
$$L=(8\pi G)^1(\mathrm{\Xi }g^{ii}\mathrm{{\rm Y}}g^{00})\sqrt{g}+L_{GR}+L_{matter}$$
After the derivation of the theory we consider quantization, some remarkable predictions, and compare the theory with other theories of gravity.
## 2 The Theory
The theory describes a classical medium in a Newtonian framework – Euclidean space and absolute time. But we prefer to present the theory in a formalism where the non-covariant terms are disguised as covariant, with the preferred coordinates considered as usual scalar fields $`X^\mu (x)`$. It is easy to transform a non-covariant Lagrangian $`L=L(T_{\mathrm{}}^{\mathrm{}},_\mu T_{\mathrm{}}^{\mathrm{}})`$ into a (formally) covariant form $`L=L(T_{\mathrm{}}^{\mathrm{}},_\mu T_{\mathrm{}}^{\mathrm{}},X_{,\nu }^\mu )`$.
The medium is described by steps of freedom typical for condensed matter theory. The gravitational field is defined by a positive density $`\rho `$, a velocity $`v^i`$, and a negative-definite symmetrical tensor field $`p^{ij}`$ which we name “pressure”. The effective metric $`g_{\mu \nu }`$ is defined algebraically by
$`\widehat{g}^{00}=g^{00}\sqrt{g}`$ $`=`$ $`\rho `$
$`\widehat{g}^{i0}=g^{i0}\sqrt{g}`$ $`=`$ $`\rho v^i`$
$`\widehat{g}^{ij}=g^{ij}\sqrt{g}`$ $`=`$ $`\rho v^iv^j+p^{ij}`$
This decomposition of $`g^{\mu \nu }`$ into $`\rho `$, $`v^i`$ and $`p^{ij}`$ is a variant of the ADM decomposition. The signature of $`g^{\mu \nu }`$ follows from $`\rho >0`$ and negative definiteness of $`p^{ij}`$.
The theory does not specify all properties of the medium, but only a few general properties – especially the conservation laws. The “material properties” of the medium, denoted by $`\phi ^m`$, remain unspecified. They are the matter fields. The complete specification – which includes the material laws of the medium – gives the theory of everything. The few general properties fixed here define a theory of gravity similar to GR. While it leaves the matter steps of freedom and the matter Lagrangian unspecified, it derives the Einstein equivalence principle.
In our covariant formalism the conservation laws may be defined as the Euler-Lagrange equations for the preferred coordinates. The related energy-momentum tensor
$$T_\mu ^\nu =\frac{L}{X_{,\nu }^\mu }$$
is not the same as in Noether’s theorem, but only equivalent. Now, we identify these conservation laws with the conservation laws we know from condensed matter theory. First, the Euler-Lagrange equation for time we identify with the classical continuity equation for the medium:
$$_t\rho +_i(\rho v^i)=0$$
(1)
The equations for the spatial coordinates we identify with the Euler equation:
$$_t(\rho v^j)+_i(\rho v^iv^j+p^{ij})=0$$
(2)
Note that we use here the identification of matter fields with material properties of the medium – we have no momentum exchange with external matter. The four conservation laws transform into the harmonic condition for the metric $`g_{\mu \nu }`$. Thus, they really look like equations for the preferred coordinates:
$$\mathrm{}X^\nu =_\mu (g^{\mu \nu }\sqrt{g})=0$$
Therefore, we assume that the conservation laws are proportional to the Euler-Lagrange equations of $`S=L`$ for the preferred coordinates $`X^\mu `$:
$$\frac{\delta S}{\delta X^\mu }(4\pi G)^1\gamma _{\mu \nu }\mathrm{}X^\nu $$
We have introduced here a constant diagonal matrix $`\gamma _{\mu \nu }`$ and a common factor $`(4\pi G)^1`$ to obtain appropriate units. With Euclidean symmetry we obtain $`\gamma _{11}=\gamma _{22}=\gamma _{33}`$. Thus, we have two coefficients $`\gamma _{00}=\mathrm{{\rm Y}},\gamma _{ii}=\mathrm{\Xi }`$. Now, the Lagrangian
$$L_0=(8\pi G)^1\gamma _{\mu \nu }X_{,\alpha }^\mu X_{,\beta }^\nu g^{\alpha \beta }\sqrt{g}$$
fulfils this property. For the difference $`LL_0`$ we obtain
$$\frac{\delta (LL_0)}{\delta X^\mu }0$$
Thus, the remaining part is not only covariant in the weak sense, but does not depend on the preferred coordinates $`X^\mu `$. But this is “strong” covariance, the classical requirement for the Lagrangian of general relativity. Thus, we can identify the difference with the classical Lagrangian of general relativity and obtain in the preferred coordinates
$$L=(8\pi G)^1\gamma _{\mu \nu }g^{\mu \nu }\sqrt{g}+L_{GR}(g_{\mu \nu })+L_{matter}(g_{\mu \nu },\phi ^m)$$
with the following modification of the Einstein equations
$$G_\nu ^\mu =8\pi G(T_m)_\nu ^\mu +(\mathrm{\Lambda }+\gamma _{\kappa \lambda }g^{\kappa \lambda })\delta _\nu ^\mu 2g^{\mu \kappa }\gamma _{\kappa \nu }$$
The last term is the full energy-momentum tensor, therefore, this equations defines a decomposition of the energy-momentum tensor into the energy-momentum tensor of matter and the energy-momentum tensor of the gravitational field defined by
$$(T_g)_\nu ^\mu =(8\pi G)^1\left(\delta _\nu ^\mu (\mathrm{\Lambda }+\gamma _{\kappa \lambda }g^{\kappa \lambda })G_\nu ^\mu \right)\sqrt{g}$$
## 3 Quantization
Most workers would agree that “at the root of most of the conceptual problems of quantum gravity” is the idea that “a theory of quantum gravity must have something to say about the quantum nature of space and time” . These problems, especially the problem of time , simply disappear in a theory with fixed Newtonian background. Problems related with energy and momentum conservation too – the Hamiltonian is no longer a constraint.
The violation of Bell’s inequality is independent evidence for a preferred frame. A preferred frame is required for compatibility with the EPR criterion of reality and Bohmian mechanics . Bell himself concludes : “the cheapest resolution is something like going back to relativity as it was before Einstein, when people like Lorentz and Poincare thought that there was an aether — a preferred frame of reference — but that our measuring instruments were distorted by motion in such a way that we could no detect motion through the aether.”
Our theory is in ideal agreement with “the present educated view on the standard model, and of general relativity, … that these are leading terms in effective field theories” – an idea introduced by Sakharov . It seems natural to assume that our medium has an atomic structure. An interpretation of $`\rho `$ as the number of “atoms” per volume leads to an interesting prediction for the cutoff:
$$\rho (x)V_{cutoff}=1.$$
It is non-covariant. For the homogeneous universe, it seems to expand together with the universe. It differs from the usual expectation that the cutoff is the Planck length $`a_P10^{33}cm`$ (cf. , ).
## 4 Comparison with other theories of gravity
Because of the simplicity of the additional terms it is no wonder that they have been already considered. Two other theories have the same Lagrangian for appropriate signs of the cosmological constants: the “relativistic theory of gravity” proposed by Logunov et al. and classical GR with some additional scalar “dark matter” fields. Nonetheless, equations are not all. There are other physical important things which makes the theories different as physical theories, like global restrictions, boundary conditions, causality restrictions, quantization concepts which are closely related with the underlying “metaphysical” assumptions.
### 4.1 Comparison with RTG
The “relativistic theory of gravity” (RTG) proposed by Logunov et al. has Minkowski background metric $`\gamma _{\mu \nu }`$. The Lagrangian of RTG is
$$L=L_{Rosen}+L_{matter}(g_{\mu \nu },\psi ^m)m_g^2(\frac{1}{2}\gamma _{\mu \nu }g^{\mu \nu }\sqrt{g}\sqrt{g}\sqrt{\gamma })$$
which de facto coincides with our theory for $`\mathrm{\Lambda }=m_g^2<0`$, $`\mathrm{\Xi }=\gamma ^{11}m_g^2>0`$, $`\mathrm{{\rm Y}}=\gamma ^{00}m_g^2>0`$.
The metaphysical context of RTG is completely different. It is a special-relativistic theory, therefore incompatible with the EPR criterion of reality and Bohmian mechanics. Another difference is the causality condition: In RTG, only solutions where the light cone of $`g_{ij}`$ is inside the light cone of $`\gamma _{ij}`$ are allowed. A comparable but weaker condition exists in our theory too: $`T(x)`$ should be a time-like function, or, $`\rho (X,T)>0`$. Note also that our theory suggests a different way of quantization: the prediction for the cutoff length $`l_{cutoff}`$ is not Lorentz-covariant.
### 4.2 Comparison with GR plus dark matter
The Lagrangian is also equivalent to GR with some dark matter – four scalar fields $`X^\mu `$. In this theory they are no longer preferred coordinates, but simply fields. Such “clock fields” in GR have been considered by Kuchar . Usual energy conditions require $`\mathrm{\Xi }>0,\mathrm{{\rm Y}}<0`$.
This GR variant allows a lot of solutions where the fields $`X^\mu (x)`$ cannot be used as global coordinates, especially solutions with non-trivial topology. They may also violate the condition that $`X^0(x)=T(x)`$ is time-like. Such solutions are forbidden in our theory. On the other hand, the infinite “boundary values” of the “fields” $`X^\mu (x)`$ are unreasonable for matter fields in GR. Another difference is that in our theory the $`X^\mu `$ are fixed background coordinates and therefore should not be quantized, while the “fields” $`X^\mu (x)`$ should be quantized.
## 5 Predictions
Using small enough values $`\mathrm{\Xi },\mathrm{{\rm Y}}0`$ leads to GR equations. Therefore it is not problematic to fit observation. It is much more problematic to find a way to distinguish our theory from GR by observation.
### 5.1 A dark matter candidate
Let’s consider the influence of the new terms on the expansion of the universe. In our theory a homogeneous universe is flat. The the usual ansatz $`ds^2=d\tau ^2a^2(\tau )(dx^2+dy^2+dz^2)`$ gives
$`3(\dot{a}/a)^2`$ $`=`$ $`\mathrm{{\rm Y}}/a^6+3\mathrm{\Xi }/a^2+\mathrm{\Lambda }+\epsilon `$
$`2(\ddot{a}/a)+(\dot{a}/a)^2`$ $`=`$ $`+\mathrm{{\rm Y}}/a^6+\mathrm{\Xi }/a^2+\mathrm{\Lambda }p`$
We see that $`\mathrm{\Xi }`$ influences the expansion of the universe similar to dark matter with $`p=\frac{1}{3}\epsilon `$.
### 5.2 Big bounce instead of big bang singularity
$`\mathrm{{\rm Y}}`$ becomes important only in the very early universe. But for $`\mathrm{{\rm Y}}>0`$, we obtain a qualitatively different picture. We obtain a lower bound $`a_0`$ for $`a(\tau )`$ defined by
$$\mathrm{{\rm Y}}/a_0^6=3\mathrm{\Xi }/a_0^2+\mathrm{\Lambda }+\epsilon $$
The solution becomes symmetrical in time, with a big crash followed by a big bang. For example, if $`\epsilon =\mathrm{\Xi }=0,\mathrm{{\rm Y}}>0,\mathrm{\Lambda }>0`$ we have the solution
$$a(\tau )=a_0\mathrm{cosh}^{1/3}(\sqrt{3\mathrm{\Lambda }}\tau )$$
In time-symmetrical solutions of this type the horizon is, if not infinite, at least big enough to solve the cosmological horizon problem (cf. ) without inflation.
### 5.3 Frozen stars instead of black holes
The choice $`\mathrm{{\rm Y}}>0`$ influences also another physically interesting solution – the gravitational collapse. There are stable “frozen star” solutions with radius slightly greater than their Schwarzschild radius. The collapse does not lead to horizon formation, but to a bounce from the Schwarzschild radius. Let’s consider an example. The general stable spherically symmetric harmonic metric depends on one step of freedom m(r) and has the form
$$ds^2=(1\frac{m}{r}\frac{m}{r})(\frac{rm}{r+m}dt^2\frac{r+m}{rm}dr^2)(r+m)^2d\mathrm{\Omega }^2$$
Let’s consider the ansatz $`m(r)=(1\mathrm{\Delta })r`$. We obtain
$`ds^2`$ $`=`$ $`\mathrm{\Delta }^2dt^2(2\mathrm{\Delta })^2(dr^2+r^2d\mathrm{\Omega }^2)`$
$`0`$ $`=`$ $`\mathrm{{\rm Y}}\mathrm{\Delta }^2+3\mathrm{\Xi }(2\mathrm{\Delta })^2+\mathrm{\Lambda }+\epsilon `$
$`0`$ $`=`$ $`+\mathrm{{\rm Y}}\mathrm{\Delta }^2+\mathrm{\Xi }(2\mathrm{\Delta })^2+\mathrm{\Lambda }p`$
Now, for very small $`\mathrm{\Delta }`$ even a very small $`\mathrm{{\rm Y}}`$ becomes important, and we obtain a non-trivial stable solution for $`p=\epsilon =\mathrm{{\rm Y}}g^{00}`$. Thus, the surface remains visible, with time dilation $`\sqrt{\epsilon /\mathrm{{\rm Y}}}M^1`$. |
warning/0001/nlin0001022.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The Painlevé analysis, as a test for integrability of partial differential equations (PDEs), was proposed by Weiss, Tabor and Carnevale in 1983 . It is a generalization of the singular point analysis for ordinary differential equations (ODEs), which dates back to the work of Sofia Kovalevskaya of 1889 . She studied the Euler-Poisson equations in the complex domain and found conditions under which the only movable singularities exhibited by the solutions were ordinary poles, leading to her discovery of a new first integral. In the late ninteenth century Paul Painlevé completely classified first order ODEs , as well as a large class of second order ODEs , on the basis that the only movable singularities their solutions exibit, are ordinary poles. This special property is today known as the the Painlevé property (see, for example ). We also say that an ODE is of Painlevé type, by which we mean that it belongs to the class of equations in Painlevé’s classification, or that it can be transformed to one of the equations in that class; therefore an ODE which has the Painlevé property. The list of ODEs, classified by Painlevé, is given in the book of Davis .
We consider a PDE to be integrable if it can be solved by an inverse scattering transform (we refer to the book , and references theirin). A PDE which is integrable possess the Painlevé property, which means that its solutions are single-valued in the neighbourhood of non-characteristic movable singularity manifolds . In this sense the method described by Weiss, Tabor and Carnevale proposes a necessary condition of integrability, also known as the Painlevé test, which is analogous to the algorithm for ODEs described by Ablowitz, Ramani and Segur which determines whether a given ODE has the Painlevé property. One seeks a solution of a given PDE (in rational form) in the form of a Laurent series (also known as the Painlevé expansion)
$`u(𝒙)=\varphi ^m(𝒙){\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}u_j(𝒙)\varphi ^j(𝒙),`$ (1.1)
where $`u_j(𝒙)`$ are analytic functions of the complex variables $`𝒙=(x_0,x_1,\mathrm{},x_{n1})`$ (we do not change notation for the complex domain), with $`u_00`$, in the neighbourhood of a non-characteristic movable singularity manifold defined by $`\varphi (𝒙)=0`$ (the pole manifold), where $`\varphi `$ is an analytic function of $`𝒙`$. The PDE is said to pass the Painlevé test if, on substituting (1.1) in the PDE, one obtains the correct number of arbitrary functions as required by the Cauchy-Kovalevsky theorem, given by the expansion coefficients in (1.1), whereby $`\varphi `$ should be one of the arbitrary functions. The coefficient in the Painlevé expansion, where the arbitrary functions are to appear, are known as the resonances. If a PDE satisfies the Painlevé test, it is usually possible to construct Bäcklund transformations and Lax pairs , which then proves the sufficient condition of integrability.
Recently some attention was given to the construction of exact solutions of nonintegrable PDEs by the use of a truncated Painlevé series . On applying the Painlevé expansion to nonintegrable PDEs one obtains conditions on $`\varphi `$ at the resonances; the singular manifold conditions. By truncating the series one usually obtains additional constraints on the singularity manifolds, leading to compatibility problems for the solution of $`\varphi `$ . It has been known for some time that the 2-dimensional Bateman equation
$`\varphi _{x_0x_0}\varphi _{x_1}^2+\varphi _{x_1x_1}\varphi _{x_0}^22\varphi _{x_0}\varphi _{x_1}\varphi _{x_0x_1}=0,`$ (1.2)
plays an important role in the Painlevé analysis of 2-dimensional nonintegrable PDEs .
In the present paper we show that the general solution of the $`n`$-dimensional Bateman equation, as generalized by Fairlie , solves the singularity manifold condition at the resonance for a class of wave equations. In the present paper we consider the $`n`$-dimensional ($`n3`$) sine-Gordon -, Liouville -, Mikhailov equation, and double sine-Gordon equation. The Painlevé test of the 2-dimensional double sine-Gordon equation was analyzed by Weiss , and resulted in the singularity constrained (1.2). Weiss pointed out that the 2-dimensional Bateman equation (1.2) can be linearized by a Legendre transformation. Moreover, it is invariant under the Moebius group and admits the general implicit solution
$`x_0f_0(\varphi )+x_1f_1(\varphi )=c,`$ (1.3)
where $`f_0`$ and $`f_1`$ are arbitrary smooth functions and $`c`$ is an arbitrary real constant. In the following section we derive the explicit relation between the singularity manifold and the 2-dimensional Bateman equation for two 2-dimensional polynomial wave equations. Finally we give some examples which demonstrate the use of our Propositions.
## 2 Propositions
Fairlie proposed the following $`n`$-dimensions Bateman equation:
$`det\left(\begin{array}{ccccc}0& \varphi _{x_0}& \varphi _{x_1}& \mathrm{}& \varphi _{x_{n1}}\\ \varphi _{x_0}& \varphi _{x_0x_0}& \varphi _{x_0x_1}& \mathrm{}& \varphi _{x_0x_{n1}}\\ \varphi _{x_1}& \varphi _{x_0x_1}& \varphi _{x_1x_1}& \mathrm{}& \varphi _{x_1x_{n1}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \varphi _{x_{n1}}& \varphi _{x_0x_{n1}}& \varphi _{x_1x_{n1}}& \mathrm{}& \varphi _{x_{n1}x_{n1}}\end{array}\right)=0.`$ (2.6)
Equation (2.6) generalizes the 2-dimensional Bateman equations (1.2) in $`n`$ dimensions. It admits the following general implicit solution
$`{\displaystyle \underset{j=0}{\overset{n1}{}}}x_jf_j(\varphi )=c,`$ (2.7)
where $`f_j`$ are $`n`$ arbitrary smooth functions.
We consider the $`n`$-dimensional generalization of the well known 2-dimensional sine-Gordon -, Liouville -, and Mikhailov equations, given respectively by
$`\mathrm{}_nu+\mathrm{sin}u=0`$
$`\mathrm{}_nu+\mathrm{exp}(u)=0`$ (2.8)
$`\mathrm{}_nu+\mathrm{exp}(u)+\mathrm{exp}(2u)=0,`$
as well as the double sine-Gordon equation in $`n`$ dimensions:
$`\mathrm{}_nu+\mathrm{sin}{\displaystyle \frac{u}{2}}+\mathrm{sin}u=0.`$ (2.9)
Here $`\mathrm{}_n`$ denotes the d’Alembert operator in $`n`$-dimensional Minkowski space, and is defined by
$`\mathrm{}_n:={\displaystyle \frac{^2}{x_0^2}}{\displaystyle \underset{j=1}{\overset{n1}{}}}{\displaystyle \frac{^2}{x_j^2}}.`$
It is well known that the wave equations (2) are integrable for $`n=2`$ (see, for example, ).
Before we state our Proposition for the singularity manifolds of those equations, we introduce some notations and a Lemma. We call the $`(n+1)\times (n+1)`$-matrix, the determinant of which defines the $`n`$-dimensional Bateman equation (2.6), the $`n`$-dimensional Bateman matrix and denote this matrix by $`B_{n+1}^n`$. The subscript of $`B`$ shows the size of the matrix while the superscript gives the dimension (the number of variables of $`\varphi `$), i.e., for the $`n`$-dimensional Bateman matrix (2.6), the associated Bateman matrix is
$`B_{n+1}^n=\left(\begin{array}{ccccc}0& \varphi _{x_0}& \varphi _{x_1}& \mathrm{}& \varphi _{x_{n1}}\\ \varphi _{x_0}& \varphi _{x_0x_0}& \varphi _{x_0x_1}& \mathrm{}& \varphi _{x_0x_{n1}}\\ \varphi _{x_1}& \varphi _{x_0x_1}& \varphi _{x_1x_1}& \mathrm{}& \varphi _{x_1x_{n1}}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ \varphi _{x_{n1}}& \varphi _{x_0x_{n1}}& \varphi _{x_1x_{n1}}& \mathrm{}& \varphi _{x_{n1}x_{n1}}\end{array}\right).`$ (2.15)
In particular the submatrices of the above $`n`$-dimensional Bateman matrix are of importance, i.e., the submatrices $`B_p^n`$, where $`3pn+1`$. These submatrices, which we call $`n`$-dimensional Bateman submatrices, are obtained by deleting rows and corresponding columns of $`B_{n+1}^n`$. We give the following
DEFINITION. Let
$`M_{x_{j_1}x_{j_2}\mathrm{}x_{j_r}}`$
denote the determinant of a Bateman submatrix, that remains after the rows and columns containing the derivatives $`\varphi _{x_{j_1}}`$, $`\varphi _{x_{j_2}}`$, …, $`\varphi _{x_{j_r}}`$ have been deleted from the $`n`$-dimensional Bateman matrix (2.15). Let
$`j_1,\mathrm{},j_r\{0,1,\mathrm{},n1\},j_1<j_2<\mathrm{}<j_r,rn2,\text{for}n3.`$
Then $`M_{x_{j_1}x_{j_2}\mathrm{}x_{j_r}}`$ are the determinants of the Bateman matrices $`B_{n+1r}^n`$. We call the equations
$`M_{x_{j_1}x_{j_2}\mathrm{}x_{j_r}}=0`$ (2.16)
the minor $`n`$-dimensional Bateman equations.
Note that the $`n`$-dimensional Bateman equation (2.6) has $`n!/[r!(nr)!]`$ minor $`n`$-dimensional Bateman equations. Consider an example: If $`n=5`$ and $`r=2`$, then there exist 10 minor Bateman equations, one of which is given by $`M_{x_2x_3}`$, i.e.,
$`det\left(\begin{array}{cccc}0& \varphi _{x_0}& \varphi _{x_1}& \varphi _{x_4}\\ \varphi _{x_0}& \varphi _{x_0x_0}& \varphi _{x_0x_1}& \varphi _{x_0x_4}\\ \varphi _{x_1}& \varphi _{x_0x_1}& \varphi _{x_1x_1}& \varphi _{x_1x_4}\\ \varphi _{x_4}& \varphi _{x_0x_4}& \varphi _{x_1x_4}& \varphi _{x_4x_4}\end{array}\right)=0.`$ (2.21)
We can now state the following
LEMMA. If $`\varphi `$ satisfies the $`n`$-dimensional Bateman equation (2.6), then it satisfies any minor Bateman equation
$`M_{x_{j_1}x_{j_2}\mathrm{}x_{j_r}}=0`$
with
$`j_1,\mathrm{},j_r\{0,1,\mathrm{},n1\},j_1<j_2<\mathrm{}<j_r,rn2,\text{for}n3.`$
Proof: By implicitly differentiating the general solution (2.7) of the $`n`$-dimensional Bateman equation (2.6), it is easily shown that any minor $`n`$-dimensional Bateman equation is satisfies by this solution. Since (2.7) is the general solution of the $`n`$-dimensional Bateman equation, the proof is concluded. $`\mathrm{}`$
We now prove
PROPOSITION 1. For $`n3`$, the singularity manifold conditions of the $`n`$-dimensional sine-Gordon -, Liouville -, and Mikhailov equations (2), are satisfied by the solution of the $`n`$-dimensional Bateman equation (2.6).
Proof: We do the proof for the sine-Gordon equation. For the Liouville - and Mikhailov equation, the proofs are similar. By the substitution
$`v(𝒙)=\mathrm{exp}[iu(𝒙)]`$
the $`n`$-dimensional sine-Gordon equation takes the following form:
$`v\mathrm{}_nv\left(_nv\right)^2+{\displaystyle \frac{1}{2}}\left(v^3v\right)=0,`$ (2.22)
where
$`\left(_nv\right)^2:=\left({\displaystyle \frac{v}{x_0}}\right)^2{\displaystyle \underset{j=1}{\overset{n1}{}}}\left({\displaystyle \frac{v}{x_j}}\right)^2.`$
The dominant behaviour of (2.22) is 2, so that the Painlevé expansion is
$`v(𝒙)={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}v_j(𝒙)\varphi ^{j2}(𝒙).`$
The resonance is at 2 and the first two coefficients in the Painlevé expansion have the following form:
$`v_0=4\left(_n\varphi \right)^2,v_1=4\mathrm{}_n\varphi .`$
We first consider $`n=3`$. The singularity manifold condition at the resonance is then given by
$`det\left(\begin{array}{cccc}0& \varphi _{x_0}& \varphi _{x_1}& \varphi _{x_2}\\ \varphi _{x_0}& \varphi _{x_0x_0}& \varphi _{x_0x_1}& \varphi _{x_0x_2}\\ \varphi _{x_1}& \varphi _{x_0x_1}& \varphi _{x_1x_1}& \varphi _{x_1x_2}\\ \varphi _{x_2}& \varphi _{x_0x_2}& \varphi _{x_1x_2}& \varphi _{x_2x_2}\end{array}\right)=0,`$
which is the 3-dimensional Bateman equation $`detB_4^3=0`$, as defined by (2.6).
Consider now $`n4`$. The condition at the resonance can be written as follows:
$`{\displaystyle \underset{j_1,j_2,\mathrm{},j_{n3}=1}{\overset{n1}{}}}M_{x_{j_1}x_{j_2}\mathrm{}x_{j_{n3}}}{\displaystyle \underset{j_1,j_2,\mathrm{},j_{n4}=1}{\overset{n1}{}}}M_{x_0x_{j_1}x_{j_2}\mathrm{}x_{j_{n4}}}=0,`$ (2.24)
where
$`j_1<j_2<\mathrm{}<j_{n3},`$
and $`M_{x_{j_1}x_{j_2}\mathrm{}x_{j_{n3}}}`$, $`M_{x_0x_{j_1}x_{j_2}\mathrm{}x_{j_{n4}}}`$ are minor $`n`$-dimensional Bateman equations, i.e., the determinants of $`4\times 4`$ Bateman matrices $`B_4^n`$. By the Lemma give above, equation (2.24) is satisfied by the solution of the $`n`$-dimensional Bateman equation (2.6). $`\mathrm{}`$
We now consider the double sine-Gordon equation in $`n`$ dimensions (2.9):
$`\mathrm{}_nu+\mathrm{sin}{\displaystyle \frac{u}{2}}+\mathrm{sin}u=0.`$
It was shown by Weiss , that for $`n=2`$ this equation does not pass the Painlevé test, and that the singularity manifold condition is given by the Bateman equation (1.2).
For $`n`$ dimensions we prove the following
PROPOSITION 2. For $`n2`$, the singularity manifold condition of the double sine-Gordon equation (2.9) is satisfied by the solution of the $`n`$-dimensional Bateman equation (2.6).
Proof: By the substitution
$`v(𝒙)=\mathrm{exp}\left[{\displaystyle \frac{i}{2}}u(𝒙)\right]`$
the rational form of the double sine-Gordon equation (2.9) is obtained as
$`v\mathrm{}v+(v_n)^2+{\displaystyle \frac{1}{4}}(v^3v)+{\displaystyle \frac{1}{4}}(v^41)=0.`$
The Painlevé expansion takes the form
$`v(𝒙)={\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}v_j(𝒙)\varphi ^{j1}(𝒙)`$
and the resonance is 2. The first two expansion coefficients are
$`v_0=4(_nv)^2,v_1={\displaystyle \frac{2}{v_0}}\mathrm{}_n\varphi {\displaystyle \frac{1}{2}}`$
For the singularity manifold condition we have to consider four cases:
Case $`n=2`$: At the resonance we obtain (1.2), i.e.,
$`detB_3^2=0.`$
Case $`n=3`$: The condition now takes the following form:
$`8detB_4^3+\left(M_{x_1}+M_{x_2}M_{x_0}\right)v_0=0.`$
Case $`n4`$: The condition at the resonance can be written as follows:
$`8\left({\displaystyle \underset{j_1,j_2,\mathrm{}j_{n3}=1}{\overset{n1}{}}}M_{x_{j_1}x_{j_2}\mathrm{}x_{j_{n3}}}\right)`$
$`+\left({\displaystyle \underset{j_1,j_2,\mathrm{}j_{n2}=1}{\overset{n1}{}}}M_{x_{j_1}x_{j_2}\mathrm{}x_{j_{n2}}}{\displaystyle \underset{j_1,j_2,\mathrm{}j_{n3}=1}{\overset{n1}{}}}M_{x_0x_{j_1}x_{j_2}\mathrm{}x_{j_{n3}}}\right)v_0=0,`$
where
$`j_1<j_2<\mathrm{}<j_{n3}<j_{n2}.`$
By the above Lemma the proof is concluded. $`\mathrm{}`$
We now consider two well known nonlinear polynomial field theory equations, the so-called nonlinear Klein-Gordon equations:
$`\mathrm{}_2u+u^k=0`$ (2.25)
with $`k=2,3`$. In light-cone coordinates, i.e.,
$`x_0{\displaystyle \frac{1}{2}}(x_0x_1),x_1{\displaystyle \frac{1}{2}}(x_0+x_1),`$
(2.25) takes the form
$`{\displaystyle \frac{^2u}{x_0x_1}}+u^k=0.`$ (2.26)
It should be noted that the 2-dimensional Bateman equation remains invariant under the light-cone coordinates. Therefore, for our purpose we can work with (2.26) instead of (2.25). In it was shown that the nonlinear Klein-Gordon equation (2.26), with $`k=3`$, does not pass the Painlevé test. We are now interested in the relation between the 2-dimensional Bateman equation (1.2) and the singularity manifold condition of (2.26) for the case $`k=2`$ as well as $`k=3`$.
We prove the following
PROPOSITION 3. The solution of the 2-dimensional Bateman equation (1.2) satisfies the singularity manifold condition of the nonlinear Klein-Gordon equation (2.26) for $`k=2`$ and $`k=3`$.
Proof: First we consider equation (2.26) with $`k=3`$, i.e.,
$`{\displaystyle \frac{^2u}{x_0x_1}}+u^3=0.`$ (2.27)
For the Painlevé expansion
$`u(x_0,x_1)=\varphi ^m(x_0,x_1){\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}u_j(x_0,x_1)\varphi ^j(x_0,x_1),`$ (2.28)
we find that the dominant behaviour is -1, the resonance is 4, and the first three expansion coefficients in expansion (2.28) are
$`u_0^2=2\varphi _{x_0}\varphi _{x_1},`$
$`u_1={\displaystyle \frac{1}{3u_0^2}}\left(u_0\varphi _{x_0x_1}+u_{0x_1}\varphi _{x_0}+u_{0x_0}\varphi _{x_1}\right),`$
$`u_2={\displaystyle \frac{1}{3u_0^2}}\left(u_{0x_0x_1}3u_0u_1^2\right),`$
$`u_3={\displaystyle \frac{1}{u_0^2}}\left(u_2\varphi _{x_0x_1}+u_{2x_1}\varphi _{x_0}+u_{2x_0}\varphi _{x_1}+u_{1x_0x_1}6u_0u_1u_2\right).`$
At the resonance we obtain the following singularity manifold condition:
$`\mathrm{\Phi }\sigma \left(\varphi _{x_0}\mathrm{\Phi }_{x_1}\varphi _{x_1}\mathrm{\Phi }_{x_0}\right)^2=0,`$ (2.29)
where $`\mathrm{\Phi }`$ is the 2-dimensional Bateman equation given by (1.2), i.e.,
$`\mathrm{\Phi }=\varphi _{x_0x_0}\varphi _{x_1}^2+\varphi _{x_1x_1}\varphi _{x_0}^22\varphi _{x_0}\varphi _{x_1}\varphi _{x_0x_1},`$
and $`\sigma `$ contains derivatives of $`\varphi `$ with respect to $`x_0`$ and $`x_1`$. The explicit form of $`\sigma `$ is not interesting for our proof. The explicit appearance of $`\mathrm{\Phi }`$ (2.29) concludes the proof for the nonlinearity $`k=3`$.
For the equation
$`{\displaystyle \frac{^2u}{x_0x_1}}+u^2=0`$ (2.30)
the singularity manifold condition is somewhat more complicated. The dominant behaviour of (2.30) is -2 and the resonance is 6. The first five expansion coefficients in the Painlevé expansion take the following form:
$`u_0=6\varphi _{x_0}\varphi _{x_1},`$
$`u_1={\displaystyle \frac{1}{\varphi _{x_0}\varphi _{x_1}+u_0}}\left(u_{0x_1}\varphi _{x_0}+u_{0x_0}\varphi _{x_1}+u_0\varphi _{x_0x_1}\right),`$
$`u_2={\displaystyle \frac{1}{2u_0}}\left(u_{0x_0x_1}+u_1^2u_{1x_1}\varphi _{x_0}u_{1x_0}\varphi _{x_1}u_1\varphi _{x_0x_1}\right),`$
$`u_3={\displaystyle \frac{1}{2u_0}}\left(u_{1x_0x_1}+2u_1u_2\right),`$
$`u_4={\displaystyle \frac{1}{\varphi _{x_1}\varphi _{x_0}+u_0}}\left(u_3\varphi _{x_0x_1}+u_{2x_0x_1}+2u_1u_3+u_{3x_1}\varphi _{x_0}+u_{3x_0}\varphi _{x_1}+u_2^2\right),`$
$`u_5={\displaystyle \frac{1}{6\varphi _{x_0}\varphi _{x_1}+2u_0}}\left(2u_1u_4+2u_4\varphi _{x_0x_1}+2u_{4x_0}\varphi _{x_1}+2u_{4x_1}\varphi _{x_0}+2u_2u_3+u_{3x_0x_1}\right).`$
At the resonance the singularity manifold condition is a PDE of order six, which consists of 372 terms (!) all of which are derivatives of $`\varphi `$ with respect to $`x_0`$ and $`x_1`$. This condition may be written in the following form:
$`\sigma _1\mathrm{\Phi }+\sigma _2\mathrm{\Psi }+\left(\varphi _{x_0}\mathrm{\Psi }_{x_1}\varphi _{x_1}\mathrm{\Psi }_{x_0}\sigma _3\mathrm{\Psi }\sigma _4\mathrm{\Phi }\right)^2=0,`$ (2.31)
where $`\mathrm{\Phi }`$ is the 2-dimensional Bateman equation (1.2), and
$`\mathrm{\Psi }=\varphi _{x_0}\mathrm{\Phi }_{x_1}\varphi _{x_1}\mathrm{\Phi }_{x_0},\mathrm{\Phi }=\varphi _{x_0x_0}\varphi _{x_1}^2+\varphi _{x_1x_1}\varphi _{x_0}^22\varphi _{x_0}\varphi _{x_1}\varphi _{x_0x_1}.`$
Here $`\sigma _1,\mathrm{},\sigma _4`$ consist of derivatives of $`\varphi `$ with respect to $`x_0`$ and $`x_1`$. Their explicit form is not interesting. By (2.31) it is clear that the general solution of the Bateman equation satisfies the singularity manifold condition for (2.30). $`\mathrm{}`$
Due to its enormous complexity in higher dimensions, we were not able to find the explicit relations between the singularity manifold for higher dimensional equations of the form
$`\mathrm{}_nu+u^k=0`$ (2.32)
and the $`n`$-dimensional Bateman equation (or minor Bateman equations). We
CONJECTURE. In $`n`$-dimensions, the solution of the $`n`$-dimensional Bateman equation (2.6) satisfies the singularity manifold condition of (2.32) for $`k=2,3`$.
Some examples of (2.32) are also given below, and these are consistent with this view.
## 3 Application
According to a conjecture by Ablowitz, Ramani and Segur , every ODE that can be obtained by a Lie symmetry reduction (similarity reduction) of a PDE, which is solvable by the inverse scattering transform method, has the Painlevé property. Some weak form of this conjecture was proved in . On the other hand, if we would consider a nonintegrable 2-dimensional PDE, then it is possible that some of the ODEs resulting by some reduction Ansatz of the PDE, may also be of Painlevé type. In particular, the reduced ODE would fullfil the necessary condition to be of Painlevè type (pass the Painlevé test for ODEs) for those Ansätze for which the new independent variable satisfies the condition on the singularity manifold of the given PDE. By the Propositions stated in the previous section, we know that the condition on the singularity manifold is satisfied by the $`n`$-dimensional Bateman equation for our class of equations. Thus, the Propositions, lead to the following
COROLLARY. The nonlinear wave equations (2), (2.9), (2.27) (2.30) can be reduced to ODEs which satisfy the necessary condition to be of Painlevé type, if and only if the new independent variables of the reduced ODEs satisfy the corresponding $`n`$-dimensional Bateman equation (1.2).
This means that if we were to reduce one of the nonintegrable $`n`$-dimensional PDEs discussed in our paper into an ODE with independent variable $`\omega `$ by, for example, an Ansatz of the form
$`u(x_0,x_1,\mathrm{},x_{n1})=f_1(x_0,x_1,\mathrm{},x_{n1})\phi (\omega )+f_1(x_0,x_1,\mathrm{},x_{n1}),`$ (3.1)
then we can easily test the necessary condition of integrability of the resulting ODE by checking whether $`\omega `$ satisfies the $`n`$-dimensional Bateman equation (2.6). This would be the same as to perform the Painlevé test on the resulting ODE. By Lie symmetry analysis of PDEs one is able to systematically construct Ansätze which reduce the PDEs to ODEs according to their Lie transformation group properties (see for example ). By the above Corollary one is now able to classify the group invariants (that are independent of $`u`$) for the given PDEs, and determine which group invariants may result in ODE’s of Painlevé type, whithout performing the Painlevé analysis on the actual reduced ODEs, but by merely checking whether the invariants satisfy the $`n`$-dimensional Bateman equation (2.6). One must note that the reduction Ansatz is not necessarily related to a classical Lie symmetry invariant. One can obtain very interesting reduction Ansätze by the use of the so-called conditional symmetries, or $`Q`$-symmetries (see for some interesting examples).
Below we give some examples of the stated Corollary. A more systematic analysis and classification of the the equations treated here, will be presented in a future paper.
EXAMPLE 1. Consider the 3-dimensional Liouville equation , i.e.,
$`\mathrm{}_3u+\lambda \mathrm{exp}(u)=0,`$ (3.2)
with the Ansatz
$`u(x_0,x_1,x_2)=\phi (\omega )2\mathrm{ln}(\alpha _0y_0\alpha _1y_1\alpha _2y_2)`$
$`\omega (x_0,x_1,x_2)=(\alpha _0y_0\alpha _1y_1\alpha _2y_2)(\beta _0y_0\beta _1y_1\beta _2y_2)^a`$ (3.3)
where $`a𝒬\backslash \{0\}`$ and
$`\alpha _0^2\alpha _1^2\alpha _2^2=\alpha _0\beta _0\alpha _1\beta _1\alpha _2\beta _2=0,`$
$`\beta _0\beta _0\beta _1\beta _1\beta _2\beta _2<0,`$
$`y_\mu =x_\mu +a_\mu ,\mu =0,1,2.`$
Here $`\omega `$, given by (3), satisfies the 3-dimensional Bateman equation $`detB_4^3=0`$, so that by the Corollary we are ensured that the reduced ODE, resulting from Ansatz (3), satisfies the necessary condition to be of Painlevé type. Ansatz (3) leads to the following ODE:
$`a^2\omega ^2{\displaystyle \frac{d^2\phi }{d\omega ^2}}+a(a1)\omega {\displaystyle \frac{d\phi }{d\omega }}+\lambda \mathrm{exp}(\phi )=0.`$ (3.4)
Equation (3.4) is of Painlevé type and admits the general solution
$`\phi (\omega )=2\mathrm{ln}\left[{\displaystyle \frac{\sqrt{\lambda }}{\sqrt{2}c_1}}\omega ^{1/a}\mathrm{cos}(c_1\omega ^{1/a}+c_2)\right];\lambda <0`$ (3.5)
$`\phi (\omega )=2\mathrm{ln}\left[{\displaystyle \frac{\sqrt{\lambda }}{\sqrt{2}c_1}}\omega ^{1/a}\mathrm{cosh}(c_1\omega ^{1/a}+c_2)\right];\lambda >0.`$ (3.6)
By (3.5) and the Ansatz (3) an exact solution of the Liouville equation (3.2) follows:
$`u(x_0,x_1,x_2)=2\mathrm{ln}\left[{\displaystyle \frac{\sqrt{\lambda }}{\sqrt{2}c_1}}\omega ^{1/a}\mathrm{cos}(c_1\omega ^{1/a}+c_2)\right]2\mathrm{ln}(\alpha _0y_0\alpha _1y_1\alpha _2y_2);\lambda <0`$
$`u(x_0,x_1,x_2)=2\mathrm{ln}\left[{\displaystyle \frac{\sqrt{\lambda }}{\sqrt{2}c_1}}\omega ^{1/a}\mathrm{cosh}(c_1\omega ^{1/a}+c_2)\right]2\mathrm{ln}(\alpha _0y_0\alpha _1y_1\alpha _2y_2);\lambda >0`$
$`\omega (x_0,x_1,x_2)=(\alpha _0y_0\alpha _1y_1\alpha _2y_2)(\beta _0y_0\beta _1y_1\beta _2y_2)^a,y_\mu =x_\mu +a_\mu ,\mu =0,1,2.`$
This example can easily be extended to $`n`$ dimensions.
EXAMPLE 2. Consider the 4-dimensional sine-Gordon equation , i.e,
$`\mathrm{}_4u+\mathrm{sin}(u)=0.`$ (3.7)
By the Ansatz
$`u(x_0,x_1,x_2,x_3)=\phi (\omega )`$
$`\omega (x_0,x_1,x_2,x_3)={\displaystyle \frac{x_2x_3(x_0+x_1)}{\sqrt{1+(x_0+x_1)^2}}}+f(x_0+x_1),`$ (3.8)
where $`f`$ is an arbitrary smooth function of its argument, (3.7) reduces to the following integrable ODE:
$`{\displaystyle \frac{d^2\phi }{d\omega ^2}}\mathrm{sin}\phi =0.`$ (3.9)
It easy to show that $`\omega `$, given by (3), satisfies the 4-dimensional Bateman equation $`detB_5^4=0`$. Equation (3.9) can be integrated in terms of Jacobi elliptic functions to obtain exact solutions of the 4-dimensional sine-Gordon equation (3.7).
EXAMPLE 3. Consider the 2-dimensional nonlinear Klein-Gordon equation
$`u_{x_0x_1}+\lambda u^3=0.`$ (3.10)
We demonstrate that by the given Corollary and the Ansatz
$`u(x_0,x_1)=h(x_0,x_1)\phi (\omega ),`$ (3.11)
where $`\omega `$ satisfies the 2-dimensional Bateman equation (1.2) i.e.,
$`x_0f_0(\omega )+x_1f_1(\omega )=c,`$
we are able to construct ODEs which pass the Painlevé test. Ansatz (3.11) leads to
$`\left({\displaystyle \frac{fgh}{(x_0\dot{f}_0+x_1\dot{f}_1)^2}}\right){\displaystyle \frac{d^2\phi }{d\omega ^2}}`$
$`+\left({\displaystyle \frac{h(\dot{f}_0f_1+f_0\dot{f}_1)}{(x_0\dot{f}_0+x_1\dot{f}_1)^2}}{\displaystyle \frac{fgh(x_0\ddot{f}_0+x_1\ddot{f}_1)}{(x_0\dot{f}_0+x_1\dot{f}_1)^3}}{\displaystyle \frac{h_{x_1}f_0+h_{x_0}f_1}{(x_0\dot{f}_0+x_1\dot{f}_1)}}\right){\displaystyle \frac{d\phi }{d\omega }}`$
$`+h_{x_0x_1}\phi +\lambda h^3\phi ^3=0.`$ (3.12)
Here $`h=h(x_0,x_1),f_i=f_i(\omega )`$, and $`\dot{f}_idf_i/d\omega `$ $`(i=0,1)`$. For example, let
$`h(x_0,x_1)={\displaystyle \frac{1}{x_0}},f_1(\omega )=1,`$
then (3) reduces to
$`\ddot{\phi }+\left(2{\displaystyle \frac{\dot{f}}{f}}{\displaystyle \frac{\ddot{f}}{\dot{f}}}\right)\dot{\phi }\left({\displaystyle \frac{\lambda \dot{f}^2}{f}}\right)\phi ^3=0.`$ (3.13)
Equation (3.13) satisfies the necessary condition to be of Painlevé type (it passess the Painlevé test for ODEs), which is in agreement with the above Corollary, as we are using the general solution of the 2-dimensional Bateman equation (1.2). Note that for $`f_0(\omega )=\omega `$ we obtain the same ODE which was obtained with a Lie symmetry analysis in . We remark that the use of the general solution (1.3) of the Bateman equation (1.2), in the construction of exact solutions of (3.10), is clearly limited. A more effective approach, to obtain exact solutions, would be to linearize the 2-dimensional Bateman equation by the Legendre transformation, as outlined by Webb and Zank . However, this is not the purpose of the present paper.
EXAMPLE 4. Consider the 4-dimensional nonlinear Klein-Gordon equation
$`\mathrm{}_4u+\lambda u^3=0,`$ (3.14)
where $`\lambda `$. Assymptotic solutions of (3.14) were constructed in by the use the Poincaré group $`P(1,3)`$ and its invariants. By composing the group invariants, we obatin the following Ansatz for (3.14):
$`u(x_0,x_1,x_2,x_3)=\phi (\omega )`$
$`\omega (x_0,x_1,x_2,x_3)=\beta _1(<\stackrel{~}{𝒑},𝒙>+a_1)\beta _2(<\stackrel{~}{𝜶},𝒙>+a_2)\beta _3(<\stackrel{~}{𝜷},𝒙>+a_3)`$ (3.15)
$`+a\mathrm{ln}\{\alpha _1(<\stackrel{~}{𝒑},𝒙>+a_1)\alpha _2(<\stackrel{~}{𝜶},𝒙>+a_2)\alpha _3(<\stackrel{~}{𝜷},𝒙>+a_3)\}.`$
Here $`<\stackrel{~}{𝒑},𝒙>p_0x_0p_1x_1p_2x_2p_3x_3,<\stackrel{~}{𝜶},𝒙>\alpha _0x_0\alpha _1x_1\alpha _2x_2\alpha _3x_3,`$ $`<\stackrel{~}{𝜷},𝒙>\beta _0x_0\beta _1x_1\beta _2x_2\beta _3x_3`$ and $`a_j`$ ($`j=0,1,2,3`$) are arbitrary real constants, whereas $`\alpha _j,\beta _j,\stackrel{~}{\alpha }_\mu ,\stackrel{~}{\beta }_\mu ,\stackrel{~}{p}_\mu `$ ($`j=1,2,3;\mu =0,1,2,3`$) are real constants which must satisfy the following conditions:
$`\beta _1^2\beta _2^2\beta _3^2=1,\alpha _1^2\alpha _2^2\alpha _3^2=\alpha _1\beta _1\alpha _2\beta _2\alpha _3\beta _3=0`$ (3.16)
$`<\stackrel{~}{𝒑},\stackrel{~}{𝒑}>=1,<\stackrel{~}{𝜶},\stackrel{~}{𝜶}>=<\stackrel{~}{𝜷},\stackrel{~}{𝜷}>=1,`$
$`<\stackrel{~}{𝜶},\stackrel{~}{𝜷}>=<\stackrel{~}{𝜶},\stackrel{~}{𝒑}>=<\stackrel{~}{𝜷},\stackrel{~}{𝒑}>=0.`$ (3.17)
Here $`\omega `$, given by (3), satisfies the 4-dimensional Bateman equation $`detB_5^4=0`$, and the reduced equation
$`{\displaystyle \frac{d^2\phi }{d\omega ^2}}+\lambda \phi ^3=0`$ (3.18)
is of Painlevé type. The general solution of (3.18) is given in terms of Jacobi elliptic functions .
EXAMPLE 5. Consider the 4-dimensional nonlinear Klein-Gordon equation
$`\mathrm{}_4u+\lambda _1u+\lambda _2u^3=0,`$ (3.19)
where $`\lambda _1,\lambda _2`$. By the invariants of the Poincaré group, and its Lie subalgebras, the following two Ansätze are, for example, possible:
$`u(x_0,x_1,x_2,x_3)=\phi (\omega _1)`$
$`\omega _1={\displaystyle \frac{c}{2}}\{<\stackrel{~}{𝜸},𝒙>^2+[<\stackrel{~}{𝜷},𝒙>+{\displaystyle \frac{1}{4}}(<\stackrel{~}{𝒑},𝒙>+<\stackrel{~}{𝜶},𝒙>)^2]^{1/2}\}`$
$`+q_1<\stackrel{~}{𝜸},𝒙>q_2[<\stackrel{~}{𝜷},𝒙>+{\displaystyle \frac{1}{4}}(<\stackrel{~}{𝒑},𝒙>+<\stackrel{~}{𝜶},𝒙>)^2],`$ (3.20)
and
$`u(x_0,x_1,x_2,x_3)=\phi (\omega _2)`$
$`\omega _2(x_0,x_1,x_2,x_3)=q_3[<\stackrel{~}{𝒑},𝒙>^2<\stackrel{~}{𝜶},𝒙>^2<\stackrel{~}{𝜷},𝒙>^2]^{1/2},`$ (3.21)
where $`<\stackrel{~}{𝒑},𝒙>\stackrel{~}{p}_0x_0\stackrel{~}{p}_1x_1\stackrel{~}{p}_2x_2\stackrel{~}{p}_3x_3`$, $`<\stackrel{~}{𝜶},𝒙>\stackrel{~}{\alpha }_0x_0\stackrel{~}{\alpha }_1x_1\stackrel{~}{\alpha }_2x_2\stackrel{~}{\alpha }_3x_3`$, and $`<\stackrel{~}{𝜷},𝒙>\stackrel{~}{\beta }_0x_0\stackrel{~}{\beta }_1x_1\stackrel{~}{\beta }_2x_2\stackrel{~}{\beta }_3x_3`$. Here $`c`$ and $`q_3`$ are arbitrary nonzero real constants, whereas the rest of the real parameters have to satisfy condition (3) and
$`<\stackrel{~}{𝜸},\stackrel{~}{𝜸}>=1,<\stackrel{~}{𝜸},\stackrel{~}{𝒑}>=<\stackrel{~}{𝜷},\stackrel{~}{𝜸}>=<\stackrel{~}{𝜶},\stackrel{~}{𝜸}>=0,q_1^2+q_2^2=q0.`$
By the above Ansätze the following ODEs are respectively obtained:
$`(2c\omega _1+q){\displaystyle \frac{d^2\phi }{d\omega _1^2}}+2c{\displaystyle \frac{d\phi }{d\omega _1}}\lambda _1\phi +\lambda _2\phi ^3=0,`$ (3.22)
$`q_3\omega _2{\displaystyle \frac{d^2\phi }{d\omega _2^2}}+2q_3{\displaystyle \frac{d\phi }{d\omega _2}}+\lambda _1\omega _2\phi \lambda _2\omega _2\phi ^3=0.`$ (3.23)
Equations (3.22) and (3.23) are not of Painlevé type, which is in agreement with the fact that $`\omega _1`$ and $`\omega _2`$ do not satisfy the 4-dimensional Bateman equation $`detB_5^4=0`$.
A systematic classification of integrable reductions of the above given multidimensional wave equations, by the use of the Propositions and Corollary stated here, will be the subject of a future paper.
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warning/0001/hep-lat0001032.html | ar5iv | text | # Review of Unquenched Results
## 1 INTRODUCTION
A major challenge in realizing the full potential of the non-perturbative regularization provided by lattice field theory is the presence of fermion fields in simulations. In many cases, such as QCD at finite density and the Hubbard model, the fermions require an improvement in algorithms before large systems can be simulated. For QCD at zero chemical potential, simulations with dynamical fermions will continue to become better controlled by faster computers, even without (hoped for) theoretical improvements.
There is over a decade of experience with full QCD simulations, although only recently has enough data become available for some extrapolations to the continuum . As with quenched simulations questions of what are satisfactory volumes and lattice spacings, what are reliable simulation lengths and what are the best extrapolations need to be answered. Removing the uncontrolled truncation of the quenched approximation is vital to achieving precision results, but answers to the above are very important.
After a brief overview of the work on algorithms presented at this conference, I will focus on results for QCD, particularly the low lying hadrons. There is evidence the quenched hadron spectrum, in the continuum and chiral limits, differs from nature . For full QCD, the light hadron spectrum is interesting in its own right, provides the basis for a determination of light quark masses and is a useful testing ground for exploring finite volume effects and simulation lengths. Also, hadronic states differing only by their parity are degenerate unless chiral symmetry is broken, so studying these other parity states in the spectrum may show sensitivity to the inclusion of the quark determinant.
Most of the techniques developed in the quenched approximation for measuring other observables can be easily adapted to full QCD simulations. Many groups are reporting on effects of dynamical fermions outside of the light hadron spectrum ($`f_B`$, etc.). Please see the reviews of these areas for information.
## 2 ALGORITHMS
Three groups reported work on the algorithms used to produce the Markov chain in numerical simulations.
### 2.1 Meron Cluster Algorithm
Wiese and collaborators (MIT) , have tackled the long standing problem of algorithms for fermion systems where the weights in the path integral are complex. Computational times exponential in the volume are needed if the weight is made real by taking its absolute value and moving the phase to an observable. By working with a continuum Euclidean time formulation, where fermions are represented by their world lines, their algorithms form clusters whose flip gives a definite sign. A second improvement insures that clusters which cancel through a sign flip and those that don’t are generated with similar probabilities. This algorithm, with its attendant improved estimator, has allowed them to simulate various fermionic systems that are intractable with conventional approaches. Application of these ideas to coupled gauge-fermion systems is being investigated.
### 2.2 Truncated Determinant Approach
Duncan, et. al. have been studying an algorithm where the fermion determinant is split, in a gauge-invariant way, into an infrared and ultraviolet part. The ultraviolet part is modeled by an effective action made up of small Wilson loops, while the infrared part is determined precisely using a Lanczos procedure. To date, they have generated a sequence of gauge fields using just the infrared part of the determinant and have found that the exact contribution of the ultraviolet part is well fit by an expansion in a small number of Wilson loops. A somewhat surprising result also comes from the generation of the gauge fields: they first update links and then put in a global accept/reject step based on the infrared part of the determinant. They find reasonable acceptance for this, even for large volumes.
### 2.3 Multiboson Algorithm
de Forcrand and collaborators have been testing an evolved version of Luscher’s multiboson algorithm in a real QCD context and comparing its performance to a standard, but highly tuned, hybrid Monte-Carlo code of the SESAM collaboration. This new work uses only 24 boson fields for a simulation with $`m_\pi /m_\rho =0.833`$, a dramatic improvement. This is achieved by an “ultraviolet filtering” of the determinant, where the determinant is filtered by an exponential term made up of weighted small Wilson loops, and a cancelling term enters the gauge action. (This filtering is similar to that used by , but here is part of an exact algorithm.) An optimized quasi-heatbath is then applied to the combined boson-gauge system. They conclude that for a $`16^3\times 24`$ lattice, with $`\beta =5.6`$ and the equivalent to $`\kappa =0.156`$ for Wilson fermions, the multiboson algorithm decorrelates the plaquette better than HMC, measuring both algorithms in units of applications of the Dirac operator.
## 3 SIMULATION KEY
To distinguish the different actions, the abbreviations below will be used (following the CP-PACS collaboration).
| gauge action | |
| --- | --- |
| P | plaquette |
| R | RG improved |
| fermion action | |
| S | staggered |
| W | Wilson |
| C | Clover |
| D | domain wall |
The RG improved gauge action is the one proposed by Iwasaki, et. al. .
In the tables of simulations parameters, three results are generally listed: $`m_\rho a`$, $`m_\pi /m_\rho `$ and $`m_\rho L`$, where $`L`$ is the spatial lattice extent and $`a`$ is the lattice spacing. The value of $`m_\rho a`$ (where available) is for the quark mass which gives the quoted $`m_\pi /m_\rho `$. Also note that a 770 MeV $`\rho `$ in a 3 fm box has $`m_\rho L=11.6`$. Since these tables should only serve as a guide to the general features of the data sets, errors are not listed and only 2 significant figures are given. Please see the original work for more details.
The hybrid Monte Carlo (HMC) algorithm has been used for the 2 flavor Wilson and domain wall and 4 flavor staggered simulations. For 2 flavor staggered simulations the ‘R’ algorithm of Gottlieb, et. al. has been employed. The trajectory counts in the tables are from the groups themselves and the definition of a trajectory differs between the groups.
## 4 STAGGERED FERMIONS
Recent staggered simulations have been done by the MILC and Columbia groups and are listed in Table 1. MILC reported hadron spectrum results last year and new results for decay constants this year . They found a continuum extrapolation of $`m_N/m_\rho `$ for staggered fermions in full QCD gives a larger value than in the quenched continuum. The Columbia group has new data for 2 and 4 flavor QCD with staggered fermions, with simulations still underway .
There are existing results from 2 flavor staggered simulations at $`\beta =5.7`$ and $`m=0.01`$ which give some information about run lengths. Figure 1 shows, starting at the top, $`m_N`$, $`m_\rho `$ and $`m_\pi `$ from the Columbia group and Fukugita, et. al. . The masses labeled “1k” are from 1,000 trajectories on a $`20^4`$ lattice, “3k” from 3,000 trajectories on $`16^3\times 32`$ and “10k” from 10,000 trajectories on $`16^3\times 40`$. The points to the left of “10k” are the 10,000 trajectory run broken up in to nine 1,000 trajectory runs (plus thermalization) . The solid horizontal lines are the “10k” masses and the dashed horizontal lines give $`\pm 5`$% of the central value.
The similar size for the errors on the 3k and 10k runs (the same analysis was used for both) indicates that long correlation times were likely missed in the 3,000 trajectory run. The larger error bars for the 1,000 trajectory runs may be due to short term noise effects. For these masses and couplings, achieving reliable errors at the few percent level likely requires more than 10,000 trajectory run.
Another important question for full QCD simulations is the volume required. We now investigate this question through the splitting between parity partners in the hadron spectrum. We will see that, at least for staggered fermions, this is a sensitive indicator of finite volume effects.
MILC studies of finite volume effects in 2 flavor dynamical staggered simulations for couplings up to $`\beta =5.6`$ show lattice volumes of 2.5-3 fermi remove finite volume effects for $`m_\pi /m_\rho 0.5`$. Their most recent $`\beta =5.6`$ results for a $`24^4\times 48`$ lattice are shown in Figure 2. There is no sign of parity doubling, consistent with their statement about finite volume effects.
Fukugita, et. al. also studied finite volume effects but at weaker coupling ($`\beta =5.7`$). Their data is plotted in Figure 3, revealing parity doubling in the $`m0`$ limit for the $`N`$ and $`N^{}`$. (The $`N^{}`$ is the staggered fermion parity partner of $`N`$.) Finite volume effects are likely distorting the baryons, but not the mesons. Which way $`m_N`$ and $`m_N^{}`$ move for larger volume is an open question, currently being addressed by $`24^3`$ simulations underway at Columbia.
The Columbia group has previously reported parity doubling for 4 flavor QCD on a $`16^3\times 32`$ volume with $`\beta =5.4`$. (These parameters give a rho mass within a few percent of the rho mass for 2 flavors on a $`16^3\times 32`$ lattice.) Figure 4 shows our results, where the $`m=0.01`$ point is from the 256-node Columbia machine, the $`m=0.015`$ point is from QCDSP and the $`m=0.02`$ result was calculated on both the 256-node machine and QCDSP, which agree within errors. Parity doubling is clear for both mesons and baryons as $`m0`$.
The Columbia group has now completed a 5,000 trajectory simulation on a $`24^3\times 32`$ volume using QCDSP. Figure 5 shows that the degeneracy between the hadrons in the $`m0`$ has gone away. There is very little change in the $`m=0.02`$ masses, but for $`m=0.01`$ $`m_\rho `$ and $`m_N`$ have dropped by $`20`$% ($`m_\rho `$: $`0.438(8)0.373(6)`$, $`m_N`$: $`0.690(21)0.574(9)`$). This clearly shows finite volume effects distorting the parity splittings and here it primarily effects the nucleon and rho masses. Also, given the large decrease in $`m_\rho `$, the larger $`24^3`$ lattice has $`m_\rho L=8.9`$ very close to 7.0 for $`16^3`$.
Four flavors likely make this finite volume effect more pronounced, but it is a warning for 3 flavor simulations. We are currently simulating with 2 flavors on a $`24^3\times 32`$ volume to see how large the effects are there. Unfortunately, these large finite volume effects are masking any information about the role of the determinant in the parity splittings.
A final message about run lengths from the 4 flavor staggered simulations is shown in Figure 6. The upper line is the pion propagator at a distance 10 lattice sites from the source for the $`m=0.01`$ simulation, plotted against trajectory number. The lower line is the same propagator for the $`m=0.02`$ simulation. The $`m=0.01`$ simulation shows fluctuations on a few thousand trajectory time scale. This is clear evidence that very long runs are needed as the quark mass is made smaller. (Long autocorrelation times for topological charge for 4 flavor staggered simulations on a $`16^3\times 32`$ volume with $`\beta =5.35`$ and $`m=0.01`$ have also been seen .)
## 5 WILSON FERMIONS
The SESAM , UKQCD and CP-PACS collaborations reported on full QCD simulations last year and UKQCD and CP-PACS have new results this year. The run parameters are given in Table 2 and 3. Both UKQCD and CP-PACS are using clover improved Wilson fermions; UKQCD uses $`C_{\mathrm{SW}}`$ determined by the Alpha collaboration and CP-PACS uses a tadpole improved value.
UKQCD chooses $`\beta `$ and $`\kappa `$ to keep the lattice spacing, as determined using the Sommer scale $`r_0`$ and the static quark–anti-quark potential, constant. The physical volume then corresponds to 1.7 fm for all lattice spacings they consider. As discussed above for staggered fermions, this volume can be expected to be rather small. They do report evidence that the potential at small $`r`$ for the dynamical simulations lies below the value for quenched simulations. In addition, plotting vector meson masses against pseudoscalar meson masses, they see a trend toward the physical $`(K,K^{}`$) value.
The CP-PACS collaboration has been involved in extensive simulations of full QCD with a variety of lattice spacings, quark masses and volumes. Their parameter choices keep the spatial size fixed at $`2.4`$ fm, using the $`\rho `$ mass at the physical $`m_\pi /m_\rho `$ value to set the scale. With their full QCD data set, they can extrapolate to the continuum, with fixed finite volume effects. To date, they only have 2,000 trajectories for their $`\beta =2.10`$ point, which experience with staggered fermions suggests may not be long enough. They are addressing this issue.
CP-PACS found that a major difficulty with the quenched hadron spectrum is the failure of a unique strange quark mass to give the physical values for the $`K`$ and $`\varphi `$ masses. They have addressed this question with their new data and one of their results is shown in Figure 7. They find evidence that the $`a0`$ extrapolation of the lattice result is much closer to the physical $`K`$ mass ($``$) than the quenched $`a0`$ value. Extrapolations of the octet and decuplet baryon masses have larger errors and definitive conclusions cannot be drawn.
Given their evidence that a single value of the strange quark mass determines both the $`K`$ and $`\varphi `$ masses, they have also determined the $`a0`$ strange quark mass at $`\mu =2`$GeV. One of their results is shown in Figure 8. They find various ways of determining the strange quark mass all agree in the $`a0`$ limit and there is a large difference between the quenched and unquenched values for this mass. Their final result is $`m_s=87(11)`$ MeV ($`\varphi `$ input) and $`m_s=84(7)`$ MeV ($`K`$ input). These values are considerably lower than those from phenomenology. They also find $`m_{u,d}=3.3(4)`$ MeV.
The CP-PACS collaboration also reported a value for the flavor singlet pseudo-scalar meson for 2 flavor QCD, referred to as the $`\eta `$ . Using a volume source without gauge fixing (the Kuramashi method) to measure the disconnected quark diagrams, the mass difference between the $`\pi `$ and $`\eta `$ can be calculated, at least when the disconnected diagrams are not far separated in time. They find $`m_\eta `$ non-zero in the chiral limit and a subsequent continuum extrapolation gives $`m_\eta =863(86)`$ MeV.
UKQCD has also reported a preliminary value for the $`\eta `$ mass for their set of dynamical lattices. They employ a variance reduction technique in finding the propagator for the disconnected diagrams and get an $`\eta `$ mass around 800 MeV in the chiral limit, with an uncontrolled systematic error.
Both groups see little problem with autocorrelation times for these topologically sensitive measurements. UKQCD finds an autocorrelation time smaller than that found by the SESAM collaboration.
## 6 DOMAIN WALL FERMIONS
The development of domain wall and overlap formulations of lattice fermions has produced a great deal of excitement for both analytical and numerical studies. These formulations offer a way of separating the chiral limit from the continuum limit and may lead to a lattice formulation of non-abelian chiral gauge theories . For QCD-like theories, quenched simulations with domain wall fermions were first done by and were reviewed last year . Concurrently, simulations of the Schwinger model , including dynamical fermions, showed that domain wall fermions produced the expected physics and were compatible with standard HMC algorithms.
Two areas where using fermions with better chiral properties are of particular interest are in simulations studying QCD thermodynamics and matrix elements . The character of the QCD phase transition is governed by the chiral symmetries of the theory. For matrix element calculations, chiral symmetry can be vital for controlling operator mixing and using chiral perturbation theory as a guide. Both of these areas are under active study with the QCDSP computers, using domain wall fermions.
The Columbia group has been studying full QCD thermodynamics with domain wall fermions. The studies to locate the transition region and set the parameters to use are detailed in . As part of this work, zero temperature scale setting calculations are also required, which we will focus on here . The full QCD, zero temperature domain wall simulations done to date are detailed in Table 4. For these dynamical simulations, a heavy bosonic field, frequently called the Pauli-Villars field, is needed to remove the bulk infinity that occurs when the extent of the fifth dimension is sent to infinity.
Scale setting calculations to support thermodynamics studies are necessarily at fairly strong coupling. A first positive result is that even for these course lattices, the HMC algorithm exhibits no problems thermalizing and evolving lattices. A clear region where the Wilson line and the chiral condensate undergo a rapid crossover is seen and scale setting calculations have been done there. An important question is how the residual quark mass, $`m_{\mathrm{res}}`$ (due to mixing of the light modes between the two walls), depends on the extent of the fifth dimension.
For the PD action, the critical coupling on an $`N_t=4`$ lattice is $`\beta =5.325`$, for $`L_s=24`$, $`m_f=0.02`$ and $`M=1.9`$. Figure 9 shows hadron masses for this $`\beta `$ determined on $`8^3\times 32`$ lattices for dynamical fermion masses of 0.02 and 0.06. By extrapolating $`m_\pi ^2`$ to zero, one finds $`m_{\mathrm{res}}=0.059(2)`$. (A similar value, within errors, can also be extracted from the Ward-Takahashi identity .) This is clearly a large mass, compared to the input mass of 0.02.
Extrapolating the rho and nucleon mass to the point where $`m_\pi ^2`$ vanishes gives $`m_\rho =1.02(7)`$, $`m_N=1.14(9)`$ and $`m_\rho /m_N=1.10(12)`$. (The errors are calculated by simple propagation of errors.) While the length of the simulations ($`1,000`$ trajectories) may be sufficient at these strong couplings, only a single volume and lattice spacing has been studied. However, it is encouraging that this ratio is much closer to its physical value than for other fermion methods at this lattice spacing.
It is important to study the effects of increasing the fifth dimension on the residual mass. Figure 10 shows the residual mass versus $`L_s`$ (the length of the fifth dimension) for simulations on $`8^3\times 4`$ lattices with the PD action at $`\beta =5.2`$ with a quark mass of 0.02 and a domain wall height of 1.9. These lattices are in the confined phase. The residual mass is determined from the axial ward identity, as discussed in . The function shown is $`m_{\mathrm{res}}^{(\mathrm{GMOR})}=0.17(2)\times \mathrm{exp}(0.026(6)\times L_s)`$. The data shows a vanishing residual mass in the large $`L_s`$ limit, but the falloff is very slow. Increasing the fifth dimension by 24 drops the residual mass by about a factor of 2.
In an effort to decrease $`m_{\mathrm{res}}`$ for a fixed $`L_s`$, the Columbia group has studied the renormalization group improved action proposed by Iwasaki. For quenched simulations, $`m_{\mathrm{res}}`$ is made smaller for a given $`L_s`$, but the behavior with $`L_s`$ may not be improved . Results for hadron masses for dynamical simulations with the RD action are shown in Figure 11. These were carried out on $`8^3\times 32`$ lattices for $`\beta =2.0`$ with $`L_s=24`$. From the behavior of $`m_\pi ^2`$ one finds $`m_{\mathrm{res}}=0.013(2)`$. At the point where $`m_\pi ^2`$ vanishes $`m_\rho =0.855(49)`$, $`m_N=1.07(14)`$ and $`m_\rho /m_N=1.25(14)`$. (Once again, these are naive errors.)
While $`m_{\mathrm{res}}`$ is smaller for the RD simulations at $`\beta =2.0`$ than the PD simulations at $`\beta =5.325`$, the physical lattice scales are different. For the RD action, the $`N_t=4`$ phase transition is at $`\beta =1.9`$. We do not have two dynamical masses for the RD action at $`\beta =1.9`$, but we can compare the pion masses at $`m=0.02`$. For the PD action at $`\beta =5.325`$, $`m_\pi =0.654(3)`$ and for the RD action at $`\beta =1.9`$, $`m_\pi =0.604(3)`$, a slight improvement. (However, the uncertainty in the determination of the critical temperature could easily account for this difference.)
A direct comparison of the pion mass for the RD action ($`\beta =2.0`$, $`8^3\times 32`$) for $`L_s=24`$ gives $`m_\pi =0.475(7)`$, while for $`L_s=48`$, $`m_\pi =0.420(10)`$. $`m_{\mathrm{res}}`$ is decreasing as $`L_s`$ is increased, although the rate is slow.
## 7 CONCLUSIONS
With the current-Teraflops scale computers, full QCD simulations have reached the point where the finite volume, long simulation time, small quark mass and continuum limits can be probed, but not concurrently. The HMC and HMD algorithms continue to be the techniques of choice, but the multiboson algorithm has been evolved to a competitive level. Staggered 4 flavor QCD shows large finite volume effects, even for $`m_\pi /m_\rho >0.5`$ and the parity splittings for light hadrons are quite volume sensitive. For currently accessible weaker coupling simulations, runs of length greater than 10,000 trajectories are probably needed.
CP-PACS has improved their results on the continuum limit of 2 flavor QCD and find continuum meson masses closer to the physical values than for the quenched case. This leads them to a determination of the strange quark mass which is lower than that expected phenomenologically. For staggered fermions, the new 2 flavor points being produced by Columbia will augment the existing continuum extrapolation of the MILC group.
Full QCD domain wall fermion simulations are straightforward, but require a large fifth dimension at strong coupling. It is encouraging that on very coarse lattices $`m_N/m_\rho `$ is much lower than for other fermion formulations. Further studies will be needed to check the scaling properties for domain wall fermions to see if the use of coarser lattices can offset the extra calculational cost of the fifth dimension.
ACKNOWLDEGEMENTS
The author would like to thank R. Burkhalter, P. de Forcrand, A. Duncan, J. Garden, S. Gottlieb, K. Kanaya, C. Michael, U. Wiese, the CP-PACS collaboration and the MILC collaboration for providing access to their results. |
warning/0001/nucl-th0001021.html | ar5iv | text | # Untitled Document
Further Comments On The Effects of Deformation on Isovector
Electromagnetic and Weak Transition Strengths
Shadow J.Q. Robinson<sup>a</sup>, L. Zamick<sup>a,b</sup>, A. Mekjian<sup>a</sup>, N. Auerbach<sup>a,c</sup>
a) Department of Physics and Astronomy, Rutgers University, Piscataway,
New Jersey 08855
b) TRIUMF,4004 Wesbrook Mall, Vancouver, British Columbia,
Canada, V6T 2A3
c) Permanent address, School of Physics, Tel-Aviv University, Tel Aviv,
Israel.
## Abstract
We present a superior proof that the results for summed strength isovector dipole, spin dipole, and orbital dipole excitations are independent of deformations at the $`\mathrm{\Delta }`$ N = 0 level. The effects of different oscillator frequencies in the x, y, and z directions are also considered.
1) INTRODUCTION
As has been previously noted , using the rotational model for $`{}_{}{}^{12}C`$ and harmonic oscillator wave functions, the results for summed strength isovector dipole, spin-dipole, and orbital dipole excitations were independent of deformation in a $`\mathrm{\Delta }N=0`$ Nilsson model. We then presented a more general proof which did not require the use of the rotational model explicitly but rather did require that the valence nucleons were all in the $`0p`$ shell and that the mean square radius of $`p_{1/2}`$ and $`p_{3/2}`$ particles were the same (as they are with harmonic oscillator wave functions).
We here present a superior proof and make several points about dipole excitations. We consider excitations from the ground state of an $`N=Z`$ open shell nucleus (like $`{}_{}{}^{12}C`$). We will assume the ground state has angular momentum $`J=0^+`$. As in the original work we consider the operators ($`rY_k^1t`$, $`r[Y^1s]_k^\lambda t`$, and $`r[Y^1\mathrm{}]_k^\lambda t`$ .). Some of these operators arise in (p,n) reactions or neutrino absorption such as $`\nu _e+^{12}C^{12}N+e^{}`$.
We show again in Table 1 the results of the summed strength in the asymptotic (oblate) limit and the spherical limit for the above operators in $`{}_{}{}^{12}C`$. The results to individual final momenta $`\lambda `$ are different in these two limits, but the total summed strength is the same in these two limits.
For the ordinary dipole operator $`rY_k^{}t_+`$, the summed strength (SUM) , multiplied by $`4\pi m\omega `$/$`\mathrm{}`$ is 27; for the spin dipole it is 20.25 and for the orbital dipole 48. We will soon explain why this is so.
2) THE NEW APPROACH
To see why the results for SUM are independent of the specific $`0p`$ configuration (or deformation) when spherical harmonic oscillators (H.O.) wavefunctions are used we note the following unique feature of dipole excitations: In the H.O. approximation there is only one excitation energy, $`1\mathrm{}\omega `$. For the other modes this is not the case. For $`E2`$ transitions, the strength of which are highly dependent on deformation there are both $`0\mathrm{}\omega `$ and $`2\mathrm{}\omega `$ excitations; for $`E3`$ we have $`1\mathrm{}\omega `$ and $`3\mathrm{}\omega `$ excitations etc.
Since for $`E1`$ transitions there is only one excitation energy involved we can relate the summed strength to the energy weighted strength E.W.S.
$$SUM=E.W.S./\mathrm{}\omega $$
(1)
The energy weighted strengths have been studied a great deal , and if we ignore, for the moment, the lack of commutivity of the potential energy with the various dipole operators, very simple results emerge.
Let us first show the electric dipole EWS referred to the center of mass, as given in Bohr and Mottelson. They write the operator $`M(E1,\mu )=e_i(\frac{NZ}{2A}t_3(i))(rY_\mu ^1)_i`$ The ’classical’ EWS for this operator is
$$EWS=\frac{9}{4\pi }\frac{\mathrm{}^2}{2M}\frac{NZ}{A}$$
(2)
Which for $`N=Z=\frac{A}{2}`$ becomes
$$EWS=\frac{9}{32\pi }A\frac{\mathrm{}^2}{M}$$
(3)
In our problem we have $`t_+`$ rather than $`t_z`$. Again going to the case of $`N=Z=\frac{A}{2}`$ , we have $`M(E1,\mu )=e_i(t_+(i))(rY_\mu ^1)_i`$ The EWS is now expressed as
$$\frac{1}{2}[EWS(+)+EWS()]=\frac{9}{8\pi }<0|[zt_{},[\frac{\mathrm{}^2}{2M}\frac{d^2}{dz^2},zt_+]]|0>$$
(4)
where EWS(+) is the energy weighted strength for a process in which a neutron is changed into a proton and EWS(-) where a proton is changed into a neutron. Using the relations
$$[\frac{d^2}{dz^2},zt_+]=2\frac{d}{dz}t_+$$
(5)
$$[z,\frac{d}{dz}]=1$$
(6)
$$[t_{},t_+]=2t_z$$
(7)
We are reduced to
$$\frac{1}{2}[EWS(+)+EWS()]=\frac{9}{8\pi }\frac{\mathrm{}^2}{M}<0|t_z+\frac{1}{2}+t_z2z\frac{d}{dz}|0>$$
(8)
We can easily compute $`<2z\frac{d}{dz}>`$ by integration by parts (given real wavefunctions).
$`<2z{\displaystyle \frac{d}{dz}}>={\displaystyle \psi 2z\frac{d}{dz}\psi }=I_T`$
$`I_T=\psi ^22z{\displaystyle }\psi (2\psi +2z{\displaystyle \frac{d\psi }{dz}}`$
$`I_T=02I_T`$
$`I_T=1`$
$`<2z{\displaystyle \frac{d}{dz}}>=1`$ (9)
This yields the simple result first derived by Lipparini and Stringari
$$\frac{1}{2}[EWS(+)+EWS()]=\frac{9}{8\pi }\frac{\mathrm{}^2}{M}<0|t_z+\frac{1}{2}t_z|0>$$
(10)
For N=Z we have
$$EWS(+)=EWS()=\frac{9}{16\pi }A\frac{\mathrm{}^2}{M}$$
(11)
since in this case, $`EWS(+)=EWS()`$. For the SUM we obtain
$$SUM=\frac{EWS(rY_k^1t_+)}{\mathrm{}\omega }=\frac{9}{16\pi }A\frac{\mathrm{}}{M\omega }$$
(12)
Finally we get
$$4\pi SUM\frac{M\omega }{\mathrm{}}=9\frac{A}{4}$$
(13)
This is the quantity given in Tables 1 and 2 of ref . For $`A=12`$ we get 27 for this quantity, confirming the results previously obtained.
3) EFFECT OF DIFFERENT FREQUENCIES IN THE X,Y,AND Z DIRECTIONS
To take deformation effects further into account we introduce different frequencies in the x, y, and z directions. It can be shown that we get the correct result by making the following replacement in Eq(12).
$$\frac{1}{\mathrm{}\omega }\frac{1}{3}(\frac{1}{\mathrm{}\omega _x}+\frac{1}{\mathrm{}\omega _y}+\frac{1}{\mathrm{}\omega _z})$$
(14)
To obtain this result we must not only consider excitations from 0p to higher shells but also excitations from 0s to 0p. Note that if the above expression (14) is expanded in terms of a deformation parameter $`\delta `$ there will be no linear terms.
To get an estimate of the size of this effect we use the self consistency conditions
$$\mathrm{\Sigma }_x\omega _x=\mathrm{\Sigma }_y\omega _y=\mathrm{\Sigma }_z\omega _z$$
(15)
where for <sup>12</sup>C in the asymptotic limit
$`\mathrm{\Sigma }_x=\mathrm{\Sigma }_y=10`$
$`\mathrm{\Sigma }_z=6`$ (16)
We define $`\omega _0`$ by $`\omega _x\omega _y\omega _z=\omega _0^3`$ and assume volume conservation, i.e. keep $`\omega _0`$ constant. We then find $`\omega _x=0.8434\omega _0`$ and $`\omega _z=1.4057\omega _0`$. We find $`\frac{1}{3}(\frac{1}{\mathrm{}\omega _x}+\frac{1}{\mathrm{}\omega _y}+\frac{1}{\mathrm{}\omega _z})=\frac{1.0275}{\mathrm{}\omega _0}`$
There is a very small change in the overall strength. However $`\frac{2}{3}`$ of the strength is shifted down to $`0.8434\mathrm{}\omega _0`$ and $`\frac{1}{3}`$ is shifted up to $`1.4057\mathrm{}\omega _0`$. (Obviously the energy weighted strength does not change in this model.)
4) SPIN DIPOLE AND ORBITAL DIPOLE MODES
We next consider the spin-dipole mode and consider the EWSR in which only the kinetic energy is taken into account
$`EWS(spinmultipole)={\displaystyle \underset{\lambda M}{}}{\displaystyle \underset{i}{}}`$
$`{\displaystyle \frac{1}{2}}<[[Y^L(i)s(i)]_M^\lambda ,[{\displaystyle \frac{p^2(i)}{2m}},[Y^L(i)s(i)]_M^\lambda ]>`$ (17)
$`={\displaystyle \underset{i}{}}{\displaystyle \underset{L,M,M_L,M_S,M_L^{^{}},M_S^{^{}}}{}}(L1M_LM_S|\lambda M)(L1M_L^{^{}}M_S^{^{}}|\lambda M)`$
$`{\displaystyle \frac{1}{2}}<[Y_{M_L}^L(i)s_{M_S}^{}(i),[{\displaystyle \frac{p^2(i)}{2m}},Y_{M_L^{^{}}}^L(i)s_{M_S^{^{}}}(i)>`$ (18)
Now since
$$\underset{M,L}{}(L1M_LM_S|\lambda M)(L1M_L^{^{}}M_S^{^{}}|\lambda M)=\delta _{M_L,M_L^{^{}}}\delta _{M_s,M_s^{^{}}}$$
(19)
We obtain
$`EWS(spinmultipole)={\displaystyle \frac{1}{2}}{\displaystyle \underset{i}{}}{\displaystyle \underset{M_L}{}}<[Y_{M_L}^L(i),[{\displaystyle \frac{p^2(i)}{2m}},Y_{M_L}^L(i)]]{\displaystyle s_{M_s}^{}s_{M_s}}>`$ (20)
$`={\displaystyle \underset{i}{}}EWS(ordinarymultipole)s(i)s(i)`$ (21)
Now $`s(i)s(i)`$ is equal to $`3/4`$ for spin $`1/2`$ particles. Hence $`EWS`$(spin multipole) $`=`$ 3/4 $`EWS`$(ordinary multipole) For the spin dipole case in which there is only one excitation energy ($`\mathrm{}\omega `$) the above relation also holds for summed strength. This was mentioned but not proven by Auerbach and Zamick.
As noted in their work and in Table 1 of the present work, the value of SUM for the spin-dipole case is 20.25 which is indeed 3/4 of the ordinary dipole 27.
For the orbital dipole case we replace $`s(i)s(i)`$ by $`l(i)l(i)`$. This differs from the spin-dipole case in the sense that $`l(i)l(i)`$ is state dependent with eigenvalue l(l+1) . The value of SUM in Table 1 (summed also over $`\lambda `$) is 48, all coming from the 0p shell. The value of SUM for the ordinary dipole coming from the 0p shell is 27-3=24. The factor of 2 is due to the fact that $`l(l+1)`$ equals two in the 0p shell.
This work was supported by the U.S. Department of Energy under grant DE-FG02-95ER-40940 and DE-FG02-96ER-40987. One of us (N.A.) is grateful for the hospitality afforded to him at Rutgers University.
Table 1 Total isovector dipole strength in $`{}_{}{}^{12}C(\frac{4\pi m\omega }{\mathrm{}}SUM(\lambda )0^+\lambda )`$
in parenthesis are the strengths due to excitations from $`0s`$.
| Dipole | $`rY_K^1t`$ | | |
| --- | --- | --- | --- |
| | Asymptotic | Spherical | |
| $`\lambda `$ | | | |
| 1 | 27 | 27 | |
| |
| |
| |
| Spin Dipole | $`r[Y^1s]^\lambda t`$ | | |
| $`\lambda `$ | | | |
| 0 | 2.25 (0.25) | 3.25 (0.75) | |
| 1 | 6.75 (0.75) | 8.25 (1.50) | |
| 2 | 11.25 (1.25) | 8.75 (0) | |
| Sum | 20.25 (2.25) | 20.25 (2.25) | |
| |
| |
| |
| Orbital Dipole | Dipole $`r[Y^1\mathrm{}]^\lambda t`$ | | |
| $`\lambda `$ | | | |
| 0 | 0 (0) | 0 (0) | |
| 1 | 14 (0) | 14 (0) | |
| 2 | 34 (0) | 34 (0) | |
| Sum | 48 (0) | 48 (0) | |
References
1. L. Zamick and N. Auerbach, Nuclear Physics A 658, 285 (1999)
2. A. Bohr and B. Mottelson, Nuclear Structure, Vol. 1 and Vol. 2 (1969) and (1975), (Benjamin, New York, 1975).
3. W. Lipparini and S. Stringari, Phys. Rep. 175, 103 (1989) |
warning/0001/hep-th0001021.html | ar5iv | text | # A Two-loop Test of Buscher’s T-duality I
## 1 Introduction
Among the wealth of different dualities relating the perturbative string theories and M-theory, the most ’ancient’ one is T-duality. At the same time T-duality was the starting point in the discovery of D-branes . In string theory T-duality can be proven in arbitrary order of the string perturbation theory , as long as the vacuum preserves conformal invariance . Generalizing the $`R\frac{1}{R}`$ symmetry of the toroidally compactified strings, Buscher gave a set of equations that describe the transformation of the Neveu-Schwarz background fields of the sigma model action, and works whenever there is an isometry. Denoting by $`g_{\mu \nu }`$ and $`b_{\mu \nu }`$ the background metric and the antisymmetric tensor field, the explicit form of the Buscher transformation is (since we work on a flat world sheet, we neglect the dilaton):
| $`\stackrel{~}{g}_{00}={\displaystyle \frac{1}{g_{00}}},\stackrel{~}{g}_{0\alpha }={\displaystyle \frac{b_{0\alpha }}{g_{00}}},\stackrel{~}{b}_{0\alpha }={\displaystyle \frac{g_{0\alpha }}{g_{00}}},`$ |
| --- |
| $`\stackrel{~}{g}_{\alpha \beta }=g_{\alpha \beta }{\displaystyle \frac{g_{0\alpha }g_{0\beta }b_{0\alpha }b_{0\beta }}{g_{00}}},\stackrel{~}{b}_{\alpha \beta }=b_{\alpha \beta }{\displaystyle \frac{g_{0\alpha }b_{0\beta }g_{0\beta }b_{0\alpha }}{g_{00}}},`$ |
(1)
where $`\stackrel{~}{g}_{\mu \nu }`$ and $`\stackrel{~}{b}_{\mu \nu }`$ denote the new background fields.
To derive this transformation Buscher used functional integral arguments, that in the meantime have become widely known, and applied in many context (see e.g. ). The idea is to gauge the aforementioned isometry of the sigma model action and impose a constraint using a Lagrange multiplier. Integrating the multiplier one recovers the original theory to start with, while after a gauge fixing and integrating over one of the original fields gives the dual theory.
It was also shown , that, at the classical level, the duality transformation rule can be recovered in an elegant way by performing a canonical transformation . This clearly shows that the models connected by the Buscher transformations are equivalent classically. At the quantum level the only loophole of the path integral argument mentioned above is that it neglects the effect of the regularization and renormalization. Though it was shown in that for conformal invariant models one has full quantum equivalence, in other words the path integral argument holds, this is not the case in the non-conformal setting. It was shown in for a deformed $`SU(2)`$ principal sigma model, and further clarified in for several different deformed $`SU(3)`$ principal sigma models, that Buscher’s formula – as applied to renormalized quantities – has to be modified to give two loop quantum equivalence.
Later it was shown that the one loop beta functions of the original and dual models always agree . Work has been done to establish the corrections to Buscher’s formulae . Advances were made from a different point of view. It is widely known that the low energy degrees of freedom of the sigma model can be described using an effective action, that contains gravity in target space. This fact constrains the possible leading terms of the low energy effective action to a computable form, that is known. It was shown that the low energy effective action consistent with the two loop sigma model beta function equations is not invariant under Buscher’s transformation. The leading part of the above action, the one loop part, nevertheless is invariant, in accordance with the one loop findings of . Using the non-invariance of the two loop action it was found how to modify Buscher’s transformation, with order $`\alpha ^{}`$ terms, such that the two loop action would remain invariant under the modified transformations. What is not clear after all is how to pull back the modification of the Buscher’s formulae found in the low energy action to the sigma model. Related work in the supersymmetric case, concentrating on the absence of the above mentioned corrections, was done in .
Any satisfactory criterion of the quantum equivalence among dually related sigma models should be based on the comparison of physical quantities as opposed to just considering beta functions. If there are global symmetries in the model then their associated conserved quantities (Noether currents) may be considered physical. The definition of physical quantities, however, is not very clear in diffeomorphic invariant sigma models without a sufficient number of isometries. To circumvent this problem the study of Weyl anomaly coefficients was suggested in . In the present paper – as an alternative – we study a thermodynamic quantity namely the free energy density in the presence of a chemical potential in the dually related sigma models. This quantity surely qualifies as physical, thus its equality in the two models gives a non trivial check on their quantum equivalence. Furthermore the free energy density can be computed perturbatively – at least in asymptotically free models, thus one can compare the two free energy densities using the first few orders of perturbation theory. The aim of this paper is to carry out this comparison in the two loop order, where the first really ‘quantum’ effects appear, thus improving the almost ‘classical’ one loop case studied earlier in , .
The paper is organized as follows: in Section 2 we give a brief review of the pertinent facts that we need from the renormalization of the $`SU(2)`$ principal $`\sigma `$ model and its dual and develop a Lagrangian with more parameters, that can accommodate both Lagrangians as special cases. In Section 3 we investigate the conserved charges that can be coupled to the models, calculate the Hamiltonian, find the ground states in the presence of the external field, and make the Lagrangians suitable for a perturbative computation. Section 4 deals with the definition of the perturbative free energy density, and its computation, using dimensional regularization. First we define the perturbative free energy density, and set up a scheme to compute it systematically to any order. The rest of Section 4 deals with the actual computation of the bare free energy up to two loops. In Section 5 we deal with the issue of renormalization and obtain the one and two loop renormalized free energies of the original and dual models, compare them, and improve them using the renormalization group. We close Section 5 with the analysis of the composite operator renormalization of the relevant operators. We make our conclusions in Section 6.
## 2 Lagrangians and T-duality
### 2.1 The deformed $`SU(2)`$ principal $`\sigma `$ model
This section has a twofold role. Primarily it is intended to give an overview of the results that we need in the rest of the paper and at the same time fix the notations. Secondarily it extends some of the earlier results in a way suitable for our applications.
In the following one parameter deformation of the $`SU(2)`$ principal $`\sigma `$-model was considered :
$$_O=\frac{1}{2\lambda }\left(\underset{a=1}{\overset{3}{}}J_\mu ^aJ^{\mu a}+gJ_\mu ^3J^{\mu 3}\right),$$
(2)
where $`J_\mu =G^1_\mu G=J_\mu ^a\tau ^a`$, and $`\tau ^a=\sigma ^a/2`$ with $`\sigma ^a`$ being the standard Pauli matrices. Thus $`G`$ is an element of $`SU(2)`$ and $`g`$ is the parameter of the deformation. From the Lagrangian (2) it is clear that the global $`SU(2)_L\times SU(2)_R`$ symmetry of the undeformed principal $`\sigma `$-model is broken to $`SU(2)_L\times U(1)_R`$ by the $`J_\mu ^3J^{\mu 3}`$ term. Setting $`g=0`$ corresponds to the principal $`\sigma `$-model, while for $`g=1`$ the $`O(3)`$ $`\sigma `$-model is obtained as it can be seen from eq. (4) below.
The authors of investigated the renormalization of $`\lambda `$ and $`g`$ in the two-loop order of perturbation theory, treating $`\lambda `$ as the coupling constant and $`g`$ as a parameter. Using the Euler angles ($`\varphi `$,$`\theta `$,$`\psi `$) to parameterize the elements of $`SU(2)`$, $`G`$ is written as
$$G=e^{i\varphi \tau ^3}e^{i\theta \tau ^1}e^{i\psi \tau ^3}.$$
(3)
Using this converts the Lagrangian of the deformed $`\sigma `$ model, which for the time being we shall call ’the original model’, into the following form:
| $`_O=`$ | $`{\displaystyle \frac{1}{2\lambda }}\{(_\mu \theta )^2+(_\mu \varphi )^2(1+g\mathrm{cos}^2\theta )+`$ |
| --- | --- |
| | $`+(1+g)(_\mu \psi )^2+2(1+g)_\mu \varphi ^\mu \psi \mathrm{cos}\theta \}.`$ |
(4)
Using the Killing vectors of the $`SU(2)_L\times U(1)_R`$ symmetry and exploiting the manifest target space covariance of the background field method it was proved in that the model is renormalizable in the ordinary sense: there is no wave function renormalization for $`\theta `$, $`\varphi `$ and $`\psi `$, while the coupling constant and the parameter got renormalized according to:
$$\lambda _0=\mu ^ϵZ_\lambda (\lambda ,g)\lambda ,g_0=Z_g(\lambda ,g)g.$$
(5)
Both in the one and in the two loop orders the residues of the single poles in $`Z_\lambda (\lambda ,g)=1y_\lambda (\lambda ,g)/ϵ+\mathrm{}`$ and $`Z_g(\lambda ,g)=1y_g(\lambda ,g)/ϵ+\mathrm{}`$ were determined:
| $`y_\lambda =`$ | $`{\displaystyle \frac{\lambda }{4\pi }}\left(1g+{\displaystyle \frac{\lambda }{16\pi }}(12g+5g^2)\right),`$ |
| --- | --- |
| $`y_g=`$ | $`{\displaystyle \frac{\lambda }{2\pi }}(1+g)\left(1+{\displaystyle \frac{\lambda }{8\pi }}(1g)\right).`$ |
(6)
Note the sign difference between our formulas (5) and (6), and the corresponding ones in . It is consequence of the fact that in our notation notation $`n=2+ϵ`$ rather than $`n=2ϵ`$ as used in .
The standard definition of the $`\beta `$ functions: $`\beta _\alpha =\mu \frac{d\alpha }{d\mu }`$, $`\beta _\gamma =\mu \frac{d\gamma }{d\mu }`$, lead to the following two-loop $`\beta `$ functions (eq. (20) in ):
| $`\beta _\lambda =`$ | $`{\displaystyle \frac{\lambda ^2}{4\pi }}\left(1g+{\displaystyle \frac{\lambda }{8\pi }}(12g+5g^2)\right),`$ |
| --- | --- |
| $`\beta _g=`$ | $`{\displaystyle \frac{\lambda }{2\pi }}g(1+g)\left(1+{\displaystyle \frac{\lambda }{4\pi }}(1g)\right).`$ |
(7)
It is easy to see, that the $`g=0`$ resp. the $`g=1`$ lines are fixed lines under the renormalization group flow, and $`\beta _\lambda `$ reduces to the $`\beta `$ function of the principal $`\sigma `$-model, resp. of the $`O(3)`$ $`\sigma `$-model on them. In the ($`\lambda 0,g<0`$) quarter of the ($`\lambda ,g`$) plane the renormalization group trajectories run into $`\lambda =0,g=1`$; while for $`g>0`$ they run to infinity. This implies that the $`g=0`$ fixed line corresponding to the principal $`\sigma `$-model is ‘unstable’ under the deformation.
The Lagrangian of the deformed $`\sigma `$-model, eq. (4), exhibits two obvious Abelian isometries that can be used to construct two different (Abelian) duals: namely the translations in the $`\varphi `$ and $`\psi `$ fields. We call the models obtained this way the ‘$`\varphi `$ dual’ respectively the ‘$`\psi `$ dual’ of the deformed $`\sigma `$ model (4).
In it was found that for the ‘$`\psi `$ dual’ model, as summarized below, the renormalization of the coupling and the parameter are equivalent to that of the original model. Therefore, in the present context we deal with the original, deformed $`SU(2)`$ principal $`\sigma `$-model and it’s ‘$`\psi `$ dual’, which we shall simply call ’the dual model’.
### 2.2 The ‘$`\psi `$ dual’ model
For the Lagrangian of the ’$`\psi `$ -dual’ model, using Buscher’s formulae , one has an expression analogous to eq. (4):
| $`_D=`$ | $`{\displaystyle \frac{1}{2\stackrel{~}{\lambda }}}\left((_\mu \theta )^2+(_\mu \varphi )^2\mathrm{sin}^2\theta +(_\mu \chi )^2+2a\mathrm{cos}\theta ϵ^{\mu \nu }_\mu \chi _\nu \varphi \right).`$ |
| --- | --- |
(8)
Here $`\chi `$ denotes the (appropriately scaled) variable dual to $`\psi `$, and ($`\stackrel{~}{\lambda },\stackrel{~}{g}`$) stands for the couplings of the dual model. One can show that $`_D`$ exhibits the expected $`SU(2)\times U(1)`$ symmetry for all values of the parameter $`a`$.
The couplings of the original, (4), and of the dual models, (8), are related (at the classical level) as a direct consequence of T-duality as
$$\stackrel{~}{\lambda }=\lambda ,a=\sqrt{1+\stackrel{~}{g}},\stackrel{~}{g}=g.$$
(9)
These relations can be maintained at the two-loop level if one performs the renormalization and, in addition, in both theories the couplings $`(g,\stackrel{~}{g})`$ are expressed in terms of the corresponding renormalization group invariant quantities, which at the end are set equal to each other . In addition one also has to take into account the freedom in the choice of the renormalization group invariant, due to the scheme dependence of the two-loop $`\beta `$ function.
The two-loop renormalization invariants, that characterize the flows under the corresponding sets of $`\beta `$ functions, for the original and dual models can be easily computed:
| $`M_O={\displaystyle \frac{g}{(1+g)^2}}\lambda ^2{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{g}{1+g}}\lambda ^3+𝚘(\lambda ^4),`$ |
| --- |
| $`M_D={\displaystyle \frac{a^21}{a^4}}\stackrel{~}{\lambda }^2{\displaystyle \frac{1}{4\pi }}{\displaystyle \frac{1}{a^2}}\stackrel{~}{\lambda }^3+𝚘(\stackrel{~}{\lambda }^4).`$ |
(10)
The next step is to express $`g`$ (resp. $`a`$) in terms of the renormalization group invariant. This can be done by inverting eq. (10) perturbatively. In both cases at leading order one has to solve quadratic equations and the next to leading order terms correct the results including the two-loop effects. Assuming that neither $`M_O`$ nor $`M_D`$ vanishes, one obtains (the leading terms were obtained in , but we shall need the whole expression):
| $`g(\lambda ,M_O)=1\pm {\displaystyle \frac{1}{\sqrt{M_O}}}\lambda {\displaystyle \frac{1}{2M_O}}\lambda ^2+𝚘(\lambda ^3),`$ |
| --- |
| $`a^2(\stackrel{~}{\lambda },M_D)=\pm {\displaystyle \frac{1}{\sqrt{M_D}}}\stackrel{~}{\lambda }{\displaystyle \frac{1}{2M_D}}\stackrel{~}{\lambda }^2+𝚘(\stackrel{~}{\lambda }^3).`$ |
(11)
The sign ambiguity can be removed by studying the renormalization group flows, as it was briefly mentioned in connection with eq. (7). In the original model, it turns out that the interesting region is the vicinity of $`g=1`$ with $`g1`$. In the dual model, since $`a=\sqrt{1+g}`$, choosing $`g1`$ one is uniquely led to the solution with the plus sign. Thus in both cases one has to consider the solution with the plus sign. Moreover if one sets $`M_O`$ and $`M_D`$ equal, $`M_O=M_D=M`$, then the classical $`a=\sqrt{1+g}`$ relation can be maintained. We note (though we didn’t display it in eq. (11)) that the two expressions differ already at the order of $`\lambda ^3`$, as expected, since $`M`$ is renormalization group invariant only up to two loop order.
Using $`g(\lambda ,M_O)`$ ($`a^2(\stackrel{~}{\lambda },M_D)`$) in the two loop beta functions of the coupling constants of the original and dual models yields a universal expression :
$$\beta _\lambda =\frac{\lambda ^2}{4\pi }\left(2+\lambda (\frac{1}{\pi }\frac{1}{\sqrt{M}})\right).$$
(12)
As far as the coupling constant renormalization is concerned, this universal beta function shows that the two models are equivalent, both are asymptotically free, and the actual value of $`M`$ effects only the two loop coefficient.
According to eq. (11) one can express $`a`$ in terms of the common renormalization group (RG) invariant. For the sake of simplicity it is useful to introduce the following notations:
$$A=\sqrt{\lambda },\alpha _1=\frac{1}{\sqrt[4]{M}},\alpha _2=\frac{1}{4\sqrt[4]{M^3}}.$$
(13)
It will turn out that $`A=\sqrt{\lambda }`$ is the proper coupling of the two models. In terms of these
$$a=\alpha _1A\alpha _2A^3+𝚘(A^5),g=1+\alpha _1^2A^2\frac{1}{2}\alpha _1^4A^4+𝚘(A^6).$$
(14)
### 2.3 Unified description
In order to make the computations more general it is useful to express both theories as particular cases of a generalized $`\sigma `$ model. To achieve this we rescale the field $`\psi \frac{1}{\sqrt{1+g}}\psi `$ and introduce $`a=\sqrt{1+g}`$ in place of $`g`$ in the original model. After some obvious changes of symbols, $`_O`$ and $`_D`$ can be described by the following unified $`\sigma `$ model:
| $`=`$ | $`{\displaystyle \frac{1}{2\lambda }}\left\{(_\mu \theta )^2+(_\mu \varphi )^2\left(1+r\mathrm{cos}^2\theta \right)+(_\mu \psi )^2+2a\omega _{\mu \nu }\mathrm{cos}\theta _\mu \varphi ^\mu \psi \right\}.`$ |
| --- | --- |
(15)
The previous models can be recovered by the special choices: $`r=g`$ and $`\omega _{\mu \nu }=\eta _{\mu \nu }`$ for the original model, and $`r=1`$ and $`\omega _{\mu \nu }=ϵ_{\mu \nu }`$ for the dual model. Expressing $`r`$ in terms of the RG invariant parameter we have:
$$r=1+2\beta A^2+\gamma A^4+𝚘(A^6),$$
(16)
where $`\beta =\alpha _1^2/2`$ and $`\gamma =\frac{1}{4}\alpha _1^4`$ in the original resp. $`\beta =0`$ $`\gamma =0`$ in the dual model. Observe that both $`\beta `$ and $`\gamma `$ are renormalization group (RG) invariant.
This unified Lagrangian (15) is more than it might appear at first sight. As it was shown in and argued above, the deformed principal $`\sigma `$ model and its dual can be viewed as being quantum equivalent from the point of view of the two loop beta functions. The unifying Lagrangian (15) can be viewed as a genuine quantum generalization of the deformed principal $`\sigma `$ model and its dual. It reduces to the latter ones at special values of the parameters $`\beta ,\gamma `$ and $`\omega _{\mu \nu }`$, and at the same time gives the corresponding beta functions. Thus the renormalization properties of the two models are encompassed in this generalized Lagrangian. Based on this we will be able to renormalize the free energy in both theories at the same time, shortening the computations and obtaining better control on the different contributions.. Checking the quantum equivalence will amount to compare the renormalized free energies computed at the special values of the parameters $`\beta ,\gamma `$ and $`\omega _{\mu \nu }`$ corresponding to the two models.
## 3 The ground state and perturbative Lagrangian
### 3.1 Outline of the method - Noether currents
So far testing the quantum equivalence of the dual models was mostly reduced to comparisons of the corresponding beta-functions . Of course there are more to test before one can ascertain about an equivalence. In this paper to test the physical equivalence between the original and dual theories we couple both of them to some particular conserved charge $`Q`$. This is accomplished by modifying the respective Hamiltonians $`H_O`$ $`(H_D)`$ to $`H_OhQ`$ $`(H_DhQ)`$, where $`h`$ is an external field (chemical potential type of parameter), having mass dimension one. The corresponding changes in the ground state energy densities (i.e. in the free energy densities, that we shall call for simplicity free energies) can be computed, at least in principle, to any order in perturbation theory. The comparison of the free energy densities, as functions of $`h`$, up to a certain order of perturbation theory, in the original and dual models, then provides a useful check whether the two models do really, physically correspond to each other. For a comparison of this type to make sense both theories must be asymptotically free (to guarantee that perturbation theory applies), and of course we have to choose $`Q`$ to be really the same.
As the global symmetry group of both the original and the dual models is $`SU(2)\times U(1)`$ we may think naively that any linear combination of an $`SU(2)`$ Noether charge and the $`U(1)`$ Noether charge, $`x^aQ^a+yQ_{U(1)}`$, can be used as the charge above $`Q`$ to couple to the Hamiltonians. However for our test we need the same $`Q`$ coupled to $`H_O`$ and $`H_D`$, thus we can choose only such charges that are mapped to themselves under the canonical transformation connecting the original and dual models. This, of course, implies that the charge $`Q`$ must stay local under the canonical transformation.
We first point out the relation between the appropriate Noether currents and charges of the global symmetries of $`_O`$ and $`_D`$. An exhaustive treatment will be given elsewhere . Let’s start remarking that under Abelian duality transformations the image of the $`U(1)`$ current, that belongs to the distinguished isometry used in duality, is a topological current built from the dual field. Thus the image of the $`U(1)_R`$ current of the deformed sigma model is the topological current of the $`\chi `$ field, $`ϵ_{\mu \nu }^\nu \chi `$, and as such, its charge should vanish on a topologically trivial 2d space-time. Therefore $`Q_{U(1)}`$ cannot be present in both $`H_O`$ and $`H_D`$.
The next simplest possibility is to use the $`U(1)`$ charges corresponding to the $`\varphi `$-translations in both the original and dual theories: $`N_0^3`$ and $`\stackrel{~}{N}_0^3`$. It can be shown (see for details) that the canonical transformation which implements Abelian duality, effectively exchanges only $`p_\psi `$ and $`\chi ^{}`$ ($`p_\chi `$ and $`\psi ^{}`$) and leaves $`p_\varphi `$, $`p_\theta `$ unchanged. It is obvious that $`\stackrel{~}{N}_0^3`$ is ’identical’ under duality to $`N_0^3`$, but it is not entirely trivial (though it is true) that the space component of $`N_\mu ^3`$ really becomes the space component of $`\stackrel{~}{N}_\mu ^3`$.
### 3.2 From the Hamiltonian to the perturbative Lagrangian
In conclusion, to test the physical equivalence of the original and dual models, in both of them we introduce the coupled Hamiltonian densities $`\overline{}_{O,D}=_{O,D}hN_0^3`$.<sup>1</sup><sup>1</sup>1One can also investigate different linear combinations of the Noether charges . Performing the inverse Legendre transformation on these quantities one obtains the $`\overline{}_{O,D}`$ Lagrange densities of the coupled models. By explicit computation one can show that this procedure of obtaining $`\overline{}_{O,D}`$ is equivalent to the following formal gauging ($`_\mu \zeta ^iD_\mu \zeta ^i`$) of the original Lagrangian $`_{O,D}`$:
$$D_\mu \psi =_\mu \psi ,D_\mu \varphi =_\mu \varphi +h\delta _{0\mu },D_\mu \theta =_\mu \theta .$$
(17)
This gauging, of course, can be done on the universal, common Lagrangian (15); the outcome being denoted by $`\overline{}`$. As it can be seen in (17), the coupling of the chemical potential explicitly breaks Lorentz invariance. This will play an important role in our analysis of the quantum equivalence.
To perform the perturbative analysis first one has to determine the classical ground state of the system, in other words the minima of the Hamiltonian $`\overline{}=\overline{}(\theta ,\varphi ,\psi ,p_\theta ,p_\varphi ,p_\psi )`$, finding the critical points and checking the Hessian’s positive definiteness. Using the actual (perturbative) expression for r, eq. (16), it is straightforward to show that in the perturbative region around $`r=1`$ (which in the original model corresponds to $`g=1`$ as opposed to the dual model where it is exactly -1) there is a two-parameter family of local minima, all of them being physically equivalent, given by: $`(p_\theta =p_\psi =0,p_\varphi =\frac{h}{\lambda },\theta =\frac{\pi }{2},\varphi =const.,\psi =const.)`$. Since $`\varphi `$ and $`\psi `$ does not appear explicitly in the Hamiltonian, we can chose for convenience the stable classical ground state, common to both models, to be given by $`\theta =\frac{\pi }{2},\varphi =\psi =0`$. We shall expand our fields around this background.
The Lagrangian, $`\overline{}`$, as emerging from eq. (15) is not suitable for a direct perturbative computation on account of the overall $`1/\lambda `$ factor. However this factor can be removed by an appropriate rescaling of the fields:
$$\stackrel{~}{\theta }=\frac{\theta }{A},\stackrel{~}{\varphi }=\frac{\varphi }{A},\stackrel{~}{\psi }=\frac{\psi }{A},$$
(18)
where $`A=\sqrt{\lambda }`$ as it was given in (13). We note here that originally the fields had no renormalization in either the original or in the dual model, consistent with the fact that the Euler angles are compact variables. The rescaling nevertheless introduces nontrivial renormalization, but this can be deduced from the renormalization of $`\lambda `$.
Deleting the tilde from the rescaled fields, we get the following result:
| $`\overline{}=`$ | $`{\displaystyle \frac{1}{2}}\{(_\mu \theta )^2+(1+r\mathrm{cos}^2(A\theta ))(_\mu \varphi +{\displaystyle \frac{h}{A}}\delta _{0\mu })^2+`$ |
| --- | --- |
| | $`+(_\mu \psi )^2+2a\omega _{\mu \nu }\mathrm{cos}(A\theta )(_\mu \varphi +{\displaystyle \frac{h}{A}}\delta _{0\mu })_\nu \psi \}.`$ |
(19)
According to the general prescription of Wick rotation, we define the Euclidian continuation of our model by $`\overline{}_E=\overline{}_M(_0i_0,_1_1)`$. Note that a similar continuation of $`hih`$ would lead to inconsistencies (like wrong signs in the propagators: $`\frac{1}{p_E^2h^2}`$ instead of$`\frac{1}{p_E^2+h^2}`$).
Before starting the perturbative expansion of the trigonometric functions we make another field transformation (redefinition): $`\theta \frac{\pi }{2\lambda }+\theta `$, guaranteeing that the minimum we expand around is: $`(\theta =0,\varphi =0,\psi =0)`$. Thus the Euclidian Lagrangian we use has the form:
| $`\overline{}=`$ | $`{\displaystyle \frac{1}{2}}\{(_\mu \theta )^2+(1+r\mathrm{sin}^2(A\theta ))(_\mu \varphi +{\displaystyle \frac{h}{A}}\delta _{0\mu })^2+`$ |
| --- | --- |
| | $`+(_\mu \psi )^2+2a\mathrm{sin}(A\theta )(\omega _{\mu \nu }_\mu \varphi _\nu \psi +{\displaystyle \frac{h}{A}}\mathrm{\Omega }_{0\nu })_\nu \psi \},`$ |
(20)
where $`r`$ and $`a`$ are given by (16) resp. (14), and the new parameters $`\omega `$ and $`\mathrm{\Omega }`$, which bear the model dependence, are: $`\omega _{\mu \nu }=\delta _{\mu \nu }`$, $`\mathrm{\Omega }_\nu =i\delta _{0\nu }`$ in the original and $`\omega _{\mu \nu }=iϵ_{\mu \nu }`$, $`\mathrm{\Omega }_\nu =ϵ_{0\nu }`$ respectively in the dual model.
Next we expand the Lagrangian $`\overline{}`$ around the classical ground state: $`(\theta =0,\varphi =0,\psi =0)`$, with $`A=\sqrt{\lambda }`$ being the relevant coupling constant. We note at this point that one could follow a different route and expand the parameters $`r`$, $`g`$ and $`a`$ in terms of the RG invariant. The motivation for this would be that at the end of the computation this has to be done anyway. This possibility and the complications that arise will be investigated elsewhere .
After some algebra the result is as follows (for the sake of simplicity we denote $`\overline{}`$ by $``$):
$$=_2A^2+_1A^1+_0++_1A+_2A^2+𝚘(A^3),$$
(21)
where
| $`_2=`$ | $`{\displaystyle \frac{1}{2}}h^2,_1=ih_0\varphi ,`$ |
| --- | --- |
| $`_0=`$ | $`{\displaystyle \frac{1}{2}}(_\mu \psi )^2+{\displaystyle \frac{1}{2}}(_\mu \varphi )^2+{\displaystyle \frac{1}{2}}(_\mu \theta )^2r{\displaystyle \frac{1}{2}}h^2\theta ^2+ah\mathrm{\Omega }_\nu \theta _\nu \psi ,`$ |
| $`_1=`$ | $`irh\theta ^2_0\psi +a\omega _{\mu \nu }\theta _\mu \varphi _\nu \psi ,`$ |
| $`_2=`$ | $`{\displaystyle \frac{1}{2}}r\theta ^2(_\mu \varphi )^2+{\displaystyle \frac{1}{6}}rh^2\theta ^4{\displaystyle \frac{1}{6}}ah\mathrm{\Omega }_\nu \theta ^3_\nu \psi .`$ |
(22)
Notice that $`_2`$ is a constant (i.e., it is independent of the fields), while $`_1(h,\phi )=ih_0\varphi `$ is a total derivative, thus it can be discarded in this non-topological sector of the theory. From $`_0`$ we see that $`\varphi `$ is a massless scalar field, while $`\theta `$ and $`\psi `$ are mixed, apart of the mixing the former is a massive scalar field with mass $`\sqrt{r}h`$, the latter is massless. The interaction of the different fields is highly non-trivial, as can be seen above, and contains infinitely many vertices. Nevertheless, these vertices are naturally separated in the weak coupling regime. We emphasize that only the first two terms, $`_2`$ and $`_1`$, are common to both models, as the model dependent parameters $`\alpha _i`$, $`\beta `$, $`\omega `$ and $`\mathrm{\Omega }`$ appear in $`_j`$ for all $`j0`$.
## 4 The free energy
Our goal is to define and compute the free energy density in perturbation theory. After setting the stage we do the explicit computations.
### 4.1 Definition of the free energy
At this point our aim is to define the free energy density perturbatively. Denoting the fields collectively by $`\phi =(\theta ,\varphi ,\psi )`$, the free energy (density) reads:
$$e^{(h)V}=\frac{𝒟\phi e^{S[h,\phi ]}}{𝒟\phi e^{S[h=0,\phi ]}},$$
(23)
where $`V`$ is the volume of the system, $`S[h,\phi ]=d^2x(h,\phi (x))`$, and $`𝒟\phi `$ denotes the functional integration over the field configurations $`\phi =(\theta ,\varphi ,\psi )`$: $`𝒟\phi =𝒟\varphi 𝒟\psi 𝒟\theta `$. The role of the denominator in (23) is to insure the correct normalization: $`(h=0)=0`$. From dimensional arguments one expects the following functional dependence: $`(h)=h^2\mathrm{\Psi }(h)`$, where $`\mathrm{\Psi }=\mathrm{\Psi }(h)`$ is a dimensionless function.
Let us note the similarity between the free energy defined above and quantum effective action (the generator of the 1PI graphs). The role of the external field is played by $`h`$, that couples to a conserved charge (composite operator) rather than an elementary field. This similarity will play a structurally simplifying role when we discuss the renormalization of the model.
For the perturbative expansion, in view of eqs. (21,23), it proves useful to introduce
$$S[h,\phi ]=\underset{i=2}{\overset{\mathrm{}}{}}S_i[h,\phi ]A^i,$$
(24)
with $`S_i[h,\phi ]=d^2x_i(h,\phi (x))`$, where $`i2`$. Using this in eq. (24) we obtain a similar expression for the free energy:
$$(h)=\underset{i=2}{\overset{\mathrm{}}{}}_i(h)A^i.$$
(25)
Our task will be to determine the first few terms in this expansion, in both models, and compare them. More precisely we determine the first six terms of $`(h)`$ and check whether they are equal.
As $`_2`$ is independent of the fields $`\phi `$, it results that $`S_2[h,\phi ]=\frac{1}{2}h^2V`$, and $`\mathrm{exp}(S_2[h,\phi ])`$ factorizes (we will come back to this) in the functional integral (23). This way one readily obtains that the first term of (25) is
$$_2(h)=\frac{1}{2}h^2.$$
(26)
Of course this is valid both in the original and dual model, implying that at leading order the perturbative free energy densities coincide. In addition, since $`_1`$ is a total derivative, it implies that $`S_1[h,\phi ]=0`$, thus
$$_1(h)=0;$$
(27)
again a model independent statement.
Thus all what remains to be dealt with is the reduced action
$$\overline{S}[h,\phi ]=\underset{i=0}{\overset{\mathrm{}}{}}S_i[h,\phi ]A^i,$$
(28)
and the reduced free energy
$$\overline{}(h)=\underset{i=0}{\overset{\mathrm{}}{}}_i(h)A^i,(h)=_2(h)+\overline{}(h).$$
(29)
Eq. (27) might also suggest that all $`_i`$ with $`i`$ odd vanishes, this would comply with the fact that the coupling $`A=\sqrt{\lambda }`$ is just an artifact of our perturbation theory, as it was $`\lambda `$ that appeared in the original Lagrangian. The vanishing of all odd power contributions in $`A`$ would indeed imply that the true coupling is in fact $`\lambda `$.
Introducing the $`Z(h)=𝒟\phi e^{\overline{S}[h,\phi ]}`$ auxiliary function then we can rewrite eq. (23) as $`e^{\overline{}(h)V}=Z(h)/Z(h=0)`$. Moreover, if $`M`$ is an operator we define the following ’expectation value’:
$$M=\frac{𝒟\phi e^{S_0[h,\phi ]}M}{𝒟\phi e^{S_0[h,\phi ]}}.$$
(30)
Expanding $`Z(h)`$ as a power series in $`A`$ we have:
| $`Z(h)=`$ | $`[\mathrm{\hspace{0.17em}1}S_1AS_2{\displaystyle \frac{1}{2}}S_1^{\mathrm{\hspace{0.17em}2}}A^2`$ |
| --- | --- |
| | $`S_3S_1S_2{\displaystyle \frac{1}{6}}S_1^{\mathrm{\hspace{0.17em}3}}A^3+\vartheta (A^4)]{\displaystyle }𝒟\phi e^{\overline{S}[h,\phi ]}.`$ |
(31)
For simplicity we have omitted to write the functional dependence of the $`\overline{S}_i[h,\phi ]`$-s. Using the identity
$$1+y_1A+y_2A^2+y_3A^3+𝚘(A^4)=e^{x_1A+x_2A^2+x_3A^3+𝚘(A^4)},$$
(32)
where
$$x_1=y_1,x_2=y_2\frac{1}{2}y_1^{\mathrm{\hspace{0.17em}2}},x_3=y_3y_1y_2+\frac{1}{3}y_1^{\mathrm{\hspace{0.17em}3}},$$
we can read off the various components of the reduced free energy density
$$e^{_0(h)V}=\frac{𝒟\phi e^{S_0[h,\phi ]}}{𝒟\phi e^{S_0[h=0,\phi ]}},$$
(33)
$$_1(h)=\frac{1}{V}S_1,$$
(34)
$$_2(h)=\frac{1}{V}\left[S_2\frac{1}{2}S_1^{\mathrm{\hspace{0.17em}2}}+\frac{1}{2}S_1^2\right].$$
(35)
We emphasize that in the above formulas we kept only those terms which depend on $`h`$. In other words we discarded the contribution of $`Z(h=0)`$, which in fact is a divergent quantity. This is consistent with the fact the that only the derivatives of the free energy are observable thermodynamical quantities.
### 4.2 Propagators
To compute the various vacuum expectation values (or correlators since we are in Euclidean space) determining $`_i(h)`$ we use dimensional regularization (with $`n=2+ϵ`$). Since
$$S_0(h)=d^2x_0(h)=\frac{1}{2}d^2x\phi ^t(x)M(x)\phi (x)$$
(36)
where
| $`\phi ^t(x)=(\theta (x),\varphi (x),\psi (x)),`$ |
| --- |
| $`M(x)=\left(\begin{array}{ccc}^2rh^2& 0& ah\mathrm{\Omega }_\nu _\nu \\ 0& ^2& 0\\ ah\mathrm{\Omega }_\nu _\nu & 0& ^2\end{array}\right),^2=_\mu _\mu ,`$ |
(37)
in order to determine the propagators of the various fields we have to invert the matrix operator $`M(x)`$. This is easily done in momentum space resulting:
$$\begin{array}{cc}G_\theta (x)=\frac{d^np}{(2\pi )^n}e^{ipx}p^2g(p),\hfill & G_\psi (x)=\frac{d^np}{(2\pi )^n}e^{ipx}(p^2rh^2)g(p),\hfill \\ G_{\theta \psi }(x)=iah\frac{d^np}{(2\pi )^n}e^{ipx}\mathrm{\Omega }pg(p),\hfill & G_{\varphi \varphi }(x)=\frac{d^np}{(2\pi )^n}e^{ipx}\frac{1}{p^2},\hfill \\ G_{\theta \varphi }(x)=0,\hfill & G_{\psi \varphi }(x)=0,\hfill \end{array}$$
(38)
where
$$g(p)=\frac{1}{p^4h^2[rp^2+(a\mathrm{\Omega }p)^2]},\mathrm{\Omega }p\mathrm{\Omega }_\nu p_\nu .$$
(39)
We have used commonly the notation $`G_{\phi \phi ^{}}(xy)=\phi (x)\phi ^{}(y)`$, with $`G_\phi G_{\phi \phi }`$. Just note in passing a few simple properties: $`G_{\theta \theta }(x)`$ is even, while $`G_{\theta \psi }(x)`$ and $`G_{\psi \psi }(x)`$ are odd, and $`G_{\psi \theta }(x)=G_{\theta \psi }(x)`$. It is also obvious that $`\varphi `$ behaves like a massless scalar.
### 4.3 Computation of $`_0(h)`$
The computation of $`_0(h)`$ involves in fact the evaluation of a functional determinant, similarly to the the case of the quantum effective action. Evaluating the Gaussian integral from eq. (33), using the identity $`detX=e^{\mathrm{Tr}\mathrm{ln}X}`$, results:
$$_0(h)=\frac{1}{2}\frac{d^np}{(2\pi )^n}\mathrm{ln}\left(1\frac{h^2(rp^2+(a\mathrm{\Omega }p)^2)}{p^4}\right).$$
(40)
A proper way to compute this expression is to take its derivative with respect to $`h`$ and solve the following initial value problem:
$$\frac{d_0(h)}{dh}=h\frac{d^np}{(2\pi )^n}\frac{rp^2+(a\mathrm{\Omega }p)^2}{p^4h^2(rp^2+(a\mathrm{\Omega }p)^2)},_0(0)=0.$$
(41)
Rescaling $`php`$, the $`h`$ dependence factorizes and we get:
$$_0(h)=\frac{h^n}{n}\frac{d^np}{(2\pi )^n}\frac{r_0p^2+(a_0\mathrm{\Omega }p)^2}{p^4h^2[r_0p^2+(a_0\mathrm{\Omega }p)^2]}.$$
(42)
Above we have made it explicit that the integral is computed in terms of the bare quantities $`r_0`$ and $`a_0`$, rather then the renormalized ones, $`r`$ and $`a`$. From now on we are going to make this distinction clear in all subsequent formulas.
This same expression was obtained in , though the initial Lagrangian differed from the one used here by a certain rescaling of the fields. A rescaling usually cannot cause major discrepancies, at least at low orders of perturbation theory, as this example also reflects.
The model dependence is manifest in (42). The integral is divergent in two dimensions more precisely it has a first order pole in $`ϵ`$. The analysis of has computed the pole term and the constant term in the $`ϵ`$ expansion of (42) in closed form, in terms of generalized hyper-geometric functions. After renormalization, the two expressions – as functions of the original (respectively dual) coupling and parameter – were not equal, but, the difference could be accounted for by the scheme dependence of the two loop beta functions. Indeed it was pointed out in that the equivalence of the original and dual $`\beta `$ functions corresponds to a perturbative redefinition of the coupling constants in the dual model:
$$\stackrel{~}{\lambda }=\lambda +\frac{\lambda ^2}{4\pi }(1+g),\stackrel{~}{g}=g+\frac{\lambda }{4\pi }(1+g)^2.$$
(43)
Implementing this redefinition in the expressions of the renormalized one loop free energy densities revealed their equality.
Here we take a different route from the one described above. Our strategy is to express the parameters that bear the model dependence in terms of the RG invariant in both models, then set these two RG invariants equal and compare the results. Evaluating $`r_0p^2(a_0\mathrm{\Omega }p)^2`$ we get $`p_0^2g_0p_1^2`$ in the original model and $`\stackrel{~}{p}_1^2g_0\stackrel{~}{p}_0^2`$ respectively in the dual model, where $`\stackrel{~}{p}_\mu =ϵ_{\mu \nu }p_\nu `$. Due to we have $`\stackrel{~}{p}^2=p^2`$, and one can perform a change of variables from $`p`$ to $`\stackrel{~}{p}`$. This way, changing also $`p_1p_0`$, one can obtain formally identical expressions in the two cases. At this point we only remark that the role of $`p_0`$ and $`p_1`$ in dimensional regularization is different.
### 4.4 Computation of $`_1(h)`$
According to (34) and (22) the computation of $`_1(h)`$ involves the following correlation functions: $`\theta (x)_\nu \psi (x)`$ and $`\theta ^2(x)_0\varphi (x)`$. Using Wick’s theorem we obtain:
$$_1(h)=0.$$
(44)
Eq. (44) has twofold meaning. It shows once again the model independence of the free energy, though we have already encountered explicit model dependence. On the other hand supports our earlier statement about the vanishing of the non-analytic corrections in $`\lambda `$.
### 4.5 Computation of $`_2(h)`$
In the case of $`_2(h)`$ we will not be able to obtain the result in a closed form, nevertheless what we can actually compute will suffice to achieve our goals. According to eq. (35) we have to compute $`S_2(h)`$, $`S_1(h)^2`$ and $`S_1(h)^{\mathrm{\hspace{0.17em}2}}`$. We can immediately quote eq. (44) and $`S_1(h)=0`$.
#### 4.5.1 Computation of $`S_2(h)`$
From eq. (22) we see that
| $`S_2(h)={\displaystyle d^nx}`$ | $`[{\displaystyle \frac{1}{2}}r_0\theta ^2(x)(_\mu \varphi )^2(x)+{\displaystyle \frac{1}{6}}r_0h^2\theta ^4(x)`$ |
| --- | --- |
| | $`{\displaystyle \frac{1}{6}}a_0h\mathrm{\Omega }_\nu \theta ^3_\nu \psi (x)].`$ |
(45)
At first sight one might want to discard the terms that are not coupled to $`h`$. Nevertheless these terms acquire $`h`$ dependence through the $`h`$-dependent propagators.
Based on Wick’s theorem for the first term we have $`\theta ^2(x)(_\mu \varphi )^2(x)=\theta ^2(x)(_\mu \varphi )^2(x)`$. Since the $`\varphi `$ propagator is in fact a Green’s function, or in other words fundamental solutions of the corresponding wave equation, we can conclude that $`(_\mu \varphi )^2(x)=\delta ^{(n)}(0)`$ (modulo equal time commutator terms), where $`\delta ^{(n)}(0)`$ is the Dirac delta distribution ’evaluated at 0’. But in dimensional regularization the latter is set to zero, yielding no contribution.
The next term in eq. (45) is also readily evaluated: $`\theta ^4(x)=3\theta ^2(x)^2=3G_\theta (0)^2`$. Evaluating the third term gives (recall that the system in finite volume: $`d^nx=V`$):
$$S_2(h)=\frac{1}{2}h^2V\left[r_0+\frac{1}{n}a_0^{\mathrm{\hspace{0.17em}2}}\mathrm{\Omega }_\nu \mathrm{\Omega }_\nu \right]G_\theta (0)^2.$$
(46)
#### 4.5.2 Computation of $`S_1(h)^2`$
The computation of $`S_1(h)^2`$ involves in fact the evaluation of double integral $`d^nxd^ny(x)(y)`$. From eq. (22) it turns out that $`S_1(h)^2`$ equals $`d^nx`$ times
| $`r_0^2h^2_0^{\mathrm{\hspace{0.17em}2}}G_\varphi (x)[G_\theta (0)^2+2G_\theta (x)^2]4ir_0ha_0\omega _{\mu \nu }_0_\mu G_\varphi (x)G_\theta (x)_\nu G_{\theta \varphi }(x)+`$ |
| --- |
| $`+a_0^{\mathrm{\hspace{0.17em}2}}\omega _{\mu \nu }\omega _{\lambda \rho }_\mu _\lambda G_\varphi (x)[_\nu G_{\theta \varphi }(x)_\rho G_{\theta \varphi }(x)G_\theta (x)_\nu _\rho G_\psi (x)+`$ |
| $`+_\nu G_{\theta \varphi }(0)_\rho G_{\theta \varphi }(0)].`$ |
(47)
A priori it is not clear at all how to obtain the overall $`h^2`$ factor required by dimensional analysis. Moreover it is also puzzling how the explicit factors of $`i`$ disappear during the computation. As we shall see shortly it is the form of the ‘tensor’ parameters $`\mathrm{\Omega }_\nu `$, $`\omega _{\mu \nu }`$ and propagators that is responsible for the correct answers.
The first term in (47) results
$$2r_0^{\mathrm{\hspace{0.17em}2}}h^2V\frac{1}{n}G_\theta (0)^2\delta _{00}.$$
(48)
The only non-trivial fact that one has to use an IR regularization of the field $`\varphi `$, the regulator’s mass we denote by $`m`$. Then the first term in this parenthesis will be proportional to
$$d^nx_0^{\mathrm{\hspace{0.17em}2}}G_\varphi (x)d^nk\frac{k_0^{\mathrm{\hspace{0.17em}2}}}{k^2+m^2}\delta ^{(n)}(k)=0.$$
(49)
The second term in (47) results
$$4r_0a_0^{\mathrm{\hspace{0.17em}2}}h^2Vi\omega _{\mu \nu }\mathrm{\Omega }_\rho \frac{d^nk_1}{(2\pi )^n}\frac{d^nk_2}{(2\pi )^n}\frac{(k_1+k_2)_0(k_1+k_2)_\mu }{(k_1+k_2)^2}k_1^{\mathrm{\hspace{0.17em}2}}g(k_1)k_{2\nu }k_{2\rho }g(k_2),$$
(50)
while the last terms of (47) give
| $`a_0^{\mathrm{\hspace{0.17em}2}}\omega _{\mu \nu }\omega _{\lambda \rho }V{\displaystyle \frac{d^nk_1}{(2\pi )^n}\frac{d^nk_2}{(2\pi )^n}}`$ | $`{\displaystyle \frac{(k_1+k_2)_\lambda (k_1+k_2)_\mu }{(k_1+k_2)^2}}g(k_1)g(k_2)k_{2\rho }`$ |
| --- | --- |
| | $`[a_0^{\mathrm{\hspace{0.17em}2}}h^2k_{1\nu }(k_1\mathrm{\Omega })(k_2\mathrm{\Omega })+k_1^{\mathrm{\hspace{0.17em}2}}k_{2\nu }(k_2^{\mathrm{\hspace{0.17em}2}}r_0h^2)].`$ |
(51)
The general structure of the integrals that appear (with one exception) is of the following form:
$$\frac{d^nk_1}{(2\pi )^n}\frac{d^nk_2}{(2\pi )^n}\frac{(k_1+k_2)_{\mu _1}(k_1+k_2)_{\mu _2}}{(k_1+k_2)^2}k_{1_{\mu _3}}k_{1_{\mu _4}}k_{2_{\mu _5}}k_{2_{\mu _6}}g(k_1)g(k_2).$$
(52)
The only integral that cannot be brought to this form has $`k_2^2g(k_2)`$ instead of $`g(k_2)`$. If $`g(k)`$ were a covariant expression in $`k`$, then the value of the integral would be given completely by the index structure. More precisely covariance would require the result to be the sum of triple products of the Kronecker delta functions. The number of independent possibilities is $`6!/(2!)^3/3!=15`$. Exploiting the obvious symmetries under the exchange of $`\mu _1\mu _2`$, $`\mu _3\mu _4`$, $`\mu _5\mu _6`$, and the less obvious exchange $`(\mu _3,\mu _4)(\mu _5,\mu _6)`$, which can be seen by a change of integration variables, the tensor structure reduces significantly to only four independent terms. Denoting $`\delta _{\mu _i\mu _j}`$ by $`(ij)`$, we get for (52)
| $`(12)(34)(56)I_1+(12)[(35)(46)+(36)(45)]I_2+(34)[(15)(26)+(16)(25)]I_3+`$ |
| --- |
| $`+(56)[(13)(24)+(14)(23)]I_3+\{(13)[(25)(46)+(26)(45)]+(14)[(25)(36)+`$ |
| $`+(26)(35)]+(15)[(24)(36)+(23)(46)]+(16)[(23)(45)+(24)(35)]\}I_4.`$ |
(53)
Using the standard methodology one can compute the unknowns $`I_1`$ through $`I_4`$, in a straightforward manner. Unfortunately $`g(k)`$ is not covariant with respect to the full $`SO(n)`$. As we pointed out, covariance was already broken at the level of the Lagrangian. If we analytically continue $`p_0`$ to $`n_0`$ dimensions and $`p_1`$ to $`n_1`$ dimensions, with $`n_0+n_1=n`$, then instead of the full $`SO(n)`$ group we get $`SO(n_0)\times SO(n_1)`$. In other words, among the $`p_0`$-s and $`p_1`$-s standard covariance arguments remain valid, and the computation sketched above makes sense. Thus if all the indices are of $`p_0`$ or $`p_1`$ type, then we can reliably compute the integrals.
In the expression of the two-loop free energy (50, 51) the free indices we have dealt above are contracted with different tensor structures. In the original model $`\mathrm{\Omega }_\nu `$ and $`\omega _{\mu \nu }`$ involve only delta functions, and this way what we said applies. On the other hand for the dual model $`\mathrm{\Omega }_\nu `$ and $`\omega _{\mu \nu }`$ are epsilon tensors, mixing the indices, and the arguments presented break down. The same is true for the sum of the terms involved, that we eventually want to compute.
There is an independent argument that shows that even if we were able to do the integrals, the result would not be reliable, due to the behavior of the $`ϵ`$ tensor in dimensional regularization. More precisely, even if the integrals in (50) and (51) were covariant, based on the definition of $`\omega _{\mu \nu }`$ and on (53), we could conclude that in the dual model $`S_1(h)^2`$, and as a result $`_2(h)`$, contains terms proportional to the product of two $`ϵ`$ tensors, with uncontracted indices. The broken covariance by the external field $`h`$, invalidates the above argument, but is highly likely that it complicates matters, rather than simplifies, and as a result we would still end up with $`ϵ_{\mu \nu }ϵ_{\alpha \beta }`$ terms. As it is well known, there is no consistent way defining such on object in dimensional regularization .
At this point the computation of $`S_1(h)^2`$ is hopeless, driving the same statement about $`_2(h)`$. Nevertheless we can do something less ambitious. Following , in order to compare the results, at the very end of the computation we want to express them in terms of the two RG invariants, which are finally set equal. In other words, based on (11) and (16) we trade the renormalized quantities $`g`$ and $`r`$ for the RG invariant $`M`$, implicitly meaning that by this we have also set the RG invariants equal. Since it is the bare $`a_0`$ that appears in (50) and (51), and also in (42), let’s investigate more closely what happens to this term during renormalization and expansion in terms of the RG invariant.
At the level of the bare quantities we aim to have $`a_0=\sqrt{1+g_0}`$. Based on this is certainly true at two loop level in the original model, while in the dual one this is more subtle. As it was shown in naively the relation $`a_0=\sqrt{1+g_0}`$ cannot be maintained at two loop level in the dual model. Nevertheless taking into account the redefinition of the dual model’s coupling and parameter, eq.(43), the above relation can be maintained at one loop order.
As $`a_0=\sqrt{1+g_0}`$ is doubtlessly valid in the leading order, the renormalization (5) amounts to $`a_0=\sqrt{1+Z_g(\lambda ,g)g}`$. But (6) shows that $`Z_g(\lambda ,g)=1+𝚘(\lambda )`$, thus to leading order we have $`a_0=\sqrt{1+g}`$, with $`g`$ the renormalized coupling. On the other hand at this point we can use (11) and conclude that, after renormalization and expanding in the RG invariant, $`a_0`$ becomes proportional to $`\sqrt{\lambda }`$.
The good news is that all the terms we were unable to compute (50) and (51) are proportional to $`a^2`$, hence are of order $`\lambda ^2`$. Since we have not even attempted to compute the $`_4(h)`$ term that is of the same order, we neglect them for the time being. Having in mind the insertion of the RG expressions for the parameters, we simply get:
$$_2(h)=h^2\frac{n2}{2n}G_\theta (0)^2.$$
(54)
Naturally this result is to be interpreted as modulo terms that will be of higher order after renormalization and expansion in the RG invariant.
## 5 Renormalization
The results from the previous section are divergent, and we used dimensional regularization to compute them. As advertised we make use of the similarities between the free energy and the quantum effective action to discuss the renormalization of the former, modeled by the renormalization of the latter (see e.g. ). In this section we follow a more or less naive renormalization procedure, and compute the renormalization group improved perturbative two-loop free energy. It will be the role of the next section to tight the loose end, and prove that what we did is indeed correct.
Since we are at second order of perturbation theory, we have both first and second order poles in dimensional regularization, as can be seen in eq. (54). Since the free energy is a physical quantity, it has to be well defined after renormalization. The recipe of this section is to use the renormalization of the deformed principal $`\sigma `$ model (5),(6), and the renormalization of its dual, to renormalize the free energy. The above renormalizations were performed using the geometric method of .
This procedure can be immediately objected since the free energy is computed in a theory that has additional terms in the Lagrangian (19), compared to the deformed principal $`\sigma `$ model (4) and its dual (8). A priori there is no reason to expect that the wavefunction and coupling constant renormalization functions are the same, though in fact they are. In the following we present a simple argument in favor of the above statement. We shall give a complete proof in section 5.5 below.
Coupling the external field amounts to the appearance of terms proportional to $`h`$. Setting $`h=0`$ we get the original models (deformed principal $`\sigma `$ model and its dual). Thus the wavefunction and coupling constant renormalization functions can at most differ from the ones of the original models by terms proportional to $`h`$. But $`h`$ is a dimensionful quantity, having mass dimension $`+1`$ (it is a super-renormalizable coupling or relevant perturbation). On the other hand our theories have no other dimensionful quantities, and the wavefunction and coupling constant renormalization functions are dimensionless. We conclude that $`h`$ cannot appear in the latter ones, proving the assertion.
In addition since $`h`$ couples to a conserved charge, it is not renormalized. The only issue that remains will be to deal with the renormalization of the operators that couple to $`h`$, viewed as composite operators. We postpone this to section 5.5.
In order to cancel the second order poles of the regularized free energy (54) we have to go beyond the computation of . Let us review what we know at this stage about the renormalization of the deformed principal $`\sigma `$ model and its dual. In (15) we have introduced a generalized Lagrangian, and argued in the last paragraph of the section that it encompasses the renormalization properties of both models. Thus we can translate the renormalization properties of the original models to those of the generalized Lagrangian. The Euler angles ($`\varphi `$,$`\theta `$,$`\psi `$) are compact and have no wavefunction renormalization. The renormalization of $`r,\lambda `$ and $`a`$ in the unified model follows from those of $`g`$ (resp. $`a`$) and $`\lambda `$ in the original models.
In (18) we rescaled the fields ($`\varphi `$,$`\theta `$,$`\psi `$), this way in (22) they have the corresponding nontrivial wavefunction renormalization. Thus $`Z_\lambda (\lambda ,M)`$, the coupling constant renormalization of $`\lambda `$, is the central object for their renormalization. (Recall that $`M`$ is the renormalization group invariant that appeared in (10)). $`Z_\lambda (\lambda ,M)`$ was computed at two loop order in perturbation theory in . The interest in was restricted to the single pole terms. On the other hand we are constrained to deal with the second order poles too. First we compute the residue of the second order pole in $`Z_\lambda (\lambda ,M)`$.
### 5.1 $`\lambda `$ at two loop
The goal is to determine the terms $`y_\lambda (\lambda ,M)`$ and $`\overline{y}_\lambda (\lambda ,M)`$ in the following expansion:
$$Z_\lambda (\lambda ,M)=1+\frac{1}{ϵ}y_\lambda (\lambda ,M)+\frac{1}{ϵ^2}\overline{y}_\lambda (\lambda ,M)+𝚘(\frac{1}{ϵ^3}).$$
(55)
The first term was basically determined in . All we need is to use (6) for $`Z_\lambda (\lambda ,g)`$ and (14) for $`g`$, in terms of the renormalization group invariant $`M`$, to obtain:
$$y_\lambda (\lambda ,M)=\frac{1}{2\pi }\lambda +\frac{1}{8\pi }(2\alpha _1^2\frac{1}{\pi })\lambda ^2.$$
(56)
Our task of computing $`\overline{y}_\lambda (\lambda ,M)`$ is highly simplified by the special properties of the non-linear $`\sigma `$ models. Following it was shown in that the generalized renormalization theory of non-linear $`\sigma `$ models lead to generalized renormalization group equations, that allow one to determine the residues of the higher order poles in a given coupling constant renormalization function like $`Z_\lambda (\lambda ,M)`$, without extra diagrammatic computations.
More precisely it is shown that in a theory (like ours) with a single dimensionless coupling constant $`\lambda `$, with mass scale parameter $`\mu `$, defined in $`n`$ dimensions (dimensionally regularized and minimally subtracted), having an expansion of the bare coupling $`\lambda _0`$ in terms of the renormalized coupling $`\lambda `$ of the form
$$\lambda _0=\mu ^{2n}\left(\lambda +\underset{\nu =1}{\overset{\mathrm{}}{}}\frac{a_\nu (\lambda )}{(n2)^\nu }\right),$$
(57)
the pole residues $`a_\nu (\lambda )`$ satisfy the recursive pole equations:
$$(1\lambda \frac{}{\lambda })a_{\nu +1}(\lambda )=(1\lambda \frac{}{\lambda })a_1(\lambda )\frac{}{\lambda }a_\nu (\lambda ).$$
(58)
We can apply this directly to compute $`\overline{y}_\lambda (\lambda ,M)`$. From the expressions following (5) we see that in our case $`a_1(\lambda )=\lambda y_\lambda (\lambda ,M)`$, with $`y_\lambda (\lambda ,M)`$ given in (56). Using (58) for $`\nu =1`$ we can determine the leading term in the expansion of $`a_2(\lambda )`$. A simple calculation shows that $`a_2(\lambda )=\frac{1}{4\pi ^2}\lambda ^3+\mathrm{}`$, implying
$$\overline{y}_\lambda (\lambda ,M)=\frac{1}{4\pi ^2}\lambda ^2+\mathrm{}.$$
(59)
### 5.2 One loop free energy
By now it is a computation on the back of an envelope to obtain the renormalized one loop free energy density. We have to use (26) for the leading term (with $`\lambda `$ as the bare coupling $`\lambda _0`$) and (42) for the one loop regularized result. The integral itself, as it was pointed out in , is hard to deal with exactly. Nevertheless, with the parameters traded for the RG invariants, and keeping only the leading contribution in $`\lambda `$ (as explained in the last paragraph of Section 4), and renormalizing the expression using (5), (55) and (56), we obtain:
$$^{1loop}(h)=h^2(\frac{1}{2\lambda }\frac{1}{8\pi }[\mathrm{ln}\left(\frac{h}{\mu }\right)^2+\gamma 1\mathrm{ln}(4\pi )]).$$
(60)
Some of the terms that are higher order in $`\lambda `$, and are not yet displayed, will be used for the computation of the $`_2`$.
We note two things here. Firstly, the one loop equivalence might be argued to be not surprising based on the experience gained in and the general one-loop beta function result of . Secondly, the above expression correctly reproduces the known one-loop free energy of the $`O(3)`$ sigma model, that arises in the $`\alpha _1=0`$ limit.
As opposed to the corresponding computation in , the procedure that leads to (60) does not require the splitting of the momentum integration in the 0-th direction and the rest. The reason is simply that due to the expansion in the RG invariant, the order $`\lambda `$ term becomes totally covariant. The non-covariance is shifted to the next order.
### 5.3 Two loop free energy
For the computation of the two loop free energy we need all the results developed. The $`𝚘(\lambda )`$ corrections arising from the leading bare term (26) are obvious in the light of (5), (55), (56) and (59) . The $`𝚘(\lambda )`$ contributions from the bare next to leading term of the one loop free energy (42) can be computed as an expansion in the RG invariant, using (14) and (16). The results are:
$$\lambda h^2\frac{\alpha _1^{\mathrm{\hspace{0.17em}2}}}{2}\left[\frac{1}{4\pi ϵ}\frac{1}{8\pi }(\gamma \mathrm{ln}(4\pi )+2\mathrm{ln}(h)1)\right],$$
(61)
where $``$ equals $`\delta _{11}`$ in the original model, and $`ϵ_{0\nu }ϵ_{0\nu }`$ in the dual one. Although, as argued above, the product of two epsilon tensors with uncontracted indices is ambiguous in dimensional regularization, the contraction of one index gives a meaningful expression. Following , we assume that $`ϵ_{\mu \alpha }ϵ_{\nu \alpha }`$ has a consistent continuation, namely:
$$ϵ_{\mu \nu }=ϵ_{\nu \mu }ϵ_{\mu \alpha }ϵ_{\nu \alpha }=\delta _{\mu \nu }.$$
(62)
These two expressions might differ depending on the regularization schemes chosen. As we have pointed out in connection with (53), we can a priori continue the 0th direction into $`n_0`$ dimensions and the 1st direction into $`n_1`$ dimensions, provided $`n_0+n_1=n`$. We checked that the final result is consistent with dimensional analysis in either case. While we know that the choice $`n_0=1`$ is a consistent scheme with the continuation of the $`ϵ`$ tensor, we cannot claim the same about the general case.
As we will see in a moment the only discrepancy in the two-loop free energies of the two models comes from (61). This way we have several choices of different schemes to see the quantum equivalence. With the notation following (61) we can chose in both models the scheme with $`\delta _{00}=\delta _{11}`$. The second choice is two different schemes in the two models, but related to each other by: $`\delta _{00}^D=\delta _{11}^O`$. Besides these naive choices we have a highly non-trivial one, with the identical choice in both models: $`\delta _{00}=1`$ and $`\delta _{11}=n1`$. Unlike for the previously mentioned ones, we know the consistency of this scheme, and we are going to work with it in the rest of the paper. As we remarked in connection with (42), it was shown in that in this case the difference in (61) is accountable for a perturbative redefinition of the coupling constant $`\lambda `$ in the dual model. Naturally this redefinition has to proceed the expansion in the RG invariant. Unlike in , where this gave correction to $`^{1loop}(h)`$, due to the expansion in the RG invariant we get contribution to $`^{2loop}(h)`$. The redefinition of the coupling constant does not effect the genuinely higher order terms as it can be seen from (43).
Thus, as we just noted, we are going to use the last scheme in what follows: $`\delta _{00}=1`$ and $`\delta _{11}=n1`$, and for definiteness we consider the case of the original model first. In the light of what we just said, for the dual model, the results (63), (64) below, must be amended by $`\frac{h^2}{8\pi }\lambda \frac{1}{\sqrt{M}}`$, to account for the perturbative redefinition of $`\stackrel{~}{\lambda }`$, eq. (43). However once this redefinition is taken into account, the two loop results below are identical in the two models. The result for the other schemes will differ from the results to be presented by terms that are some number times $`\lambda h^2\frac{\alpha _1^{\mathrm{\hspace{0.17em}2}}}{8\pi }`$.
Using the bare result (54) and summing it with the corresponding contributions described, we obtain a cancelation of both the first and second order poles, and obtain a finite expression. The computation is somewhat tedious, but straightforward, the result is simply:
| $`^{2loop}(h)={\displaystyle \frac{h^2}{2\lambda }}{\displaystyle \frac{h^2}{8\pi }}[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^2+\gamma 1\mathrm{ln}(4\pi )]`$ |
| --- |
| $`{\displaystyle \frac{h^2\lambda }{16\pi ^2}}\{[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^2+\gamma {\displaystyle \frac{1}{2}}\mathrm{ln}(4\pi )]\pi \alpha _1^{\mathrm{\hspace{0.17em}2}}[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^2+\gamma 1\mathrm{ln}(4\pi )]\}.`$ |
(63)
It simplifies a bit if we use instead of minimal subtraction (MS scheme) the $`\overline{MS}`$ scheme: $`\mathrm{ln}\mu \mathrm{ln}\mu +(\gamma \mathrm{ln}(4\pi ))/2`$. Then the free energy reads:
| $`^{2loop}(h)=`$ | $`{\displaystyle \frac{h^2}{2\lambda }}{\displaystyle \frac{h^2}{8\pi }}[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^21]`$ |
| --- | --- |
| | $`{\displaystyle \frac{h^2\lambda }{16\pi ^2}}\{[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^2{\displaystyle \frac{1}{2}}]\pi \alpha _1^{\mathrm{\hspace{0.17em}2}}[\mathrm{ln}\left({\displaystyle \frac{h}{\mu }}\right)^21]\}.`$ |
(64)
### 5.4 Improvement of the perturbation theory
We can take advantage of the asymptotic freedom of our theory and calculate the RG improvement of the perturbative result. This gives the asymptotic expansion of the free energy for large values of the external fields.
Physical quantities depend on the renormalized coupling $`\lambda `$, the renormalized parameter $`g`$, and the dimensionful scale parameter (or subtraction point) $`\mu `$, in such a way that the action of the renormalization group (RG) operator
$$𝒟=\mu \frac{}{\mu }+\beta _\lambda (\lambda ,g)\frac{}{\lambda }+\beta _g(\lambda ,g)\frac{}{g},$$
(65)
vanishes on them. As the free energy density takes the form $`(h)=h^2\mathrm{\Psi }(\lambda ,g,\mu ,h)`$ and the external field $`h`$ is not renormalized, the function $`\mathrm{\Psi }`$ is renormalization invariant $`𝒟\mathrm{\Psi }=0`$. Since we are interested in the behaviour of the free energy density for large values of $`h`$ we write $`h=h_0e^x`$ where $`h_0`$ is fixed and let $`x\mathrm{}`$. Standard RG considerations then give that
$$(h)=h^2\mathrm{\Psi }(\lambda (x),g(x),\mu ,h_0),$$
(66)
where $`\lambda (x)`$ and $`g(x)`$ are the running coupling and parameter. In our final result for the two loop free energy density (64) the parameter $`g`$ is eliminated in favor of the two loop RG invariant. As a consequence, up to this order, $`\mathrm{\Psi }=\mathrm{\Psi }(\overline{\lambda }(x),\mu ,h_0)`$ depends only on $`\overline{\lambda }(x)`$, which is a solution of
$$\frac{d}{dx}\overline{\lambda }(x)=\beta _\lambda (\overline{\lambda }(x),\alpha _1),\overline{\lambda }(x=0)=\lambda ,$$
(67)
where the beta function has the following expansion:
$$\beta _\lambda (\lambda )=b_0\lambda ^2b_1\lambda ^3b_2\lambda ^4+\mathrm{}$$
(68)
with (see eq. (12))
$$b_0=\frac{1}{2\pi },b_1=\frac{1}{4\pi }\left(\frac{1}{\pi }\alpha _1^{\mathrm{\hspace{0.17em}2}}\right)=\frac{1}{4\pi ^2}\left(1\frac{p}{2}\right),p=\frac{2\pi }{\sqrt{M}}.$$
(69)
Thus we can go on with the RG analysis as if we had only one coupling constant $`\overline{\lambda }(x)`$. The expression of the (RG invariant) $`\mathrm{\Lambda }`$ parameter is
$$\frac{\mathrm{\Lambda }}{\mu }=e^{\frac{1}{b_0\lambda }}\lambda ^{\frac{b_1}{b_0^{\mathrm{\hspace{0.17em}2}}}}e^c\left[1+\left(\frac{b_1^{\mathrm{\hspace{0.17em}2}}}{b_0^{\mathrm{\hspace{0.17em}3}}}\frac{b_2}{b_0^{\mathrm{\hspace{0.17em}2}}}\right)\lambda +𝚘(\lambda ^2)\right],$$
(70)
or more conveniently
$$\mathrm{ln}\frac{\mathrm{\Lambda }}{\mu }=\frac{1}{b_0\lambda }\frac{b_1}{b_0^{\mathrm{\hspace{0.17em}2}}}\mathrm{ln}\lambda +c+\left(\frac{b_1^{\mathrm{\hspace{0.17em}2}}}{b_0^{\mathrm{\hspace{0.17em}3}}}\frac{b_2}{b_0^{\mathrm{\hspace{0.17em}2}}}\right)\lambda +𝚘(\lambda ^2),$$
(71)
where $`c`$ is a constant of integration. We define $`\mathrm{\Lambda }\mathrm{\Lambda }_{\overline{\mathrm{MS}}}`$ by choosing $`e^c=\left(\frac{1}{2\pi }\right)^{\frac{b_1}{b_0^2}}`$ as it simplifies some of the forthcoming expressions.
As the next step of the RG analysis we note, that using expression (71), it is customary to define an effective coupling $`\alpha (h)`$ by the following transcendental equation :
$$b_0\mathrm{ln}\frac{h}{\mathrm{\Lambda }_{\overline{\mathrm{MS}}}}=\frac{1}{\alpha }+\frac{b_1}{b_0}\mathrm{ln}\alpha .$$
(72)
The important property of this effective coupling is that it depends on the physical quantity $`s=\mathrm{ln}\left(\frac{h}{\mathrm{\Lambda }_{\overline{\mathrm{MS}}}}\right)`$, moreover one can express the running coupling $`\overline{\lambda }`$ in terms of the effective coupling $`\alpha `$ perturbatively:
$$\overline{\lambda }=\alpha (1+\xi _1\alpha +\xi _2\alpha ^2+\mathrm{}),$$
(73)
where
$$\xi _1=b_0\left(\mathrm{ln}\frac{h_0}{\mu }c\right),\xi _2=\xi _1^{\mathrm{\hspace{0.17em}2}}+\frac{b_1}{b_0}\xi _1\frac{b_1^{\mathrm{\hspace{0.17em}2}}}{b_0^{\mathrm{\hspace{0.17em}2}}}+\frac{b_2}{b_0}.$$
(74)
To use the $`𝚘(\lambda )`$ part of the free energy density effectively we need the large $`s`$ asymptotic expansion of the effective coupling up to $`𝚘(s^3)`$:
$$\alpha =\frac{2\pi }{s}\left(1+\frac{A\mathrm{ln}s}{s}+\frac{B}{s}+\frac{C\mathrm{ln}^2s}{s^2}+\frac{D\mathrm{ln}s}{s^2}+\frac{E}{s^2}\right),$$
(75)
where the coefficients $`A,\mathrm{}E`$ can be obtained from (72).
Now improving (64) by the RG, i.e. using eq. (73-75) in (64, 66), results, when the dust settles, the following formula for the asymptotic (large $`s`$) behaviour of the two loop free energy:
| $`^{2loop}(h)=`$ | $`{\displaystyle \frac{h^2}{4\pi }}(s+(1{\displaystyle \frac{p}{2}})\mathrm{ln}s{\displaystyle \frac{1}{2}}+(1{\displaystyle \frac{p}{2}})^2{\displaystyle \frac{\mathrm{ln}s}{s}}`$ |
| --- | --- |
| | $`+{\displaystyle \frac{1}{s}}[(1{\displaystyle \frac{p}{2}})^28\pi ^3b_2+{\displaystyle \frac{p1}{4}}]+𝚘({\displaystyle \frac{\mathrm{ln}^2s}{s^2}})).`$ |
(76)
Unfortunately $`b_2`$ has not yet been computed in the literature, nevertheless we can say a lot about it. Due to the expansion in the RG invariant, it will have terms coming from the lower order beta function coefficients, like in (6), and the genuine three loop coefficient evaluated at $`g=1`$. For simplicity we denote this last term $`b^{(3)}=b^{(3loop)}(g=1)`$. But this is exactly the three loop beta function coefficients of the $`O(3)`$ $`\sigma `$-model: $`b^{(3)}=b_3^{O(3)}`$. This beta function has been computed in , and in our convention for the coupling constant $`b^{(3)}`$ takes the form $`b^{(3)}=\frac{5}{32\pi ^3}`$. Carrying one step further the computation of (12) results:
$$b_2=\frac{1}{8\pi }\left(\alpha _1^{\mathrm{\hspace{0.17em}4}}\frac{3}{\pi }\alpha _1^{\mathrm{\hspace{0.17em}2}}\right)+b^{(3)}=\frac{1}{16\pi ^3}\left(\frac{p^2}{2}3p\right)+\frac{5}{32\pi ^3}.$$
(77)
Accordingly, as a final result we have:
$$^{2loop}(h)=\frac{h^2}{4\pi }\left(s+(1\frac{p}{2})\mathrm{ln}s\frac{1}{2}+(1\frac{p}{2})^2\frac{\mathrm{ln}s}{s}+\frac{1}{s}\frac{3p2}{4}+𝚘(\frac{\mathrm{ln}^2s}{s^2})\right).$$
(78)
### 5.5 Renormalization of the composite operators
In this section we give a solid foundation to the results of the previous renormalization procedure. More precisely we consider the renormalization of the composite operators that arise via the coupling of the chemical potential. Our attitude is similar to the standard procedure of mass renormalization in QCD: initially the mass of the light quarks is set to zero, then the bilinear mass operator $`\overline{\psi }\psi `$ is added, and its effect is accounted by its renormalization as a composite operator. We followed the same ideology so far, neglecting the fact that the renormalization of the coupled theory is different from that of the uncoupled. It is now that we remedy this.
Let us reconsider the computation of the free energy (23). In (24) we started with the bare Lagrangian. Instead we have to use the renormalized one, with counterterms coming from the wavefunction renormalization of the fields, coupling constant renormalization, and the renormalization of the composite operators, all these coming with a natural grading:
$$S[h,\phi ]+\mathrm{\Delta }S[h,\phi ]=\underset{i=2}{\overset{\mathrm{}}{}}\left(S_i[h,\phi ]+\mathrm{\Delta }S_i[h,\phi ]\right)A^i,$$
(79)
with the renormalized quantities in the right hand side. $`\mathrm{\Delta }S_i[h,\phi ]`$ commonly denotes all the counterterms of order $`i`$. The original Lagrangian that has to be employed reads:
$$^0=\frac{1}{2}(_\mu \psi )^2+\frac{1}{2}(_\mu \varphi )^2+\frac{1}{2}(_\mu \theta )^2+Aa\omega _{\mu \nu }\theta _\mu \varphi _\nu \psi +A^2\frac{1}{2}r\theta ^2(_\mu \varphi )^2.$$
(80)
$`\mathrm{\Delta }S_2`$ amounts simply to the multiplicative renormalization we considered in (55). Since multiplied by a factor of $`1/\lambda `$, it gives rise to a term that is independent of $`\lambda `$ (already used to renormalize the one loop free energy), one that is $`𝚘(\lambda )`$ (also used), and higher order terms:
$$\mathrm{\Delta }S_2=\mathrm{\Delta }S_2^0+\mathrm{\Delta }S_2^1\lambda +𝚘(\lambda ^2).$$
(81)
$`\mathrm{\Delta }S_1=0`$ as already $`S_1=0`$. A priori $`\mathrm{\Delta }S_0`$ contains terms from the from the wavefunction renormalization of the fields, but these are independent of $`h`$, and are canceled by the denominator in (23). Thus $`\mathrm{\Delta }S_0=\delta _{𝒪^0}𝒪^0+\delta _{𝒪^1}𝒪^1`$, where $`𝒪^0=r\frac{1}{2}h^2\theta ^2`$ and $`𝒪^1=ah\mathrm{\Omega }_\nu \theta _\nu \psi `$. The renormalization of $`𝒪^0`$ (and the similar operators that appear in the Lagrangian proportionally to $`h`$) has two contributions: one from the renormalization of $`h`$ (that is zero as discussed in the previous section) and one from the renormalization of the composite operator $`\theta ^2`$. We expect that $`\delta _{𝒪^0}=1+𝚘(\lambda ^2)`$, and this expectation can be confirmed by a short explicit computation.
In order to compute $`\delta _{𝒪^0}`$ we need a Green’s function involving $`𝒪^0`$. For simplicity let’s consider the one-point function of $`\theta ^2`$: $`\theta ^2(x)`$. As a renormalization prescription we normalize the one-point function according to the tree level value, and determine $`\delta _{\theta ^2}`$ from the condition of preserving the above normalization. At $`𝚘(\lambda )`$ there are two diagrams that contribute coming from the $`𝚘(A)`$ and $`𝚘(A^2)`$ vertices. The latter one has a value proportional to
$$d^nyG(xy)^2^2G(0)=d^nyG(xy)^2\delta (0)=0$$
(82)
in dimensional regularization. The same result is obtained if the diagram is evaluated in momentum space, where the masslessness of the fields requires additional IR regularization (as exploited already in a previous section).
The other diagram is readily proportional to $`\lambda a^2`$. The remaining integral can be computed but we don’t need the result for what follows, because once again, as we express this contribution in terms of the RG invariants we shall have a dependence proportional to $`\lambda ^2`$. Thus we have $`\delta _𝒪=1+𝚘(A^3)`$ at least. As a consistency check we quote that the same result is obtained by considering for example the $`\theta ^2(x)\theta (y)\theta (y^{})`$ Green’s function.
During the computation we employed the $`\frac{1}{2}h^2\theta ^2`$ operator as a mass term for $`\theta `$, though it is a composite operator as discussed above. The motivation for this can be given as follows: assume that $`\frac{1}{2}h^2\theta ^2`$ is a perturbation, and expand it perturbatively with the rest of the terms in (31). The difference is that $`\frac{1}{2}h^2\theta ^2`$ is independent of $`\lambda `$, and as such any term in its power series expansion is of the same order, and has to be summed. In other words, any term in the perturbative expansion of $``$ will be multiplied by the full expansion of the exponential of $`\frac{1}{2}h^2\theta ^2`$, that can be resummed. The resummation on the other hand is equivalent with the corresponding mass term $`h`$ for the $`\theta `$ field. From the point of view of the original $`h=0`$ theory this is a non-perturbative result.
Next we consider the renormalization of the second operator $`𝒪^1`$ and the one appearing at the next level $`𝚘(A)`$: $`𝒪^2=\theta ^2_0\varphi `$. These operators will have their composite operator renormalization functions: $`\delta _{𝒪^1}`$ and $`\delta _{𝒪^2}`$. Since the undeformed ($`h=0`$) theory has only interactions proportional to $`\lambda `$ or to $`a\sqrt{\lambda }`$, and the latter ones must appear at least twice for a finite contribution, eq. (80), we conclude that the renormalization effects due to these vertices are at least of order $`\lambda `$ and resp. $`\lambda a^2`$: $`\delta _{𝒪^1}=1+z^1a^2\lambda +\mathrm{}`$ resp. $`\delta _{𝒪^2}=1+z^2\lambda +\mathrm{}`$. In the spirit outlined above we have to introduce the new terms $`z^1𝒪^1`$ and $`z^2𝒪^2`$ into the action (79) as terms contributing to $`\mathrm{\Delta }S[h,\phi ]`$, and account for their contribution. As these new terms are proportional to $`\lambda `$, they contribute to $`\mathrm{\Delta }S_2[h,\phi ]`$. Thus we expect new contribution to $`_2`$. Based on (31) these are proportional to $`z^1𝒪^1`$ resp. $`z^2𝒪^2`$. But we already know from (44) that these are zero. It is easy to see that all the other combinations of these two operators with the rest of the operators will give higher order contributions to the free energy.
It is obvious that the product of $`z^1𝒪^1`$ (resp. $`z^2𝒪^2`$) with $`𝒪^1`$ (resp. $`𝒪^2`$) will give non-zero contribution to $`_3`$, but as far as $`_2`$ is concerned there is no deviation from the results we obtained.
At the next level (i.e. at $`𝚘(A^2)`$) we have the composite operators $`𝒪^3=\theta ^4`$ and $`𝒪^4=\theta ^3_\nu \psi `$. Consider first $`𝒪^3=\theta ^4`$. In order to compute $`\delta _{𝒪^3}`$ we fix the normalization of the correlation function $`\theta ^4(x)`$ to the tree level value: $`\theta ^4(x)=3\overline{G}(0)^2`$. We compute $`\delta _{𝒪^3}`$ by computing the first corrections to the correlation function, and insisting as above that the above normalization remains valid. Using the available interaction terms in (80) we see that using the vertices $`\theta _\mu \varphi _\nu \psi `$ resp. $`\theta ^2(_\mu \varphi )^2`$ give the first nontrivial contribution proportional to $`\lambda ^4`$ resp. $`\lambda ^2`$, while their combination is of order $`\lambda ^3`$. Thus we conclude that $`\delta _{𝒪^3}=1+𝚘(\lambda ^2)`$. In the case of $`𝒪^4`$ after simple considerations it is similarly obtained that $`\delta _{𝒪^4}=1+𝚘(\lambda ^2)`$, and we are done.
From the above analysis we conclude that the more or less naive renormalization procedure employed in the first place is completely adequate, and the additional composite operator renormalizations do not disturb the results as far as the $`_2`$ is concerned.
## 6 Conclusions
In this paper we have computed the change in the ground state energy density of a deformed principal $`SU(2)`$ sigma model, and one of its T-duals, in the presence of an external field. The computations has been carried out at two loop order in perturbation theory. Perfect agreement has been found in the following sense: there were renormalization schemes in the two models that yielded the same expressions for the two loop free energy densities.
We defined the free energy using the conserved charge corresponding to a special symmetry of the model. Care had to be made at the choice of the charge, since we wanted it to exist as a genuine symmetry in both the original and dual model, while it is known that the charge corresponding to the isometry used for the duality becomes topological. We performed the computation for the simplest choice, but other possibilities are also investigated at least in the one loop order , . We computed the relevant diagrams and performed the renormalization using dimensional regularization. Due to asymptotic freedom, we could improve the perturbative result using the renormalization group equations.
This result strengthens the confidence in the findings of , that were mainly based on beta function computations. More precisely this way we have given much stronger evidence that in the case of constant $`g_{00}`$, at least in the case of the deformed principal $`SU(2)`$ sigma model, the naive Buscher formula that relates the original model and its T-dual, gives a true quantum equivalence. As it was pointed out and exemplified in , the constancy of $`g_{00}`$ is no guarantee for two loop quantum equivalence. Based on this one might expect that a finer test would detect discrepancy even in the case where the beta function arguments show no sign of it. Though the free energy test that we performed does not prove the lack of discrepancies, it hints to the non-existence of these at least at two loop order. It would be interesting to analyze the pertinent $`SL(3)`$ example of from the free energy point of view, though, at this point it seems to be too big of a computational challenge.
Acknowledgments
We would like to thank J. Balog, P. Forgács for valuable discussions throughout this work. R.L.K. would also like to thank F.P. Esposito, D. Kastler, D. Maison, F. Mansouri, T. Nagy, M.E. Peskin, G. Pócsik, A Rebhan, A. Slavnov, P. Surányi, L.C.R. Wijewardhana, and L. Witten for discussions at different stages of this work. This work was supported in part by the Hungarian National Science Fund (OTKA) under T029802, and by the Ministry of Education under FKFP 0178/1999. R.L.K. was also supported in part by DOE grant DOE-FGO2-84ER40153. |
warning/0001/quant-ph0001066.html | ar5iv | text | # Efficient factorization with a single pure qubit and 𝑙𝑜𝑔𝑁 mixed qubits
## Abstract
It is commonly assumed that Shor’s quantum algorithm for the efficient factorization of a large number $`N`$ requires a pure initial state. Here we demonstrate that a single pure qubit together with a collection of $`log_2N`$ qubits in an arbitrary mixed state is sufficient to implement Shor’s factorization algorithm efficiently.
The discovery of a quantum algorithm for the efficient factorization of large numbers has started a rapid development of quantum information processing . Following this ground-breaking result a number of experimentally realizable proposals for the implementation of quantum computers have been made, for example, in ion trap systems or Nuclear Magnetic Resonance (NMR) schemes . These systems are distinguished by a low decoherence rate combined with a comparatively high gate speed and therefore promise the possibility of executing many quantum gates. While noise in these systems can be made small in principle, it nevertheless imposes limitations to the maximal size of the computation and to the achievable quality (e.g. the purity) of the initial state of the quantum computer. It would therefore be interesting to see whether a quantum computation necessarily requires the preparation of an initial state of high purity, or whether some parts of the quantum computer may be left in a mixed state. Such a result would be of particular interest in NMR systems in which it is difficult to prepare physically pure quantum states of nuclear spins.
The use of mixed states in quantum algorithms has had little discussion as yet. Note, however, the work of Schulman and Vazirani in which they demonstrated that, starting from a set of qubits each in a thermal state, one can obtain a certain number of pure qubits using a quantum algorithm. These were then envisaged to be used for a quantum computation, while all the other qubits which are in a mixed state are discarded. If the initial states are in a thermal mixture at high temperature, the number of mixed quantum states and quantum gates required to obtain even a single pure qubit is very high. It would greatly enhance the efficiency of this approach if it would be possible to reduce the necessary number of pure qubits as much as possible at the expense of employing some of the mixed qubits in the actual quantum computation. Recently, Knill and Laflamme have investigated the power of quantum computation when only a single pure qubit together with a supply of maximally mixed states is available. They were able to construct a problem that such a system can solve more efficiently than the best currently known classical algorithm.
It would be interesting to see whether these ideas can be extended to other problems of practical relevance. In this paper we demonstrate that a single pure qubit together with an initial supply of $`\mathrm{log}_2N`$ qubits in an arbitrarily mixed state is sufficient to implement Shor’s algorithm for the factorization of the number $`N`$ efficiently. This is the smallest number of pure states that can achieve this task. We also demonstrate that the efficiency of the modified algorithm is essentially independent of the degree of mixing of the $`\mathrm{log}_2N`$ qubits.
We proceed by outlining the problem addressed in Shor’s algorithm, followed by the formulation of Shor’s algorithm introduced in . Then we will describe the necessary modifications to this algorithm, that will allow it to be executed using a single pure qubit and $`\mathrm{log}_2N`$ qubits in a maximally mixed state.
The basis of Shor’s algorithm is a classical order finding method which, recast as a quantum algorithm, can be executed in polynomial time, requiring only a polynomial amount of additional classical computation to compute the factors of $`N`$. The factors of a number $`N=pq`$ can, with high probability, be found if the period or order, $`r`$, (the lowest positive integer $`x0`$ such that $`f_a(x)=1`$ ) of the element $`a`$ in the space of the function $`f_a(x)=a^x\text{mod}N,`$ is known. Then, provided $`a`$ is coprime to $`N`$ (which can be checked classically in polynomial time using Euclid’s algorithm), there is a high probability that $`\text{gcd}(a^{\frac{r}{2}}\pm 1,N)`$ yields a factor of $`N`$, where $`\text{gcd}(\alpha ,\beta )`$ denotes the greatest common divisor of $`\alpha `$ and $`\beta `$ which, again, can be determined efficiently using Euclid’s algorithm .
We begin by examining the formulation of Shor’s algorithm as given in and use it as a basis to demonstrate the main result of this paper. First of all we introduce the transformation $`U_a|x=|ax\text{mod}N`$ where $`x=0,\mathrm{},N1.`$ Provided $`a`$ is coprime to $`N`$ this is a unitary transformation and has eigenvectors
$$|\psi _j=\underset{k=0}{\overset{r1}{}}e^{\frac{2\pi ijk}{r}}|a^k\text{mod}Nj=0,\mathrm{},r1$$
(1)
with corresponding eigenvalues $`e^{\frac{2\pi ij}{r}}`$. Given one of these eigenvectors we can apply $`U_a`$ to it and the value of $`r`$ will be encoded in the phase, $`e^{\frac{2\pi ij}{r}}`$. This, however, is a global phase which we cannot measure so instead we can use the ”phase-kickback” technique requiring the conditional unitary transformation given by
$`cU_a|0|x`$ $`=`$ $`|0|x;cU_a|1|x=|1|ax\text{mod}N.`$ (2)
The effect of applying the controlled unitary transform to the state $`(|0+|1)|\psi _j`$ is
$$cU_a(|0+|1)|\psi _j=(|0+e^{\frac{2\pi ij}{r}}|1)|\psi _j$$
(3)
’kicking’ the ’global’ phase shift acquired on the second qubit into a relative phase in the first qubit. We can now perform measurements on the first qubit which will allow us to estimate $`r`$, however, we cannot create the eigenstates of $`U_a`$ without knowledge of $`r`$. Instead one can use the fact that $`_{j=0}^{r1}|\psi _j=|1`$ and conditionally apply $`U_a`$ to the state $`|1`$ (which obviously requires no knowledge of $`r`$) in the second qubit
$$cU_a(|0+|1)|1=\underset{j=0}{\overset{r1}{}}(|0+e^{\frac{2\pi ij}{r}}|1)|\psi _j.$$
(4)
This state is, of course, entangled, so when we make measurements on the first qubit we will get an estimate of $`e^{\frac{2\pi ij}{r}}`$, with $`j`$ (which corresponds to an eigenstate) selected at random.
How do we estimate this phase and the value of $`r`$ accurately? The network in Fig. 1 will give us, with a sufficient probability, the best $`L`$-bit estimate of the value of $`2^Lj/r`$ .
As the algorithm proceeds it uses the controlled $`U_a,U_a^2,U_a^{2^2},\mathrm{},U_a^{2^L}`$ transformations to produce the ’kicked’ phases $`e^{\frac{2\pi ij}{r}},e^{\frac{2^2\pi ij}{r}},e^{\frac{2^3\pi ij}{r}},\mathrm{},e^{\frac{2^{L1}\pi ij}{r}}`$ into the upper ’control’ qubits. The remaining operations on the control qubits realise the quantum inverse Fourier transform. A measurement on each of these qubits produces a binary number $`c=_{i=0}^{L1}2^im_i`$ such that with a finite probability $`c/2^L`$ is the best estimate of $`j/r`$ for some integer $`j`$ again selected at random on measurement.
The first modification to this algorithm comes when we notice that the gates within the Fourier transform are applied sequentially on the qubits. Thus instead of performing the entire transform and then making measurements on all control qubits afterwards we may apply the single qubit (Hadamard) operation to the first qubit and then measure it. The operations (controlled phase shifts) controlled by this first qubit are then replaced by single qubit operations given the result of the measurement on the first. This ’semi-classical’ modification preserves the probabilities of all measurement results.
Taking this further we need only insist on one control qubit and the remaining $`\mathrm{log}_2N`$ qubits as we can ’recycle’ the control qubit after each measurement (Fig. 2): we perform all the necessary operations of the first control qubit including measurements, followed by all the operations of the second control qubit on the same physical qubit system given the results of previous measurements, and so on .
We can, therefore, already implement Shor’s algorithm with $`1+\mathrm{log}_2N`$ pure qubits that is, one control qubit and $`\mathrm{log}_2N`$ of the remaining qubits. We will find later that we can also replace the $`\mathrm{log}_2N`$ pure qubits with $`\mathrm{log}_2N`$ maximally mixed qubits and find the order $`r`$ efficiently (see also ). To see why this is the case we first need to examine the unitary transformation $`U_a`$ more closely.
The unitarity of the transform together with the fact that it maps a ’number’ state $`|x`$ to a ’number’ state $`|ax\text{mod}N`$ means that on repeated application of $`U_a`$ periodic sequences are induced on all the numbers $`x=0,1,\mathrm{},N1`$, that is, there is an $`R(x)`$ such that $`U_a^{R(x)}|x=|x`$. We may write the members of all possible sequences as $`|ga^x\text{mod}N\text{for some }g\text{ and }x\text{.}`$ For example, for $`a=2`$ and $`N=15`$ on repeated application of $`U_a`$ the possible sequences are
$`g=1:`$ $`|1|2|4|8|1`$ (5)
$`g=3:`$ $`|3|6|12|9|3`$ (6)
$`g=5:`$ $`|5|10|5`$ (7)
$`g=7:`$ $`|7|14|13|11|7.`$ (8)
It is the first of these sequences (with $`g=1`$) whose number of members is what we previously called the ’order’, $`r`$, of $`a`$ modulo $`N`$ and it is this period that we need to find to factorize $`N`$. However, there is a relationship between the order of the sequence with $`g=1`$ and the orders of all the other sequences with $`g1`$. We will label each of the different sequences by $`d`$ and the number of members in each sequence by $`r_d`$. $`U_a`$ obeys the condition $`U_a^r=I`$ so it is clear that $`r_d|r`$, that is, the orders of all the sequences divides that of the sequence with $`g=1`$. In fact we will find that nearly all of the numbers $`0,1,\mathrm{},N1`$ are contained within a sequence that has the same order as the first sequence. We can find a lower bound on the probability that for a number $`g0,1,\mathrm{},N1`$ the state $`|ga^x\text{mod}N`$ is contained within a sequence of order $`r`$.
Theorem 1 Given two prime numbers $`p`$ and $`q`$ we define r as the lowest positive integer x such that $`a^x10\text{mod}(pq)`$ for an arbitrary integer $`a`$. Then $`ga^xg0\text{mod}(pq)`$ with $`x<r`$ for at most $`p+q1`$ values of $`g`$ in the interval $`0apq1`$.
Proof: If $`\text{gcd}(g,pq)=1`$ then $`g(a^x1)0\text{mod}pqa^x10\text{mod}pq`$ and therefore $`x=r`$. There are $`(p1)(q1)`$ positive integers less than and coprime to $`pq`$, which proves the theorem
We can now see that the probability, $`P_r`$, of picking $`g`$ at random such that the lowest $`x`$ for which $`ga^xg\text{mod}(pq)`$ is $`r`$, is $`P_r(pq(p+q1))/pq=(p1)(q1)/pq`$ which approaches unity as $`p`$ and $`q`$ become large.
This tells us that if we set up an algorithm that actually finds the order of a random sequence we still have a good chance that this order is in fact $`r`$.
The $`r`$ eigenstates of $`U_a`$ in equation 1 are orthogonal superpositions of the members of the sequence with $`g=1`$. In exactly the same way we can form the remaining $`Nr`$ eigenstates of $`U_a`$ as orthogonal superpositions of members of each of the other sequences. We write these as
$$|\psi _{j_d}^d=\underset{k=0}{\overset{r_d1}{}}e^{\frac{2\pi ij_dk}{r_d}}|g_da^k\text{mod}N$$
(9)
where $`d`$ labels the sequence and $`j_d=0,\mathrm{},r_d1`$ the eigenstates of $`U_a`$ within the sequence $`d`$. $`|g_d`$ is the lowest member of the $`d`$th sequence. Each eigenstate has corresponding eigenvalue $`e^{2\pi ij_d/r_d}`$ so using the same phase estimation techniques allows us to estimate $`j_d/r_d`$ given the state $`|\psi _{j_d}^d`$. Again, this requires knowledge of the sequences induced by $`U_a`$ so instead we may perform the phase estimation technique on the maximally mixed state
$$\frac{\mathrm{𝟏}}{N}=\frac{1}{N}\underset{k=0}{\overset{N1}{}}|kk|=\frac{1}{N}\underset{d}{}\underset{j_d=0}{\overset{r_d1}{}}|\psi _{j_d}^d\psi _{j_d}^d|.$$
(10)
Phase estimation now estimates the value of $`j_d/r_d`$ for $`j_d`$ and $`d`$ chosen at random but as we have seen above nearly all the orders $`r_d`$ are equal to $`r`$.
Note that in Shor’s original algorithm the $`\mathrm{log}_2N`$ qubits encode a phase change into the control qubits which is quantum mechanically correlated to eigenstates of $`U_a`$ our modification encodes a phase change which is classically correlated to the eigenstates. This includes not only the group of eigenstates consisting of superpositions of elements in the first sequence (see Eq. 5) but groups of eigenstates consisting of superpositions of elements in each of the other sequences. However by theorem 1 most of these sequences have the same order and will encode the value $`r_d=r`$ into the control qubits. This makes it intuitively clear that the algorithm is still efficient. Note however that although the $`\mathrm{log}_2N`$ mixed qubits are only classically correlated to the pure qubit, entanglement still exists in the system: one can partition the system into two halves one containing some mixed qubits and the other containing the remaining mixed qubits and the pure qubit. Then it can be checked, that this bipartite system can have negative partial transpose and is therefore entangled .
In the following we will prove strictly that this modified version of Shor’s algorithm is indeed still efficient for order finding. Shor’s algorithm requires $`O\left(\mathrm{log}\mathrm{log}r\right)`$ repetitions for it to have a high chance of finding the order whereas the mixed state Shor’s algorithm uses exactly the same resources as Shor’s original algorithm but requires
$$O\left(\frac{pq}{(p1)(q1)}\mathrm{log}\mathrm{log}r\right)$$
(11)
repetitions for it to have a high chance of finding the order which, in the limit $`p,q\mathrm{}`$, is equally as efficient as Shor’s algorithm. For simplicity we will prove this efficiency result for a mixed state algorithm with $`L`$ control qubits. For the reasons outlined above the result will be identical using a single pure control qubit. The proof follows very closely that of Shor .
Pick an $`L`$ such the $`N^2<t=2^L<2N^2`$. The initial state of our system with all the control qubits grouped into the first state is
$$\rho _{ini}=\frac{1}{Nt}\underset{a=0}{\overset{t1}{}}\underset{b=0}{\overset{t1}{}}|ab|\underset{d}{}\underset{j_d=0}{\overset{r_d1}{}}|\psi _j^d\psi _j^d|.$$
(12)
Application of the controlled $`U_a,U_a^2,\mathrm{},U_a^{2^{L1}}`$ gates and the inverse Fourier transform yields the state
$`\rho _2={\displaystyle \frac{1}{Nt^2}}{\displaystyle \underset{d}{}}{\displaystyle \underset{j_d=0}{\overset{r_d1}{}}}{\displaystyle \underset{a,b,k,l=0}{\overset{t1}{}}}e^{2\pi ia\left(\frac{j_d}{r_d}\frac{k}{t}\right)}`$ (13)
$`e^{2\pi ib\left(\frac{j_d}{r_d}\frac{l}{t}\right)}|kl||\psi _{j_d}^d\psi _{j_d}^d|.`$ (14)
We now make a measurement on the first state. The probability that the result $`c`$ is obtained is
$`P(c)={\displaystyle \frac{1}{Nt^2}}{\displaystyle \underset{d}{}}{\displaystyle \underset{j_d=0}{\overset{r_d1}{}}}\left|S\right|^2,S={\displaystyle \underset{a=0}{\overset{t1}{}}}e^{2\pi ia\left(\frac{j_d}{r_d}\frac{c}{t}\right)}.`$ (15)
$`S`$ is just an arithmetic progression and $`|S|^2`$ can easily be bounded by
$$|S|^2>\frac{4t^2}{\pi ^2}\text{for}\left|\frac{j_d}{r_d}\frac{c}{t}\right|<\frac{1}{2t}.$$
(16)
Because $`t>N^2`$ this is a sufficient condition that given $`c/t`$ there is only one fraction $`j_d/r_d`$ with $`r_d<N`$ such that the above condition is obeyed. For a given measurement result $`c`$ there are at least $`(p1)(q1)/r`$ corresponding values of $`r_d`$ with $`r_d=r`$ by theorem 1. So the probability that $`c/t`$ is the best estimate of a fraction with denominator $`r`$ is
$$P^{}(c)>\frac{1}{Nt^2}\underset{d}{}\underset{j_d=1}{\overset{r_d1}{}}|S|^2>\frac{4(p1)(q1)}{N\pi ^2r}.$$
(17)
We now require that the numerator, $`j_d`$, is coprime to $`r`$ otherwise cancellation of common factors will occur in $`j_d/r`$. There are $`\varphi (r)`$ values of $`j_d`$ which are less than and coprime to $`r`$, where $`\varphi `$ is Euler’s totient function . Thus the probability that we can calculate $`r`$ is $`P>4(p1)(q1)\varphi (r)/Nr\pi ^2.`$ Using a theorem by Hardy and Wright (theorem 328) that $`\varphi (r)/r>\delta /\mathrm{log}\mathrm{log}r`$ for some constant $`\delta `$ we find that the number of times that we need run the algorithm to have a high chance of finding the period, $`r`$ is given by Eq. (11).
We have thus found that one pure qubit and a supply of maximally mixed qubits is sufficient to implement Shor’s algorithm, requiring no more resources in terms of quantum operations or physical systems than the algorithm operating on pure quantum states. This implies that the algorithm presented here is a ’true’ quantum algorithm, achieving an exponential speedup using only polynomial resources. This may be suprising as the degree of mixing of the state of the computer is high. However, the mixing decrease as the algorithm proceeds but never below a mixture of $`N/r_d`$ eigenstates where $`r_d`$ is the measured period. Furthermore, it should be noted that despite this strong degree of mixing the quantum computer actually evolves into an entangled state. It is this entanglement that appears to be responsible for the computational speedup.
Maximally mixed states are intuitively a less ’costly’ resource than pure states but, in fact, we do not need to require maximally mixed states: we could equally well use any random state (mixed or pure) on which to perform the controlled $`U_a`$ operations. The average efficiency over all these states would then be as we have shown in this paper. In particular thermal states of nuclear spins (e.g. in NMR), where the occupation of the ground state is only slightly greater than that of the first excited state, would change the efficiency of this algorithm by only a small amount leaving it an efficient algorithm. This ability of highly mixed states to support efficient quantum computation points towards the possibility of the implementation of true quantum computation for example in NMR systems.
The authors would like to thank Kevin Buzzard for valuable advice on number theory and S. Bose, R. Jozsa, R. Laflamme, M. Nielsen and C. Zalka for helpful comments. This work is supported by the EPSRC, the Leverhulme Trust, and the EU project EQUIP. |
warning/0001/astro-ph0001359.html | ar5iv | text | # Gravitational microlensing source limb-darkening and limb-polarization, I: angle-averaged amplification functions
## 1 Introduction
The most common interest in microlensing events is as a probe of the lensing object itself. Recently, however, there has been increasing interest in such events as probes of the otherwise unresolved sources \[Valls-Gabaud 1998, Mao and Witt 1998, Gaudi and Gould 1999\]. When the lens transits the source (or nearly so), it breaks any rotational symmetry, and this gives us access to the surface brightness, as well as polarization \[Simmons et al. 1995, Simmons et al. 1995\] and chromaticity \[Valls-Gabaud 1995, Valls-Gabaud 1998\] information. Previous work on extended source effects has concentrated on the forward problem and generally either performed the required calculations numerically rather than analytically, or used an approximate form of the amplification function. Also, much of the work on source effects has relied on the high amplification provided by binary lens causic crossings, rather than the amplification of a single lens.
The gravitational lensing forward problem – that of predicting centroid motion and magnification for a given set of source parameters – is relatively easy. The problem is also, however, typically poorly conditioned, in the sense that there will be a broad range of limb-darkening or limb-polarization functions which could plausibly correspond to the observed signal in a microlensing event. This means that a parameter-fitting approach to recovering these functions is very dangerous.
We can make progress by expressing the problem explicitly as an inverse problem (IP), and using the technology of IP methods to analyse precisely what information can be recovered for a given set of observations.
That is the subject of, and motivation for, a forthcoming paper \[Gray and Coleman 2000\]; here I am concerned with identifying the angle-averaged amplification functions as IP kernels, and obtaining analytic expressions for them. As well as facilitating the IP analysis of the problem, these analytic kernels can help in the treatment of the forward problem, since they can be evaluated more efficiently than by a numerical integration, and with high accuracy over their entire domain.
In Sect. 2, I define the amplification functions as IP kernels, in Sect. 3, I integrate them about the source’s centre, and in Sect. 4, I present the results of that integration.
## 2 Amplification functions as inverse problem kernels
The geometry of a microlensing event is as shown in Fig. 1.
The gravitational lens amplification function is \[Schneider et al. 1992\]
$$A(\xi )=\frac{1}{2}\left(\xi +\frac{1}{\xi }\right),$$
(1)
where
$$\xi =\left(1+\frac{4}{\zeta ^2}\right)^{1/2},\zeta ^2=r^2+s^22rs\mathrm{cos}(\chi \varphi ).$$
Denote the intensity by $`I(r)`$ and the Stokes parameter by $`Q(r,\chi )=P(r)\mathrm{cos}2\chi `$, where $`P(r)`$ is the polarization of the stellar surface, and we are assuming that the surface brightness is rotationally symmetric. In the case of a microlensing event, we cannot resolve details of the lensed source, and must therefore measure integrals over the source surface. We immediately obtain
$`I(s(t),\varphi (t))`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}I(r)\stackrel{~}{A}_I(r;s,\varphi )dr`$ (2)
$`Q(s(t),\varphi (t))`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}P(r)\stackrel{~}{A}_Q(r;s,\varphi )dr`$ (3)
where the amplification kernels are
$`\stackrel{~}{A}_I`$ $`=`$ $`r{\displaystyle _0^{2\pi }}A(r,\chi ;s,\varphi )d\chi `$ (4)
$`\stackrel{~}{A}_Q`$ $`=`$ $`r{\displaystyle _0^{2\pi }}\mathrm{cos}2\chi A(r,\chi ;s,\varphi )d\chi .`$ (5)
Note that the kernel $`\stackrel{~}{A}_I`$ is a factor $`2\pi r`$ times the angular average of the amplification function, and the functions $`I(s,\varphi )`$ and $`Q(s,\varphi )`$ have the dimensions of flux rather than intensity.
Equations (2) and (3) are in the form of an inverse problem. We address the inverse problem in a forthcoming paper \[Gray and Coleman 2000\]. The evaluation of the integrals $`\stackrel{~}{A}_{I,Q}`$ is rather hard, and I describe it in this paper.
## 3 Integration of the amplification functions
Write
$$z=\mathrm{exp}i(\chi \varphi ),\mathrm{d}z=iz\mathrm{d}\chi ,$$
(6)
so that
$$\zeta ^2=r^2+s^2rs\left(z+\frac{1}{z}\right).$$
Define
$$a_1=\frac{r^2+s^2}{2rs},a_2=\frac{4+r^2+s^2}{2rs}$$
(7)
and
$`x_1`$ $`=`$ $`a_1+\sqrt{a_1^21}=\{\begin{array}{cc}r/s& ,rs\hfill \\ s/r& ,r<s\hfill \end{array}`$ (10)
$`x_2`$ $`=`$ $`a_2+\sqrt{a_2^21}.`$ (11)
Defining $`\overline{x}_i=1/x_i`$, we have
$$x_i+\overline{x}_i=2a_i,x_i\overline{x}_i=1.$$
It is easy to show that
$$0<\overline{x}_2<\overline{x}_11x_1<x_2.$$
Rewriting the expression for $`\xi `$, we obtain
$$\xi \left(1+\frac{4}{\zeta ^2}\right)^{1/2}=\left(\frac{(zx_2)(z\overline{x}_2)}{(zx_1)(z\overline{x}_1)}\right)^{1/2}.$$
(12)
We now evaluate the integral
$$\stackrel{~}{I}=_C\frac{1}{z}\left(\xi +\frac{1}{\xi }\right)dz,$$
(13)
where $`C`$ is the contour shown in Fig. 2. The integrand has poles at $`z=x_{1,2}`$, $`z=\overline{x}_{1,2}`$ and $`z=0`$.
The contour encloses a single singularity at $`z=0`$. We have $`\xi (z=0)=1`$, so that
$$\stackrel{~}{I}=2\pi i\times Res(\stackrel{~}{I}(z=0))=4\pi i.$$
(14)
For contour 1, substitute $`z=\mathrm{e}^{i\psi }`$, for $`\psi [0,2\pi ]`$. Then $`\stackrel{~}{I}^{(1)}=i_0^{2\pi }(\xi +1/\xi )d\psi `$, and the substitution $`\psi =\chi \varphi `$ produces
$$\stackrel{~}{I}^{(1)}=2i_\varphi ^{\varphi +2\pi }A\left(\xi (\chi ,\varphi )\right)d\chi =\frac{2i}{r}\stackrel{~}{A}_I(r;s,\varphi ).$$
(15)
Contours 5 and 6 cancel, and with the substitution $`z=\overline{x}_2+\rho \mathrm{e}^{i\theta }`$, it is clear that $`\stackrel{~}{I}^{(3)}0`$ as $`\rho 0`$.
Now turn to contours 2 and 4. By substituting $`z=\overline{x}_2+\sigma `$, $`\sigma =0(\overline{x}_1\overline{x}_2)`$, into $`\stackrel{~}{I}^{(4)}`$, substituting $`z=\overline{x}_2+\sigma \mathrm{e}^{i2\pi }`$, $`\sigma =(\overline{x}_1\overline{x}_2)`$, into $`\stackrel{~}{I}^{(2)}`$, and noting that $`0<\sigma <\overline{x}_1\overline{x}_2`$ (so that $`|\sigma |<|\overline{x}_2x_2|`$, $`|\overline{x}_2x_1|`$, and $`|\overline{x}_2\overline{x}_1|`$, so that the phases of the corresponding square-rooted factors in $`\stackrel{~}{I}^{(2)}`$ are unaffected by the factor of $`\mathrm{e}^{i2\pi }`$), we can see that
$$\stackrel{~}{I}^{(2)}=\stackrel{~}{I}^{(4)}.$$
(16)
Now substituting $`z=x`$ ($`x`$ real) directly into Eqn. (13), we obtain
$$\stackrel{~}{I}^{(4)}=i(I_1I_2),$$
(17)
where
$`I_1`$ $`=`$ $`{\displaystyle _{\overline{x}_2}^{\overline{x}_1}}{\displaystyle \frac{1}{x}}\left({\displaystyle \frac{(x_2x)(x\overline{x}_2)}{(x_1x)(\overline{x}_1x)}}\right)^{1/2}dx`$ (18)
$`I_2`$ $`=`$ $`{\displaystyle _{\overline{x}_2}^{\overline{x}_1}}{\displaystyle \frac{1}{x}}\left({\displaystyle \frac{(x_2x)(x\overline{x}_2)}{(x_1x)(\overline{x}_1x)}}\right)^{1/2}dx.`$ (19)
Using the notation of \[Carlson 1988\],
$$[p_1,\mathrm{},p_n]\underset{i=1}{\overset{n}{}}(a_i+b_it)^{p_i/2}\mathrm{d}t,$$
we may rewrite these as
$`I_1`$ $`=`$ $`[+1,1,1,+1,2]`$ (20)
$`I_2`$ $`=`$ $`[1,+1,+1,1,2]`$ (21)
with $`a_i=(\overline{x}_2,\overline{x}_1,x_1,x_2,0)`$ and $`b_i=(+1,1,1,1,+1)`$ in both cases. We will evaluate these integrals when we return to them below, in Sect. 4.
Thus, collecting Eqns. (14), (15), (16) and (17), we obtain the angle-averaged amplification function as
$$\stackrel{~}{A}_I(r;s)=r(2\pi +I_1I_2).$$
(22)
This has been confirmed by direct numerical integration of the integrand.
Turning now to Eqn. (5), we may again substitute $`z=\mathrm{exp}i(\chi \varphi )`$, and obtain
$$\mathrm{cos}2\chi =\frac{1}{2}\left(z^2\mathrm{e}^{i2\varphi }+z^2\mathrm{e}^{i2\varphi }\right).$$
Now evaluate the integral
$$\stackrel{~}{Q}=_C\left(z\mathrm{e}^{i2\varphi }+z^3\mathrm{e}^{i2\varphi }\right)\left(\xi +\frac{1}{\xi }\right)dz,$$
(23)
with the same contour as above, and with $`\xi `$ as in Eqn. (12). This has a third-order pole at $`z=0`$, which means that the residue is
$$Res(\stackrel{~}{Q})=a_1=\frac{1}{2}\frac{\mathrm{d}^2}{\mathrm{d}z^2}\left[\xi +\frac{1}{\xi }\right]_{z=0}=(a_1a_2)^2$$
(24)
with $`a_i`$ as defined in Eqn. (7) above. Thus
$$\stackrel{~}{Q}=2\pi i\mathrm{e}^{i2\varphi }(a_1a_2)^2.$$
(25)
Much of the calculation goes through as before. Substituting $`z=\mathrm{e}^{i\psi }`$, we obtain,
$$\stackrel{~}{Q}^{(1)}=\frac{4i}{r}\stackrel{~}{A}_Q(r;s,\varphi ).$$
(26)
Contours 5 and 6 cancel, and contour 3 makes zero contribution in the $`\rho 0`$ limit. Similarly,
$$\stackrel{~}{Q}^{(2)}=\stackrel{~}{Q}^{(4)}=i\mathrm{e}^{i2\varphi }(Q_1+Q_2)+i\mathrm{e}^{i2\varphi }(Q_3+Q_4),$$
(27)
where
$`Q_1`$ $`=`$ $`[+1,1,1,+1,+2]`$ (28)
$`Q_2`$ $`=`$ $`[1,+1,+1,1,+2]`$ (29)
$`Q_3`$ $`=`$ $`[+1,1,1,+1,6]`$ (30)
$`Q_4`$ $`=`$ $`[1,+1,+1,1,6]`$ (31)
with $`a_i`$ and $`b_i`$ as above.
Assembling equations (25), (26) and (27), we obtain
$`\stackrel{~}{A}_Q(r;s,\varphi )`$ $`=`$ $`{\displaystyle \frac{r}{2}}[\mathrm{e}^{i2\varphi }(Q_1+Q_2)+\mathrm{e}^{i2\varphi }(Q_3+Q_4)`$
$`\pi \mathrm{e}^{i2\varphi }(a_1a_2)^2].`$
However, this integral should be real, so the imaginary part must be zero:
$$Q_1+Q_2+Q_3Q_4+\pi (a_1a_2)^20.$$
Thus, the final expression for the angle-averaged polarization amplification function is
$`\stackrel{~}{A}_Q(r;s,\varphi )`$ $`=`$ $`r\mathrm{cos}2\varphi (Q_1Q_2)`$
$`=`$ $`r\left(2{\displaystyle \frac{l^2}{s^2}}1\right)(Q_1Q_2).`$
This also has been confirmed by numerical integration.
Eqn. (22) and Eqn. (3) are the principal results of this paper.
In Fig. 3,
I show the amplification function $`\stackrel{~}{A}_I(r;s)`$, with a singularity along the line $`r=s`$. Note particularly the breadth of the kernel after the singularity: although the angle-averaged amplification is close to $`1`$ away from a well-defined peak at $`r=s`$, the factor of $`r`$ in Eqn. (4) means that there are contributions to $`I(t)`$ in Eqn. (2) from a broad range of $`I(r)`$, a situation especially severe for cases where the source function $`I(r)`$ extends significantly beyond $`r=1`$ (ie, those cases where $`R_{}>R_E`$). The breadth and asymmetry of the kernel is what makes the recovery of the source function so problematic. Similarly, Fig. 4
shows the amplification kernel $`\stackrel{~}{A}_Q(r;s,\varphi )`$. This has the same factor of $`r`$ as $`A_I(r;s)`$, but because the underlying angle-averaged function is close to zero away from $`r=s`$, this has a less damaging effect, so that the polarization signal should be easier to recover (which is fortunate, since that signal is so much harder to detect than the intensity signal). The kernel is still, however, both broad and asymmetric.
## 4 Amplification functions as elliptic integrals
The integrals defined by equations (20) and (21), and equations (28) to (31), are elliptic integrals. Carlson \[Carlson 1988\] provides a set of recurrence relations to reduce such integrals to a small set of elementary integrals, which are in turn expressed in terms of a set of functions $`R_J`$, $`R_F`$ and $`R_D`$ which are more symmetrical alternatives to the traditional Legendre elliptic integrals, and which can be evaluated efficiently, and to high accuracy, across their entire domain. A suitable algorithm \[Gray 2000\] can do this (algebraically exhausting) reduction mechanically, to obtain
$`{\displaystyle \frac{\stackrel{~}{A}_I(r;s)}{r}}2\pi `$ $`=`$ $`I_1I_2`$ (34)
$`=`$ $`2{\displaystyle \frac{x_1x_2}{\sqrt{x_1x_2}}}R_F`$
$`+{\displaystyle \frac{4}{3}}\lambda _2(1+\lambda _1)\sqrt{x_1x_2}`$
$`\times \left(x_2R_{J}^{}{}_{1}{}^{}{\displaystyle \frac{R_{J}^{}{}_{2}{}^{}}{x_2}}\right)`$
$`=`$ $`2{\displaystyle \frac{x_1x_2}{\sqrt{x_1x_2}}}K(\kappa )`$
$`+4{\displaystyle \frac{1+\lambda _1}{x_2}}\sqrt{x_1x_2}(\mathrm{\Pi }_2\mathrm{\Pi }_1),`$
$`{\displaystyle \frac{\stackrel{~}{A}_Q(r;s,\varphi )}{r\mathrm{cos}2\varphi }}`$ $`=`$ $`Q_1Q_2`$ (35)
$`=`$ $`{\displaystyle \frac{q_DR_D+q_FR_F+q_JR_{J}^{}{}_{2}{}^{}}{2(x_1x_2)^{3/2}}}`$
$`=`$ $`{\displaystyle \frac{q_eE(\kappa )+q_kK(\kappa )+q_\pi \mathrm{\Pi }_2}{2(x_1x_2)^{3/2}}}.`$ (36)
Here
$`\lambda _1`$ $`=`$ $`{\displaystyle \frac{x_2x_1}{x_1\overline{x}_2}}>0,`$
$`\lambda _2`$ $`=`$ $`{\displaystyle \frac{\overline{x}_1\overline{x}_2}{x_2\overline{x}_1}}>0,`$
$`\kappa ^2`$ $`=`$ $`\lambda _1\lambda _2(0<\kappa ^21);`$
the coefficients $`q_a`$ are
$`q_D`$ $`=`$ $`{\displaystyle \frac{(x_1+x_2)(x_1x_2+1)(x_2x_1)^2}{3(x_1x_21)}},`$
$`q_F`$ $`=`$ $`{\displaystyle \frac{(x_1x_21)(x_2x_1)^2}{x_2}},`$
$`q_J`$ $`=`$ $`{\displaystyle \frac{(x_2^21)(x_2x_1)^3}{3x_2^2}},`$
$`q_e`$ $`=`$ $`{\displaystyle \frac{3q_D}{\kappa ^2}}=(x_1+x_2)(x_1^2x_2^21),`$
$`q_k`$ $`=`$ $`{\displaystyle \frac{3q_d}{\kappa ^2}}+q_F+{\displaystyle \frac{3q_J}{\lambda _1}}`$
$`=`$ $`(x_2^3+3x_1x_2^2+x_1+x_2)(x_1x_21),`$
$`q_\pi `$ $`=`$ $`{\displaystyle \frac{3q_J}{\lambda _1}}={\displaystyle \frac{(x_2^21)(x_1x_21)(x_2x_1)^2}{x_2}};`$
we have used the abbreviations
$`R_F`$ $`=`$ $`R_F(0,1\kappa ^2,1),`$
$`R_D`$ $`=`$ $`R_D(0,1\kappa ^2,1),`$
$`R_{J}^{}{}_{i}{}^{}`$ $`=`$ $`R_J(0,1\kappa ^2,1,1+\lambda _i),`$
$`\mathrm{\Pi }_i`$ $`=`$ $`\mathrm{\Pi }({\displaystyle \frac{\pi }{2}},\lambda _i,\kappa );`$
and $`\mathrm{\Pi }(\varphi =\pi /2,n,k)`$, $`K(k)`$ and $`E(k)`$ are the Legendre forms of the complete elliptic integrals, as defined in \[Gradshteyn and Ryzhik 1994, 8.111\], or \[Abramowitz and Stegun 1965, 17.2\] (with $`k=\mathrm{sin}\alpha `$). These are related to the Carlson forms of the integrals by
$`R_J(0,1\kappa ^2,1;1+\lambda )`$ $`=`$ $`{\displaystyle \frac{3}{\lambda }}\left[K(\kappa )\mathrm{\Pi }(\pi /2,\lambda ,\kappa )\right]`$
$`R_D(0,1\kappa ^2,1)`$ $`=`$ $`{\displaystyle \frac{3}{\kappa ^2}}[K(\kappa )E(\kappa )]`$
$`R_F(0,1\kappa ^2,1)`$ $`=`$ $`K(\kappa ).`$
### 4.1 Singularies and asymptotic behaviour of $`\stackrel{~}{A}_I`$ and $`\stackrel{~}{A}_Q`$
We can write Eqn. (34) in terms of Heuman’s Lambda function \[Abramowitz and Stegun 1965, 17.4.39\] as follows:
$`I_1I_2`$ $`=`$ $`2{\displaystyle \frac{x_2x_1}{\sqrt{x_1x_2}}}K(\kappa )`$ (37)
$`+2\pi [\mathrm{\Lambda }_0(\mathrm{arcsin}(1+\lambda _1)^{1/2}\backslash \alpha )`$
$`(\lambda _1\lambda _2)].`$
The advantage of this is that we can now easily isolate the singularity at $`r=s`$, where we have $`x_1=\overline{x}_1=1`$, $`x_2>1`$, and thus $`\kappa =1`$. The $`\mathrm{\Lambda }_0`$ function has no singularities, and the coefficient of $`K(\kappa )`$ is finite there, so the only singularity is at $`K(\kappa =1)`$ where \[Abramowitz and Stegun 1965, 17.3.26\]
$$\underset{\kappa 1}{lim}\left(K\frac{1}{2}\mathrm{ln}\frac{16}{1\kappa ^2}\right)=0.$$
Thus the leading order term in Eqn. (34) at $`r=s`$ is
$`{\displaystyle \frac{\stackrel{~}{A}_I(rs;s)}{r}}2\pi `$ $`=`$ $`I_1I_2(rs)`$ (38)
$``$ $`{\displaystyle \frac{x_2x_1}{\sqrt{x_1x_2}}}\mathrm{ln}{\displaystyle \frac{16}{1\kappa ^2}}.`$
We can also confirm the behaviour as $`r0`$ and $`r\mathrm{}`$. For $`r>s`$, we have
$$\frac{x_2}{x_1}=1+\frac{1}{2}\frac{(4+s^2)(4s^2)}{r^2}+O(r^4)$$
(39)
and
$$\sqrt{\frac{x_2}{x_1}}\sqrt{\frac{x_1}{x_2}}=\frac{1}{2}\frac{(4+s^2)(4s^2)}{r^2}+O(r^4).$$
(40)
It follows from Eqn. (39) that both $`(1+\lambda _1)`$ and $`(1+\lambda _2)`$ go to $`1`$ as $`r\mathrm{}`$, so that the difference of Lambda functions in Eqn. (37) goes to zero. The factor $`\kappa 0`$ in this limit, and $`K(0)`$ is finite, but the coefficient of $`K`$ goes to zero, from Eqn. (40), so $`I_1I_20`$, and $`\stackrel{~}{A}_I2\pi r`$, as expected.
For $`r<s`$, both $`x_1`$ and $`x_2`$ diverge as $`r0`$, but $`x_2/x_1=1+4/s^2+O(r^2)`$, so that $`\kappa 0`$. Both $`K`$ and $`\mathrm{\Lambda }_0`$ are finite here, so that $`I_1I_2`$ does not diverge, and the singularity is confined to the coordinates $`x_i`$.
We may now move on to the difference $`Q_1Q_2`$. After rewriting $`\mathrm{\Pi }_2`$ in terms of $`K`$ and $`\mathrm{\Lambda }_0`$, the coefficient of $`K(\kappa )`$ in Eqn. (36) is $`(q_k+q_\pi /(1+\lambda _2))/2(x_1x_2)^{3/2}`$. This is a rather messy expression in general, but at $`r=s`$, where $`x_1=1`$, its value is $`2(\sqrt{x_2}1/\sqrt{x_2})`$, so that the leading order term in Eqn. (36) is
$`{\displaystyle \frac{\stackrel{~}{A}_Q(rs;s,\varphi )}{r\mathrm{cos}2\varphi }}`$ $`=`$ $`Q_1Q_2(rs)`$ (41)
$``$ $`\left(\sqrt{x_2}{\displaystyle \frac{1}{\sqrt{x_2}}}\right)\mathrm{ln}{\displaystyle \frac{16}{1\kappa ^2}}.`$
One can draw the same conclusion directly from Eqn. (35) by using the useful asymptotic expansions in \[Carlson and Gustafson 1994\], specifically relations (26), (34) and (44) in that paper.
Since the singularity in $`\stackrel{~}{A}_I`$ and $`\stackrel{~}{A}_Q`$ is no worse than logarithmic, we may numerically evaluate integrals involving these by using Gaussian quadrature with a log weight.
## 5 Conclusion
I have obtained analytic angular integrals of the microlensing amplification function, for the case of a rotationally symmetric source. This avoids the need to use approximate methods to obtain this expression, and means that they can be evaluated more efficiently than using general numerical integrations. Also, we are able to make analytic statements about the leading-order behaviour of the integrals along their $`r=s`$ singularity, and so use such asymptotic approximations in further treatments.
This also means that the dependence of the observed flux on the limb-darkening function, and of the observed polarization on the limb-polarization function, can be expressed as integral equations. Thus the problem of recovering the latter from the former can be viewed as a classic inverse problem, which can be analysed in detail using the sophisticated techniques developed for such problems. This is the subject of a forthcoming paper \[Gray and Coleman 2000\]. |
warning/0001/hep-ex0001022.html | ar5iv | text | # Recent Low 𝒙 and Diffractive Collider Data
## 1 Introduction
Since the HERA and Tevatron colliders have been operational, abundant data have become available that are sensitive to proton structure at low parton-$`x`$. Data on photon structure from HERA and LEP have been similarly impressive. This latest generation of colliders has pushed back the limits of our understanding of QCD considerably. There is not space here to do justice to all low $`x`$ data. Instead, three particularly topical areas that were discussed at the 1999 Durham Phenomenology Workshop are singled out.
## 2 Factorisation in Hard Diffraction
The HERA and Tevatron experiments have now produced abundant high quality data on the ‘single diffractive’ processes $`\gamma ^{}pXp`$ and $`\overline{p}pX\stackrel{_{_{()}}}{p}`$ at low momentum transfer. Hard scales, provided for example by a highly virtual photon (in the $`\gamma ^{}p`$ case) or final state jets or heavy quarks, encourage the use of perturbative QCD as a tool with which to understand the parton level dynamics. The development of techniques which simultaneously describe the $`\gamma ^{}p`$ and $`\overline{p}p`$ single dissociation processes is a major current issue in hadron phenomenology.
The generic diffractive process at HERA of the type $`epeXp`$ is illustrated in figure 1a. A photon of virtuality $`Q^2`$ interacts with a proton at a $`\gamma ^{}p`$ invariant mass $`W`$ and squared four momentum transfer $`t`$ to produce a dissociating photon system $`X`$ of invariant masses $`M__X`$, the proton remaining intact. In the corresponding process at the Tevatron, the photon is replaced by an anti-proton, with either of the beam particles dissociating.Two further variables are usually introduced; the fraction of the proton momentum that is exchanged to the system $`X`$ is denoted $`\xi `$<sup>2</sup><sup>2</sup>2At HERA, $`\xi `$ is usually referred to as $`x_{_{IP}}`$. Here $`\xi `$ is used to make explicit the correspondence with the equivalent variable at the Tevatron., whilst $`\beta =x/\xi `$ is the fraction of the exchanged momentum carried by the quark coupling to the photon.
A QCD factorisation theorem has recently been proved for a general class of semi-inclusive processes in deep-inelastic scattering (DIS), which include the single diffractive process . This implies that a concept of ‘diffractive parton distributions’ can be introduced , expressing proton parton probability distributions under the constraint of an intact final state proton with particular values of $`\xi `$ and $`t`$. The cross section for diffractive DIS can then be expressed as
$`\sigma ^{\gamma ^{}pXp}(\xi ,t,x,Q^2){\displaystyle \underset{i}{}}f_{i/p}(\xi ,t,x,Q^2)\widehat{\sigma }_{\gamma ^{}i}(x,Q^2),`$
where $`f_{i/p}(\xi ,t,x,Q^2)`$ are the diffractive parton distributions, evolving with $`x`$ and $`Q^2`$ according to the DGLAP equations at fixed $`\xi `$ and $`t`$, and $`\widehat{\sigma }_{\gamma ^{}i}(x,Q^2)`$ are parton interaction cross sections.
The phenomenology of soft hadronic interactions suggests that it is possible to introduce a universal factorisable pomeron exchange with a flux factor dependent only on $`\xi `$ and $`t`$. With this additional assumption of ‘Regge factorisation’, the framework of diffractive parton distributions can be used to define parton distributions for the pomeron , which should describe all hard diffractive scattering processes. The diffractive DIS cross section can then be written as
$`\sigma ^{\gamma ^{}pXp}(\xi ,t,\beta ,Q^2)f_{\mathrm{I}\mathrm{P}/\mathrm{p}}(\xi ,t){\displaystyle \underset{i}{}}f_{i/\mathrm{I}\mathrm{P}}(\beta ,Q^2)\widehat{\sigma }_{\gamma ^{}i}(\beta ,Q^2).`$
The validity of this second hypothesis for diffractive DIS, incorporating both QCD and Regge factorisation, has been extensively tested at HERA. Measurements of the total cross section for diffractive deep-inelastic scattering, usually presented in the form of a $`t`$-integrated diffractive structure function $`F_2^{D(3)}(\beta ,Q^2,\xi )`$ have shown that, to the present level of accuracy, the factorisation between the $`\xi `$ and the ($`\beta `$,$`Q^2`$) dependence is obeyed.<sup>3</sup><sup>3</sup>3The deviations from this factorisation shown to be present in can be explained in full when a sub-leading exchange ($`f`$, $`\omega `$, $`\rho `$ and / or $`a`$ trajectory) is introduced. Parton distributions for the pomeron have been extracted from QCD analyses of the $`\beta `$ and $`Q^2`$ dependence of $`F_2^D`$ using the DGLAP evolution equations. All such extractions yield parton distributions which are heavily dominated by gluons at low scales, the gluon density remaining large even at high fractional momentum. Figure 1b then represents the dominant process at leading order of QCD. A gluon carrying a fraction $`z`$ of the pomeron momentum undergoes boson-gluon fusion ($`\gamma ^{}gq\overline{q}`$) with the virtual photon.
Monte Carlo models based on the parton distributions extracted from $`F_2^D`$ describe HERA diffractive final state data well . The most stringent tests come from diffractive dijet and open charm cross sections, as both are sensitive to the magnitude as well as the shape of the gluon distribution. Recent data from H1 on diffractive dijet electroproduction are shown in figure 2a. The factorisable partonic pomeron model (labelled “res $`\mathrm{I}\mathrm{P}`$”) gives a reasonable description of the measurement. Similarly good agreement is found with ZEUS data on diffractive charm electroproduction , though a diffractive $`D^{}`$ measurement from H1 in a slightly different kinematic region suggests deviations from factorisation . With this single exception, HERA data support the hypothesis that both QCD and Regge factorisation can be applied to all hard diffractive processes in DIS.
There are good reasons to believe that the QCD factorisation theorem for diffractive DIS cannot be extended to hard diffraction in hadron-hadron interactions . The factorisation hypothesis for $`p\overline{p}`$ scattering has now been tested in some detail by taking parton distributions extracted from $`F_2^D`$ data at HERA and using them to predict cross sections for hard diffractive processes at the Tevatron. By now it is clear that this approach universally predicts cross sections well in excess of those measured.
One example is a measurement of the fraction $`N^{\mathrm{diff}}/N^{\mathrm{incl}}`$ of all dijet events that arise from the single dissociation process $`\overline{p}pXp`$, where the intact final state proton has $`0.035<\xi <0.095`$ and is scattered at $`|t|<1\mathrm{GeV}^2`$. From this ratio, the quantity
$`F_{\mathrm{JJ}}^D={\displaystyle \frac{N^{\mathrm{diff}}}{N^{\mathrm{incl}}}}(x_{\overline{p}})\left\{x_{\overline{p}}g(x_{\overline{p}})+{\displaystyle \frac{4}{9}}\left[q(x_{\overline{p}})+\overline{q}(x_{\overline{p}})\right]\right\}_{\overline{p}}`$
is formed, where $`x_{\overline{p}}`$ is the Bjorken scaling variable for the antiproton and $`x_{\overline{p}}g(x_{\overline{p}})+\frac{4}{9}\left[q(x_{\overline{p}})+\overline{q}(x_{\overline{p}})\right]`$ represent the (known) effective parton densities in the anti-proton after allowing for the leading order colour factor of $`4/9`$. Assuming factorisation is valid, the resulting quantity should correspond to the effective parton densities of the pomeron;
$`F_{\mathrm{JJ}}^D=\left\{\beta g(\beta )+{\displaystyle \frac{4}{9}}\left[q(\beta )+\overline{q}(\beta )\right]\right\}_{\mathrm{I}\mathrm{P}}f_{\mathrm{I}\mathrm{P}/\mathrm{p}}(\xi ).`$
The quantity $`F_{\mathrm{JJ}}^D`$ is shown as a function of $`\beta `$ in figure 2b and is compared with predictions based on parton densities extracted from $`F_2^D`$ by H1 . The predictions are scaled down by a factor of 20, illustrating the size of the factorisation breaking effects. At least for $`\beta \stackrel{<}{_{}}0.3`$, the data and prediction are also rather different in shape.
Another recent measurement from CDF is the fraction of visible beauty production that is attributable to diffraction, which yields the result
$`{\displaystyle \frac{\sigma _{\overline{b}b}^{\mathrm{diff}}}{\sigma _{\overline{b}b}^{\mathrm{incl}}}}=0.62\pm 0.19(\mathrm{stat}.)\pm 0.16(\mathrm{syst}.),`$
whereas the predictions on the basis of diffractive parton densities extracted from $`ep`$ data are at the level of $`10\%`$.
In a complementary analysis, Alvero et al have extracted diffractive parton distributions from $`F_2^D`$ and photoproduction dijet data from HERA and made predictions for various Tevatron measurements. Similarly large discrepancies are found when predicting the rate of $`W`$ and dijet production as components of the system $`X`$ in the process $`\overline{p}pX\stackrel{_{_{()}}}{p}`$ , with even larger differences for dijet production in the double pomeron exchange process $`p\overline{p}pX\overline{p}`$.
Something beyond the simplest Regge and QCD factorisation assumptions is clearly required to describe simultaneously diffractive data from HERA and the Tevatron. The pertinent question now is whether it is possible to build a phenomenological model of this breakdown of factorisation. One possibility is that where beam remnants are present on both sides of a rapidity gap, rescattering takes place, tending to destroy the gap. A very interesting place to study this possibility is in photoproduction, where both factorisable (direct photon) and non-factorisable (resolved photon) interactions may be expected to be present. A first study can be found in .
## 3 Total Cross Sections
Total hadron-hadron cross sections are well described over a very wide energy range by two component Regge fits , corresponding (via the optical theorem) to the exchange of the pomeron and a sub-leading ($`\rho `$, $`\omega `$, $`f`$, $`a`$) trajectory in the elastic amplitude. The intercept of the leading pomeron trajectory is most accurately determined from the high energy rise in the $`\overline{p}p`$ cross section. Other total cross sections such as $`\pi ^\pm p`$ match this scheme well, though no data exist at centre of mass energies $`\stackrel{>}{_{}}30\mathrm{GeV}`$.
In the case where one or both of the interacting hadrons is replaced by a photon, arguments have been made that the presence of a bare photon coupling in addition to the vector meson dominance hadronic component may lead to a faster rise of the total cross section with energy than is the case for pure hadron-hadron scattering. Eikonalised minijet models , incorporating semi-hard QCD interactions whilst avoiding the eventual violation of unitarity associated with simple Regge pole models, can be made to fit the available data . HERA data on the total $`\gamma p`$ cross section at centre of mass energy $`W_{\gamma p}200\mathrm{GeV}`$ are consistent with the simple Regge pole model, though the systematic errors are large and no strong conclusion is yet possible. The increasingly precise data from LEP on the total $`\gamma \gamma `$ cross section, which may be expected to rise faster even than the $`\gamma p`$ cross section, may shed some light on this issue.
Both L3 and OPAL have measured the total $`\gamma \gamma `$ cross section in the region $`10\stackrel{<}{_{}}W_{\gamma \gamma }<100\mathrm{GeV}`$. For the data used, both electrons and many final state hadrons are lost down the beam-pipe, making the kinematics difficult to constrain. The data are shown, together with lower energy fixed target data, in figure 3a. The LEP data clearly show the high energy rise with $`W`$ observed in the $`\overline{p}p`$, $`pp`$ and $`\gamma p`$ cross sections. A simple factorisation law of the type $`\sigma _{\gamma \gamma }=\sigma _{\gamma p}^2/\sigma _{pp}`$ describes the data remarkably well.
It is not yet clear whether the rise with $`W`$ is faster than that observed for total hadron-hadron cross sections; the OPAL data are consistent with the pomeron intercept describing soft hadronic interactions whereas the L3 result is significantly larger. The results are rather sensitive to the assumptions on $`\alpha _{_{\mathrm{I}\mathrm{R}}}(0)`$. As can be seen from figure 3a, a model based on minijets also gives a reasonable description of the data, as do the Schuler and Sjöstrand and PHOJET models, which attempt to make smooth transitions between the photon in its hadronic and point-like manifestations.
Improved data are required before a firm conclusion can be reached concerning the possible anomalous behaviour of $`\sigma _{\gamma \gamma }^{\mathrm{tot}}`$. The main source of error in the measurements arises from the model dependence of the acceptance corrections, with results different at the level of $`20\%`$ obtained when PHOJET or PYTHIA is used for the corrections. The principal reason for this is the different treatments of the diffractive channels in the two models. Any constraints that can be placed on the diffractive processes in $`\gamma \gamma `$ scattering will improve the total cross section measurement considerably. Processes involving the quasi-elastic production of vector mesons are likely to be the easiest to measure, due to the well known decay angular distributions. L3 have taken the first steps towards measurements of the ‘quasi-elastic’ ($`\gamma \gamma \rho ^0\rho ^0`$) and ‘single dissociation’ ($`\gamma \gamma \rho ^0X`$) processes. A measurement of the $`t`$ distribution of the single dissociation process is shown in figure 3b. Fitting the data to the usual exponential parameterisation $`d\sigma /dte^{bt}`$ yields a slope parameter in the region $`b2`$. Information of this sort provides very useful input to soft physics models and should ultimately reduce the model dependence uncertainties on the total cross section.
## 4 Searches for BFKL Dynamics
The BFKL evolution equation, which resums terms where large logarithms of the form $`\mathrm{ln}1/x`$ multiply the coupling constant, must represent a valid approximation to parton dynamics in some region of low $`x`$ phase space. The search for evidence for BFKL behaviour is one of the principle current experimental activities in low $`x`$ physics. Although it has been shown that introducing BFKL effects can improve the description of $`F_2`$ at low $`x`$ , it is not yet accepted that anything more than standard DGLAP evolution is required to describe current inclusive DIS data. Exclusive final state measurements may ultimately produce the clearest BFKL signatures. Some of the more promising areas of study are discussed below.
BFKL and DGLAP evolution have rather different implications for the details of the parton ladder governing low $`x`$ DIS processes (figure 1c). In the DGLAP case, one expects an ordering in virtuality ($`k_t`$) of the partons in the ladder, leading to rapidity ordering of the transverse momenta of outgoing partons. The BFKL scheme has no such strong ordering and therefore results in anomalously large high $`p__T`$ hadron yields away from the photon vertex, for example at central rapidity. The central region of the $`\gamma ^{}p`$ frame corresponds to the forward region of laboratory rapidity.
Both H1 and ZEUS have studied the production of jets in this difficult forward region . H1 have also measured the cross section for high $`p__T`$ forward $`\pi ^0`$ production . Similar conclusions are reached in each case. The ZEUS forward jet data are shown in figure 4a. The data cannot be described by standard DGLAP models (labelled LEPTO and HERWIG). Only models that do not impose strong transverse momentum ordering are able to describe the data. One example is the ARIADNE model , based on the colour dipole model and simulating BFKL ordering. However, the lack of strong $`k__T`$ ordering can also be modelled through the introduction of partonic structure to the virtual photon. This can be implemented in the RAPGAP Monte Carlo model, giving a successful description of all forward region data produced at HERA to date. Thus the final state data from the forward region at HERA demonstrate that something more than the simplest DGLAP model of the low-$`x`$ parton ladder is required. However, resolved virtual photons provide an alternative mechanism to BFKL to restore a good description of the data. Work on events with large rapidity separations between pairs of jets, just beginning at the Tevatron, may help to resolve some of these ambiguities.
In appropriate kinematic regions, total, elastic and diffractive cross sections may all be describable in terms of the amplitude for elastic parton-parton scattering via the exchange of gluon ladders, evolving according to BFKL dynamics. BFKL calculations are most reliable where large scales are present at both vertices . One example is the total $`\gamma ^{}\gamma ^{}`$ cross section . Where both photons have sufficiently high virtuality, measurement conditions at LEP are favourable and first data have appeared . The data suggest a relatively strong energy dependence, which may be consistent with BFKL predictions. However the present data can be described equally well by non-BFKL QCD models, for example those involving virtual photon structure .
Another process where large scales are present at both vertices is diffractive scattering at large $`|t|`$, where precision data are starting to appear from HERA and the Tevatron. The quasi elastic process $`\gamma pVY`$ where $`V`$ denotes a vector meson and $`Y`$ is a proton or low mass proton excitation has been measured for $`V=J/\psi `$, $`\rho `$ and $`\varphi `$ . The results in the relatively low $`|t|`$ regions accessed to date are mixed, only the $`J/\psi `$ fully conforming to the BFKL predictions.
The classic high $`|t|`$ diffractive process is the production of dijets separated by a rapidity gap, implying a net colour singlet exchange. Here, the magnitude of $`t`$ is close to the jet $`E__T^2`$ and is thus very large. The size of the cross section is usually quantified as the fraction of all dijet events that have a rapidity gap between the jets. Clear signals have been observed at large jet pseudorapidity separation $`\mathrm{\Delta }\eta `$ both in photoproduction at HERA and in $`p\overline{p}`$ interactions at the Tevatron . The gap fraction at large $`\mathrm{\Delta }\eta `$ decreases with centre of mass energy, being around 0.1 at $`\sqrt{s}200\mathrm{GeV}`$ at HERA, 0.025 at $`\sqrt{s}=630\mathrm{GeV}`$ at the Tevatron and 0.01 at $`\sqrt{s}=1800\mathrm{GeV}`$ at the Tevatron. This trend is opposite to that naively expected from BFKL calculations. However, it seems likely that rapidity gap destruction due to reinteractions of beam remnants plays an important role. Two very different models of these effects both predict a rapidity gap survival probability that falls with centre of mass energy in a manner that qualitatively resembles that in the data.
The Tevatron gap fractions have been measured as a function of jet $`E__T`$ as well as $`\mathrm{\Delta }\eta `$ (see figure 4b). The gap fraction is found to be flat or slowly rising with $`E__T^{\mathrm{jet}}`$, which matches predictions based on the creation of rapidity gaps by soft colour interactions in otherwise standard dijet events . It has been demonstrated that if rapidity gap destruction effects are included, BFKL dynamics can also describe these data .
All of the measurements discussed above can be interpreted in terms of BFKL effects, yet none conclusively demonstrates the need for BFKL at present colliders. Data from the upgraded Tevatron and HERA may allow us to resolve this question.
## References |
warning/0001/hep-th0001066.html | ar5iv | text | # Untitled Document
hep-th/0001066 TIFR/TH/00-02 KCL-TH-00-02
Brane-Antibrane Constructions
Sunil Mukhi<sup>1</sup>E-mail: mukhi@tifr.res.in, nemani@tifr.res.in, tong@mth.kcl.ac.uk, Nemani V. Suryanarayana<sup>1</sup> and David Tong<sup>2</sup>
<sup>1</sup> Tata Institute of Fundamental Research,
Homi Bhabha Rd, Mumbai 400 005, India
<sup>2</sup> Department of Mathematics, Kings College,
The Strand, London, WC2R 2LS, UK
ABSTRACT
In type II string theories, we examine intersecting brane constructions containing brane-antibrane pairs suspended between 5-branes, and more general non-BPS constructions. The tree-level spectra are obtained in each case. We identify various models with distinct physics: parallel brane-antibrane pairs, adjacent pairs, non-adjacent pairs, and configurations which break all supersymmetry even though any pair of branes preserves some supersymmetry. In each case we examine the possible decay modes. Some of these configurations turn out to be tachyon-free, stable non-BPS states. We use T-duality to map some of our brane constructions to brane-antibrane pairs at ALE singularities. This enables us to explicitly derive the spectra by the analogue of the quiver construction, and to compute the sign of the brane-antibrane force in each case.
January 2000
1. Introduction
The study of unstable branes in type II superstring theories has made considerable progress over the last two years\[1\]. The relevant unstable branes are of two types: D$`p`$-brane-antibrane pairs, ($`p`$ even for type IIA, odd for type IIB) and unstable D$`p`$-branes with $`p`$ odd for type IIA, even for type IIB. In both cases, the instability is signalled by the presence of a tachyon, and it is possible to identify a variety of decay modes. These assemble themselves into interesting sequences that terminate with stable BPS D-branes.
The unstable branes and their decay modes form a beautiful and fundamental structure, which has been interpreted in terms of K-theory\[2,,3\]. This structure can then be used to study more complicated situations such as orientifolds, orbifolds and Calabi-Yau compactifications. In these cases it often happens that one discovers novel stable non-BPS states (see Refs.\[1,,4,,5\] and references therein).
In the present paper, we attempt to generalise these elegant discoveries to situations where unstable D-branes are suspended between other branes. As we will see, along with many phenomena that are familiar from the study of infinite or toroidally compactified unstable branes, there are also new constraints and novel physical situations that have no counterpart in the simpler models.
All the brane constructions that we study will have completely broken supersymmetry. However, they are states of type II string theories, and therefore are endowed with special properties that arise from the fact that the underlying theory is supersymmetric. One interesting class of models that we define has broken supersymmetry despite the absence of brane-antibrane pairs or single unstable branes. In these models one arranges at least three different types of branes together, in such a way that each pair preserves some supersymmetry, but all the branes together break all supersymmetry.
Even in supersymmetric models, the study of suspended branes suffers from some uncertainties, as it is often hard to reliably extract the spectrum of light states. These uncertainties are all the more severe in the present case, as the brane configurations lack supersymmetry which would have classified the possible states into multiplets. As most of the models we study involve D-branes suspended between NS5-branes, we will find it useful to dualise away the NS5-branes into ALE geometry\[6\] and then study the model using quiver techniques\[7\], where many of the relevant quantities can in fact be reliably computed.
Although this paper will deal for the most part with unstable configurations, it is meaningful to study their spectra at weak coupling by working at string tree-level. Here one can identify the presence or absence of potential decay modes related to tachyon condensation, even though the configuration actually gets destabilised after loop effects are taken into account. All discussions of unstable configurations should be understood in the light of this comment.
The plan of this paper is as follows. In Section 2, we describe systems of parallel D-brane-antibrane pairs suspended between NS5-branes and D5-branes. These are closest in behaviour to the noncompact parallel brane-antibrane pairs, though we identify some differences that arise due to boundary conditions at the ends of the D-branes. In Section 3 we look at unstable uncharged D-branes suspended between NS5-branes and examine some possible decay modes. In Section 4 we consider systems of parallel NS5-branes with a D-brane stretched across one segment and a $`\overline{\mathrm{D}}`$-brane stretched across the adjacent segment. In this case the pair cannot annihilate. We examine some related models and argue that in general these pairs will repel each other. We also consider configurations with more NS5-branes, where the brane-antibrane pair is non-adjacent. In Section 5 we introduce the “Borromean branes”, configurations of branes in which supersymmetry is broken by all the branes together but not by any given pair of branes. This has the interesting consequence that there is no perturbative tachyon, and the configuration can potentially be stable. In Section 6 we describe the duality which relates suspended branes between NS5-branes to fractional branes at ALE singularities. Although this duality has been described and used before in the literature, we give a slightly different and very explicit derivation, which will hopefully make it somewhat clearer. In Section 7 we apply this duality to analyse parallel and adjacent brane-antibrane pairs from the point of view of quiver theory. In Section 8 we compute the spectra of open strings in these configurations by constructing boundary states for the relevant fractional branes. We also use this formalism to compute the forces between different pairs of fractional branes, confirming some of the speculations made in earlier sections.
We will not provide a review of various fundamental aspects of brane-antibrane dynamics that will be made use of in this paper. For this, the reader is advised to consult Refs.\[1,,4,,5\].
2. Suspended Parallel Brane-antibrane Pairs
Consider, in type IIA string theory, a model with a pair of NS5-branes extending along the directions $`(x^1,x^2,x^3,x^4,x^5)`$, and located at $`(x^6,x^7,x^8,x^9)=(0,0,0,0)`$ and $`(L_6,0,0,0)`$ respectively. Thus they are parallel and separated by a finite distance $`L_6`$ along $`x^6`$. Between these, we suspend a D4-brane and a $`\overline{\mathrm{D}}4`$-brane along the $`x^6`$ direction. They extend along the directions $`(x^1,x^2,x^3)`$ and can be separated from each other along $`(x^4,x^5)`$ (Fig.(2.1)). Without the antibrane, this model would belong to the class of brane constructions analysed in Ref.\[8\].
Fig.(2.1): A $`\mathrm{D}4\overline{\mathrm{D}}4`$ pair suspended between parallel NS5-branes in type IIA.
We expect that the common 3+1-dimensional world-volume supports a non-supersymmetric field theory. The spectrum of light states, including possible tachyons, can be deduced as follows. First consider an infinitely extended D4-$`\overline{\mathrm{D}}4`$ pair. As is well-known\[1\], the open-string states have four sectors, corresponding to the Chan-Paton factors:
$$a=\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right),b=\left(\begin{array}{cc}0& 0\\ 0& 1\end{array}\right),c=\left(\begin{array}{cc}0& 1\\ 0& 0\end{array}\right),d=\left(\begin{array}{cc}0& 0\\ 1& 0\end{array}\right)$$
The first two sectors come from strings beginning and ending on the same brane(antibrane). Together, these give a GSO-projected spectrum of $`U(1)\times U(1)`$ vector multiplets, containing two 4+1-dimensional gauge fields $`A_\mu ^{(I)}`$, two sets of five massless scalar fields $`X_i^{(I)}`$ and two sets of 4 massless Majorana fermions $`\chi _A^{(I)}`$, where $`i=1,\mathrm{},5`$, $`A=1,\mathrm{},4`$ and $`I=1,2`$. The index $`I`$ here labels the Chan-Paton factors $`a,b`$. In the other two sectors, associated to open strings going from the brane to the antibrane and vice-versa, there is a complex tachyon $`T`$ carrying charges $`(1,1)`$ under $`U(1)\times U(1)`$, and another set of massless Majorana fermions $`\lambda _A^{(I)}`$. (This time the index $`I`$ labels the Chan-Paton factors $`c,d`$.) These correspond to “anti-GSO” states that arise because of the fact that the two branes carry opposite RR charges.
Together, the infinite brane-antibrane pair breaks all 32 supersymmetries of the bulk theory. Hence we expect to find 32 massless Goldstone fermions (“goldstinoes”) on the world-volume of this pair. Actually we have 64 massless fermions, from the $`\chi _A^{(I)}`$ and $`\lambda _A^{(I)}`$. One way to identify which of these are goldstinoes is to use the fact that when we quotient type IIA string theory, with a $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair, by the symmetry $`(1)^{F_L}`$, we end up with type IIB string theory with an unstable D4-brane in it. This unstable brane carries a real tachyon and precisely 32 massless Majorana fermions. In this case it is clear that all the fermions must be goldstinoes.
It follows that in the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair, the goldstinoes must be those fermions which survive the $`(1)^{F_L}`$ quotient ($`(1)^{F_L}`$ has no twisted sectors on open-string states, so those fermions which survive the quotient and are goldstinoes in the final theory must have been goldstinoes before taking the quotient too). This quotient acts on the Chan-Paton factors as conjugation by $`\sigma _1`$, hence the surviving Chan-Paton factors on the unstable D4-brane are $`1,\sigma _1`$. It follows that $`\chi _A^+=\chi _A^{(1)}+\chi _A^{(2)}`$ and $`\lambda _A^+=\lambda _A^{(1)}+\lambda _A^{(2)}`$ are the goldstinoes on the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair.
Now let us go back to the model of a $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair suspended between NS5-branes and work out the spectrum. Again we have a $`U(1)\times U(1)`$ gauge theory, whose spectrum is a truncation of that on the infinite $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair. The truncation arises from the boundary conditions at the two ends of the 4-branes.
For the massless fields in the GSO-projected sector (associated to Chan-Paton factors $`1,\sigma _3`$) the truncation is well-known: the massless scalars that are collective coordinates for broken translation invariance are projected out whenever they correspond to directions along which the 4-branes cannot move. Since the NS5-branes on which they end are located at fixed values of $`x^7,x^8,x^9`$, it follows that the scalars $`X_7^{(I)},X_8^{(I)},X_9^{(I)}`$ are projected out. Along with these, the gauge field component $`A_6^{(I)}`$ and half the associated fermions, say $`\chi _3^{(I)},\chi _4^{(I)}`$ are projected out. The result is a set of massless fields: $`A_\mu ^{(I)},X_i^{(I)},\chi _A^{(I)}`$ where $`\mu =0,\mathrm{},3`$, $`i=4,5`$ and $`A=1,2`$. This is a pair of vector multiplets of $`𝒩=2`$ supersymmetry in 3+1 dimensions.
In the anti-GSO sector, supersymmetry is clearly not available to guide us in finding the spectrum that survives. We shall work in the limit in which the separation between NS5-branes is very much smaller than the string length scale, ensuring that the higher tachyon Kaluza-Klein modes have positive mass-squared and so, as for other fields, we restrict attention to the constant modes. It is plausible, and we will confirm this later, that the complex tachyon survives the boundary conditions. Recall that for an infinite $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair, the tachyon is associated to the NS sector ground state for open strings stretching between the brane and the antibrane. Thus we are claiming that this ground state is not projected out by the boundary conditions. It is physically evident that the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair suspended between parallel NS5-branes should be unstable, just as an infinite, or wrapped, $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair would be. Hence we expect it to contain a tachyon. Another fact that supports this conclusion is that the $`U(1)\times U(1)`$ gauge fields under which the tachyon is charged do survive the boundary conditions, as we have seen. In a later section we will show explicitly in a T-dual version of this model that the tachyon is indeed present.
Finally we turn to the anti-GSO sector fermions. In the case of infinite $`\mathrm{D4}\overline{\mathrm{D}}4`$ pairs, we had altogether 32 anti-GSO fermions $`\lambda _A^{(I)}`$, of which 16 were goldstinoes. Now on the suspended $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair we expect to find altogether 16 goldstinoes (because they break all 16 of the supersymmetries that are preserved by the NS5-branes). But so far we have found 8 goldstinoes $`\chi _A^+`$, $`A=1,2`$ and another 8 non-goldstino fermions $`\chi _A^{}=\chi _A^{(1)}\chi _A^{(2)}`$, $`A=1,2`$. The remaining 8 goldstinoes must therefore be $`\lambda _A^+`$, $`A=1,2`$. The symmetry between $`\lambda _A^{(1)}`$ and $`\lambda _A^{(2)}`$ then suggests that the combinations $`\lambda _A^{}=\lambda _A^{(1)}\lambda _A^{(2)}`$, $`A=1,2`$ also survive as (non-goldstino) massless fermions.
To summarize, we have a tachyonic $`U(1)\times U(1)`$ gauge theory on the world-volume of the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair, with a pair of gauge fields $`A_\mu ^{(I)}`$, massless scalars $`X_i^{(I)}`$ and massless fermions $`\chi _A^\pm ,\lambda _A^\pm `$ of which the $`+`$ superscripts correspond to goldstinoes. Finally there is a complex tachyon $`T=T_1+iT_2`$.
Now we are in a position to ask how this configuration can decay. The first and most elementary process is that the tachyon can go to its minimum. In this case, the brane-antibrane pair annihilates completely. While one of the $`U(1)`$ gauge symmetries gets Higgsed in the process, the other $`U(1)`$ gets “confined” by the mechanism discussed in Ref.\[9\]. At the same time, the value of the potential at the tachyonic minimum, $`V(T_0)`$, is expected to cancel the tension of the annihilating branes. Thus we end up with a BPS configuration consisting of just a pair of parallel NS5-branes. Note that both the brane tensions and the tachyon potential scale by a common factor of $`L_6`$, the separation between the NS5-branes, hence apart from this overall scale factor, we expect $`V(T_0)`$ to be independent of the coupling constant of the 3+1 dimensional field theory. This is consistent with a recent analysis showing that upto such a factor the tachyon potential is universal, independent of the background\[10\].
A less elementary decay mode would be condensation of a tachyonic vortex\[11,,12\]. For infinite or wrapped $`\mathrm{D}p\overline{\mathrm{D}}p`$ pairs, the complex tachyon can condense in a topologically stable vortex, while the relative $`U(1)`$ gauge field under which it is charged carries a unit of magnetic flux. The result is a stable BPS D$`(p2)`$ brane. Looking at the field content of the world-volume theory on the suspended $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair, we see that the same configuration is allowed here, providing the resulting D2-brane is extended in the $`x^6`$ direction. The configuration of Fig.(2.1) can therefore decay into a configuration of a BPS D2-brane suspended between parallel NS5-branes.
The picture of this process on the world-volume of the NS5-branes is interesting. A D4-brane ends as a 3-brane in the NS5-brane world-volume. This 3-brane behaves as a vortex since it has co-dimension 2. The $`\overline{\mathrm{D}}4`$ brane similarly behaves as an antivortex. The process where the tachyon goes to its minimum corresponds to simple annihilation of the vortex-antivortex pair. On the other hand, the more complex process in which the tachyon condenses as a vortex, corresponds to the 3-brane anti-3-brane vortices annihilating into a 1-brane<sup>1</sup> This discussion may be a little confusing since there are two types of vortices involved. The tachyonic vortex is localized in, say, the $`(x^1,x^2)`$ directions, while the 3-brane vortices in the NS5-branes are localized in the $`(x^4,x^5)`$ directions.. This 1-brane in the NS5 world-volume is just the end of a D2-brane.
Next let us examine a decay mode which would be allowed for infinite or wrapped $`\mathrm{D4}\overline{\mathrm{D}}4`$ pairs but turns out to be forbidden in the present situation. It has been noted\[13\] that any $`\mathrm{D}p\overline{\mathrm{D}}p`$ pair can also decay by creating an electric flux on its world-volume. In this case, the decay product is a fundamental string (F-string). If this were possible, the configuration of Fig.(2.1) would decay into a configuration where an F-string is suspended between two NS5-branes. Clearly, this is impossible since F-strings cannot end on NS5-branes. Hence we must show that such an electric flux is prohibited in the world-volume theory discussed above for Fig.(2.1).
This follows from the fact that the boundary conditions imposed by NS5-branes project out the component $`A_6^{(I)}`$ from each of the two 5-dimensional gauge fields. On $`\mathrm{D4}\overline{\mathrm{D}}4`$ pairs without boundaries, this component would dualize to form part of a magnetic 2-form gauge potential, under which tachyonic strings arising from stretched D2-branes are charged (although this phenomenon is, of course, inherently non-perturbative). These magnetic tachyons could then condense, forming an electric world-volume flux\[13\] (the dual of the usual vortex condensation where an electric tachyon condenses giving rise to a magnetic flux). This process corresponds to the annihilation of the pair into a fundamental string. Once the NS5-branes are put in as boundaries, components of the magnetic 3-form field strength $`H_{\mu \nu \rho }`$ with $`\mu ,\nu ,\rho 6`$ are lost, along with the possibility of an electric flux in the $`x^6`$ direction. However, one is still left with the decay mode in which the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pairs decay into a fundamental string oriented parallel to the NS5-branes in either the $`x^1,x^2`$ or $`x^3`$ directions. As with similar decays in the case of infinitely extended branes, the resulting configuration preserves some fraction of supersymmetry which, in the present situation, is $`1/4`$ of the original 32 supercharges.
The fact that the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair cannot decay into a fundamental string stretched between NS5-branes can also be understood in terms of the theory on the NS5-brane worldvolume. In this language we start with a 3-brane vortex-antivortex pair. This cannot annihilate into a point-like object (representing the end of a fundamental string), since a fundamental string carries $`B_{NS,NS}`$ charge and therefore its endpoint must be an electrically charged particle – but the NS5-brane of type IIA string theory does not carry 1-form gauge fields, and therefore it does not support electrical charges.
For our next example, consider a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair suspended between NS5-branes in type IIB string theory (Fig.(2.2)). The system is very similar to the one above. The NS5-branes have the same locations as before, while the 3-branes extend along $`(x^1,x^2,x^6)`$ and are located at definite positions in $`(x^3,x^4,x^5)`$ and at the origin in $`(x^7,x^8,x^9)`$. Again, without the antibrane this model would be a familiar one — it belongs to the class of brane constructions analysed in Ref.\[14\].
Fig.(2.2): A $`\mathrm{D}3\overline{\mathrm{D}}3`$ pair suspended between parallel NS5-branes in type IIB.
The world-volume theory now lives in 2+1 dimensions. The relevant fields from the GSO-projected sector are a pair of $`U(1)`$ gauge fields, a set of scalar triplets $`X_i^{(I)}`$, $`i=3,4,5`$, and a set of massless fermions $`\chi _A^{(I)}`$, $`A=1,\mathrm{},4`$. Note that this time the fermions are two-component spinors. From the anti-GSO sector we again expect a complex tachyon $`T`$ and a set of massless fermions $`\lambda _A^{(I)}`$, $`A=1,\mathrm{},4`$. In this model, too, the possibilities are that the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair annihilates completely (with the tachyon going to the minimum of the potential), or there is a tachyonic vortex resulting in annihilation of the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair into a D-string stretched in the $`x^6`$ direction. As in the model of Fig.(2.1), the loss of $`A_6`$ due to boundary conditions ensures the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair can only decay into an F-string that lies parallel to the NS5-branes.
In this model, it is natural to wonder whether more general decays to (p,q)-strings occur. S-duality ensures that such a decay is allowed in the case of infinitely extended $`\mathrm{D3}\overline{\mathrm{D}}3`$ pairs, with the electric and magnetic tachyons winding p and q times respectively before condensing. Notice that the restrictions on the stable decay products discussed above are precisely those of supersymmetry. The corresponding supersymmetry restriction for a (p,q)-string is a configuration in which the string “kinks” between the two NS5-branes\[15\]. Unfortunately, we do not have enough understanding of the non-perturbative magnetic tachyon dynamics to say whether such a decay actually occurs.
The third model that we want to consider is quite different. Here we have a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair suspended between two D5-branes (Fig.(2.3)).
Fig.(2.3): A $`\mathrm{D}3\overline{\mathrm{D}}3`$ pair suspended between parallel D5-branes in type IIB.
In this case, the gauge field component $`A_6`$ is retained while the components $`A_\mu `$, $`\mu =0,1,2`$ are projected out\[14\]. The scalar $`A_6`$ may be shifted by an integer multiple of $`2\pi /L_6`$ by performing a “large” gauge transformation, which acts upon any charged fields by shifting each Kaluza-Klein mode into the next. Thus the scalar $`A_6`$ becomes periodic<sup>2</sup> A well-known example of this fact is that the moduli space of a D-string stretched between D3’s is the moduli space of a single monopole, which is $`R^3\times S^1`$, with periodicity $`2\pi /L`$. and, as a result, this model has a magnetic rather than electric gauge group. Although the model of Fig.(2.3) is the S-dual of the model of Fig.(2.2), we cannot make any definite use of this fact since we are dealing with unstable non-BPS brane configurations. S-duality relates Fig.(2.2) at weak coupling to Fig.(2.3) at strong coupling and vice-versa, while we are interested in both models at weak coupling. The dynamics of the model of Fig.(2.3) at strong coupling is rather clear: there is a magnetic tachyon, corresponding to quantization of a D-string connecting the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair, and this can condense with an electric vortex to give an F-string suspended between the D5-branes. This process is S-dual to the condensation of an electric tachyon via a magnetic vortex in the model of Fig.(2.2), and the resulting configuration (which is BPS) is likewise S-dual to a D-string suspended between parallel NS5-branes.
The first statement that we can make about the model of Fig.(2.3) at weak coupling is that the boundary conditions imposed by the D5-branes on the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair must forbid the decay of this pair into a D-string stretched between the two D5-branes. It is interesting to consider all the available possibilities that lead to this conclusion.
One possibility might be that the complex tachyon is projected out along with the electric gauge field. In this case the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair suspended between two D5-branes would be stable for sufficiently close D5-branes (since it is the constant mode of the tachyon that suffers projection, one could still have a tachyonic instability for well-separated D5-branes). This would mean that suspending a $`\mathrm{D3}\overline{\mathrm{D}}3`$ between D5-branes gives rise to a stable, non-BPS state for some values of the separations. However, the present system is T-dual to a D-string-anti-D-string pair suspended between parallel D3-branes. In that system, the end-point of the D-string behaves as a magnetic monopole in the D3-brane world-volume, so the string-antistring pair is really like a monopole-antimonopole pair. Such a system can certainly annihilate to the vacuum and we therefore expect a tachyon to be present in the spectrum.
A second possibility is that, although the electric gauge field under which the tachyon was charged is projected out by the boundary conditions at the D5-brane, the complex tachyon itself survives. In this scenario, the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair can annihilate into the vacuum by having the tachyon go to a constant minimum as usual. However, there also exists the possibility of the tachyon winding before condensing, potentially forming the forbidden configuration of a D-string stretched between D5-branes. One might hope that the resulting vortex solution has divergent energy since finite-energy vortices cannot exist in standard field theories without a gauge field. In this case, we would not be able to identify the resulting configuration as a D-string. Moreover, the D-string charge coming from the Chern-Simons coupling $`FC_{RR}^{(2)}`$ would also be absent without a gauge field.
However, there are some reasons to be dissatisfied with this scenario. Finite-energy vortices without gauge fields are possible if one allows higher-derivative couplings in the field theory, which are certainly present because of stringy corrections. Moreover, besides the Chern-Simons term described above, there is another proposed coupling\[16\] of the form $`d(TDT^{})C_{RR}^{(2)}`$ on a brane-antibrane pair, which could give rise to the induced D-string charge without the help of a gauge field. Finally, it is believed\[1\] that a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair can decay by condensing a real unstable tachyonic kink, into an unstable D2-brane. In turn, this can decay, by condensing a real stable tachyonic kink, into a BPS D-string. This two-step process could again lead to the forbidden D-string, so (barring some mechanism that we do not presently understand that inhibits one of the steps), the hypothesis that the complex tachyon survives seems unlikely.
The final possibility is that one real component of the tachyon is projected out, along with the electric gauge field. The surviving field content of the suspended $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair would then be very similar to that of a single unstable D-brane: a U(1) gauge field and scalars, and a neutral real tachyon. This has the advantage that the pair still has two decay modes, out of the original three: it can decay into the vacuum, by condensation of this real tachyon into its constant minimum, or it can decay into the unstable D2-brane suspended between D5-branes by condensation of a real tachyonic kink. Only vortex condensation (whether in one step or two steps) is ruled out since there is only one tachyon.
This possibility seems the most elegant, because it inhibits only the decay mode that is definitely forbidden on grounds of charge conservation, into a D-string. Moreover, it suggests that an unstable D2-brane of type IIB suspended between D5-branes is stable. We will return to this point in the next section.
There is one remaining decay mode to consider in which the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair decays into a D-string lying parallel to the D5-branes. Such a decay occurs if either a real or complex constant tachyonic mode survives and the topology required to realise these strings as solitonic solutions lies in the periodicity of $`A_6`$. The 3-dimensional field theory on the D3-branes includes the coupling $`|A_6|^2|T|^2`$ which ensures that when the tachyon condenses, the vacuum lies at $`A_6=0`$. However, $`A_6`$ may wind as, say, $`x_1`$ ranges from $`\mathrm{}`$ to $`+\mathrm{}`$, resulting in a string like configuration stretched in the $`x_2`$ direction. This configuration has non-zero flux $`F_{16}`$ and so, by the usual Chern-Simons term, $`FC_{RR}^{(2)}`$, is identified as a D-string.
There is another way to see the appearance of strings in this case<sup>3</sup> We thank Kimyeong Lee and Piljin Yi for explanations of this point. that dates back to the work of Polyakov \[17\]. The periodic scalar, $`A_6`$, may be dualised in favour of a 3-dimensional gauge field, $`\stackrel{~}{A}`$. The tachyon couples magnetically to $`\stackrel{~}{A}`$ and, by the dual-Meissner effect, condensation of the tachyon leads to a vacuum in which objects charged electrically under $`\stackrel{~}{A}`$ are linearly confined. The electric flux lines associated with this confinement are then identified with the kinks described above. Similar systems, in which confining flux lines have a description in terms of classical solitons, have been considered recently in the supersymmetric context in \[18\].
For our final model, consider a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair suspended between an NS5-brane at one end and a D5-brane at the other. The NS5-brane fills the directions $`(x^1,x^2,x^3,x^4,x^5)`$ and is located at $`(x^6,x^7,x^8,x^9)=(0,0,0,0)`$, while the D5-brane fills $`(x^1,x^2,x^7,x^8,x^9)`$ and is located at $`(x^3,x^4,x^5,x^6)=(0,0,0,L_6)`$. The $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair fills the directions $`(x^1,x^2)`$, stretches from $`x^6=0`$ to $`x^6=L_6`$ and is located at $`(x^3,x^4,x^5,x^7,x^8,x^9)=(0,0,0,0,0,0)`$ (Fig.(2.4)).
Fig.(2.4): A $`\mathrm{D}3\overline{\mathrm{D}}3`$ pair suspended between an NS5 and a D5-brane in type IIB.
From Ref.\[14\], we know that the massless states in the GSO-projected sector, corresponding to open strings on the brane or antibrane alone, get projected out at one or the other end so there is neither an electric nor a magnetic gauge group. For the anti-GSO sector, we again expect, as above, that either a real tachyon or the whole complex tachyon survives at the D5 end. Whichever of these two possibilities is realised will govern the dynamics of the system, since the whole complex tachyon survives at the NS5-brane end. Note that this model is self-dual under S-duality, so its strong-coupling behaviour will be the same as its weak-coupling behaviour with the roles of the two bounding 5-branes exchanged.
3. Suspended Unstable D-branes
The type IIA string has unstable, non-BPS D$`p`$-branes for $`p=1,3,5,7,9`$\[1\]. Hence one can consider a brane construction in which one of these is suspended between NS5-branes. For example, let us start with an unstable D3-brane suspended in this way (Fig.(3.1)).
Fig.(3.1): An unstable D3-brane suspended between parallel NS5-branes in type IIA.
The directions of the branes are exactly as for the model in Fig.(2.2). Indeed, the configuration of Fig.(3.1) can be obtained by starting with Fig.(2.2) and taking the quotient of the configuration and the bulk theory together by $`(1)^{F_L}`$. This is possible since the brane construction in Fig.(2.2) is invariant under $`(1)^{F_L}`$, which in turn holds since $`(1)^{F_L}`$ interchanges a D3-brane with a $`\overline{\mathrm{D}}3`$-brane in type IIB, and preserves the NS5-brane.
The spectrum on the world-volume theory of Fig.(3.1) is a truncation of that on a single unstable D3-brane of type IIA. This unstable brane, if it is infinite or wrapped, has a single gauge field $`A_\mu `$, a set of massless scalars $`X_i`$ and a set of massless fermions $`\chi _A`$, where $`\mu =0,\mathrm{},3`$, $`i=4,\mathrm{},9`$ and $`A=1,\mathrm{},4`$. This is a GSO-projected spectrum and arises in the sector with identity Chan-Paton factor. In addition, there is another sector, with Chan-Paton factor $`\sigma _1`$ and anti-GSO projection, where one finds a real tachyon and another set of massless fermions $`\lambda _A`$, $`A=1,\mathrm{},4`$. This time, all the 32 fermions $`\chi _A`$ and $`\lambda _A`$ are goldstinoes of spontaneously broken supersymmetry. Once the unstable brane is bounded by NS5-branes, it breaks only 16 supersymmetries (the other 16 are already broken by the NS5-branes which act as a “background” from the point of view of the D3-brane). Thus we expect 16 massless fermions $`\chi _A`$ and $`\lambda _A`$ where now $`A=1,2`$ and all these fermions are goldstinoes. The rest of the light spectrum is made up by the gauge field $`A_\mu `$ and three scalars, $`X_i`$, $`i=3,4,5`$, along with a real tachyon $`T`$.
There is another way to obtain the configuration of Fig.(3.1): start with the configuration of Fig.(1.1), in type IIA theory, and allow the real part of the (complex) tachyon on the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair to develop a kink. This is not a topologically stable solution, and therefore does not cause the $`\mathrm{D4}\overline{\mathrm{D}}4`$ to condense to a stable object. Instead it creates an unstable D3-brane of type IIA at the point where the kink is located\[1\]. To get the directions appropriate to Fig.(3.1), the kink in Fig.(1.1) must be along the $`x^3`$ direction.
Since one way of obtaining Fig.(3.1) was to quotient the configuration of Fig.(2.2) by $`(1)^{F_L}`$, one is tempted to ask whether a similar quotient can be carried out on Figs.(2.1) and (2.3). For Fig.(2.1) this is indeed possible, and one ends up with the unstable D4-brane of type IIB suspended between NS5-branes (Fig.(3.2)). However, for Fig.(2.3) the story is quite different. The action of $`(1)^{F_L}`$ does not preserve D5-branes, hence it is not a symmetry of Fig.(2.3) and one cannot quotient by it.
Fig. (3.2): An unstable D4-brane suspended between parallel NS5-branes in type IIB.
Finally one can ask the reverse question: what happens when we quotient Fig.(3.1) by $`(1)^{F_L}`$? It is known that this action takes the unstable D3-brane of type IIA to the stable BPS D3-brane of type IIB. Hence it is reasonable to expect that quotienting Fig.(3.1) by $`(1)^{F_L}`$ gives us a single BPS D3-brane stretched between NS5-branes in type IIB theory. Likewise, we could take this quotient on Fig.(1.1) and we would end up with a stable D4-brane of type IIA suspended between two NS5-branes. Thus, the stable BPS configurations analysed in Refs.\[8,,14\] arise as $`(1)^{F_L}`$ quotients of the unstable configurations considered in this section.
Since unstable branes do not carry charges, one may ask if they can always be consistently suspended between any pair of stable branes. Consider a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair suspended between NS5-branes in type IIB. Clearly, condensation of a real tachyonic kink can give rise to an unstable D2-brane suspended between the NS5-branes. The tachyon that condenses is the real part of an electric tachyon on the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair, which as we have argued, survives the boundary conditions imposed by the NS5-brane. Thus it must be consistent to have an unstable D2-brane suspended between NS5-branes, analogous to the system in Fig.(3.2). (This configuration can also be produced by starting with a $`\mathrm{D2}\overline{\mathrm{D}}2`$ pair in type IIA, suspended between NS5-branes, and quotienting by $`(1)^{F_L}`$.) Moreover, the tachyon on this D2-brane, although uncharged, is actually the imaginary part of the electric tachyon on the original $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair. Thus, this tachyon is also not projected out by the NS5-brane, and it can undergo further kink condensation leading to a stable BPS configuration, a D-string suspended between the NS5-branes.
If we replace the NS5-branes by D5-branes, we certainly cannot produce such a configuration by starting from type IIA and quotienting by $`(1)^{F_L}`$, since type IIA has no BPS D5-branes. Thus the only way to produce this configuration will be by kink condensation starting from a $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair stretched between the D5-branes. This leads to two possibilities, which are linked to the two possibilities considered in the previous section. If the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair retains a complex tachyon, then condensation of a real kink will give rise to the unstable D2-brane suspended between D5-branes. However, as noted in the previous section, this does give rise to a potential paradox. The D2-brane would in turn be tachyonic, and its decay could potentially produce a BPS D-string suspended between D5-branes, which we know to be inconsistent.
A more plausible scenario, in the light of our discussion in the previous section, arises if the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair between D5-branes retains only one real tachyon. In that case, condensation of a tachyonic kink would lead to a D2-brane suspended between D5-branes, which in turn would be stable because there is no second tachyon. An amusing feature of this scenario is that instead of the usual 3-step chain: $`\mathrm{D3}\overline{\mathrm{D}}3\mathrm{D2}\mathrm{D1}`$, where the successive elements have two tachyons, one tachyon and no tachyon, we find a shorter chain: $`suspended\mathrm{D3}\overline{\mathrm{D}}3suspended\mathrm{D2}`$, where the two elements have respectively one tachyon and no tachyon.
In this situation, note that the suspended D2-brane not only has no 2-brane charge (which is the defining property of an unstable D-brane) but also no induced lower-form charges. Such charges can only arise via world-volume Chern-Simons couplings involving a tachyon and various even-form RR potentials, so once the tachyon is projected out, this leaves a totally uncharged object.
4. Suspended Adjacent Brane-Antibrane Pairs
In this section we introduce three parallel NS5-branes in type IIA theory and consider a system in which a D4-brane and a $`\overline{\mathrm{D}}4`$-brane end on a common NS5-brane from opposite sides. The configuration is that of Fig.(4.1).
Fig.(4.1): Adjacent $`\mathrm{D}4\overline{\mathrm{D}}4`$ pairs in type IIA.
This type of configuration is much harder to analyse. Even for the BPS analogue (with a D-brane on either side of an NS5-brane), it is difficult to reliably extract the spectrum, and one has to use the fact that one can continuously interpolate to a different configuration of intersecting branes where perturbation theory is reliable\[14\]. Alternatively one can use duality, as we will do in subsequent sections.
The question we want to address is whether this configuration has an instability, and if so, what is its nature and to what final configuration does the system tend. On physical grounds, we might expect the system to be tachyonic and unstable. If there is an attractive force between the brane-antibrane pair, then they will tend to line up as in Fig.(4.2).
Fig.(4.2): Lining up of adjacent $`\mathrm{D}4\overline{\mathrm{D}}4`$ pairs in type IIA.
On the other hand, if the force is repulsive then they will move apart. From the point of view of one of the pair, the other will go to infinity.
Rather than immediately studying the above problem, we turn first to a related system that is more amenable to calculation: that of an adjacent $`\mathrm{D1}\overline{\mathrm{D}}1`$ pair suspended between D3-branes. This configuration is related to the above by S- and T-dualities. However, the need to perform an S-duality, relating strong and weak coupling regimes, means that the physics of this system is not necessarily the same as that of Fig.(4.1). Nevertheless, in a later section we will return to the configuration of Fig.(4.1) and find the same behaviour.
The D3-branes support an $`𝒩=4`$, $`d=3+1`$-dimensional $`U(3)`$ gauge theory. The D-string is a magnetic monopole of charges $`(1,1,0)`$ under the $`U(1)\times U(1)\times U(1)`$ Cartan subalgebra of $`U(3)`$, while the $`\overline{\mathrm{D}}`$-string is a monopole of charges $`(0,1,1)`$. In terms of the $`SU(3)`$ subgroup of $`U(3)`$ (neglecting the centre-of-mass factor), the configuration is like a monopole of one $`U(1)`$ and an anti-monopole of the other. With this description, one may calculate the force between the two D-strings by treating the monopoles as point particles \[19\]. In fact, the result may be seen quite simply, as both D-string and $`\overline{\mathrm{D}}`$-string appear as positive charges on the middle D3-brane. We therefore expect them to repel, and this is indeed the case.
Let us now be both more quantitative and more general. Consider an arbitrary simple gauge group $`𝒢`$, with a single adjoint Higgs field<sup>4</sup> Of course, the $`U(N)`$ gauge theory on D3-brane has 6 adjoint Higgs fields. We assume that all branes are co-linear. The force between two D-strings (as opposed to D-string and anti-D-string) when this is not the case has been calculated in \[20\]. $`\varphi `$, which acquires a vacuum expectation value (VEV), $`\varphi =𝐡𝐇`$, where $`𝐇`$ is the rank $`r`$ dimensional Cartan basis and $`𝐡`$ is an $`r`$-dimensional vector. We assume the VEV is such that $`𝒢`$ is broken to the maximal torus, $`U(1)^r`$.
The vector $`𝐡`$ lies in the root space and therefore determines a fundamental Weyl chamber from which we define the simple roots, $`\alpha _j`$, $`j=1,\mathrm{},r`$. Recall that these are the roots that satisfy $`𝐡\alpha _j>0`$ and have the property that any other root is a linear combination of the $`\alpha `$’s with either all positive or all negative coefficients.
Magnetic monopoles are configurations with magnetic field given asymptotically by,
$$B_i=𝐠𝐇\frac{r_i}{4\pi r^3}$$
The magnetic charge vector $`𝐠`$ is forced by topological considerations to lie in the co-root lattice, $`𝐠=4\pi _jn_j\alpha _j^{}/e`$, where $`n_j`$ are integers and $`e`$ is the gauge coupling constant. Whether a given topological charge sector is to be considered as a monopole or anti-monopole depends on the sign of $`𝐡𝐠`$, as can be seen in the most general form of the Bogomol’nyi equation,
$$B_i=\mathrm{sign}(𝐡𝐠)D_i\varphi $$
In particular, all charges equal to linear combinations of positive co-roots are monopoles, while those equal to linear combinations of negative co-roots are anti-monopoles. We will be interested in charges that are given by the sum of a positive co-root and negative co-root. Notice that for $`r>1`$, there are non-trivial topological sectors with this property. This contrasts with the case of $`SU(2)`$ gauge group where a monopole-antimonopole pair necessarily lies in the sector with zero topological charge.
The method used to calculate the long-range force between two static monopoles of magnetic charges $`𝐠_1`$ and $`𝐠_2`$ separated by a distance $`r`$, is well known \[19\]. Treating the monopoles as point particles, they interact through two massless fields: the gauge field and the Higgs field. We treat each in turn. Firstly, the magneto-static potential is given by
$$V_{em}=𝐠_1𝐠_2\frac{1}{4\pi r}$$
Secondly, the potential due to the massless scalar field is determined by approximating the configuration of two monopoles by a simple superposition of the individual solutions. Using Eqn. (4.1) and (4.1), the asymptotic form of the Higgs field of the second monopole is given by,
$$\varphi =𝐡𝐇\mathrm{sign}(𝐡𝐠)\frac{1}{4\pi r}𝐠𝐇$$
The potential is calculated by examining the energy of the first monopole which, when isolated, is given by $`M=𝐡𝐠_1`$. However, in the presence of the second monopole, the effective mass of the first becomes,
$$M_1^{eff}=\left(𝐡\mathrm{sign}(𝐡𝐠_2)\frac{1}{4\pi r}𝐠_2\right)𝐠_\mathrm{𝟏}$$
which, for large separation $`r1/𝐡`$, has an expansion as
$$\begin{array}{cc}\hfill M_1^{eff}& =𝐡𝐠_1\mathrm{sign}(𝐡𝐠_1)\mathrm{sign}(𝐡𝐠_2)\frac{𝐠_1𝐠_2}{4\pi r}\hfill \\ & =M_1+V_{Higgs}\hfill \end{array}$$
The full potential between two well separated monopoles of charges $`𝐠_1`$ and $`𝐠_2`$ is therefore
$$V=V_{em}+V_{Higgs}=\frac{𝐠_1𝐠_2}{4\pi r}\left(1\mathrm{sign}(𝐡𝐠_1)\mathrm{sign}(𝐡𝐠_2)\right)$$
Let us now illustrate this formula with a few simple examples. Firstly, consider an $`SU(2)`$ gauge theory with the VEV $`h>0`$. We have two possibilities: $`g_1=g_2=1`$ or $`g_1=g_2=1`$. In the first situation, the potential vanishes. This is simply the well-known cancellation of the force between two BPS monopoles. In the second situation however, the magneto-static and Higgs potentials combine to produce a negative potential, reflecting the attractive force between a monopole-antimonopole pair. This is the relevant situation for Fig.(2.3).
Now applying Eqn.(4.1) to the case at hand, we choose $`𝐡`$ such that the simple roots of $`SU(3)`$ are $`(1,1,0)`$ and $`(0,1,1)`$. The sector we are interested in consists of a monopole of charge $`𝐠_1=(1,1,0)`$, and another of charge $`𝐠_2=(0,1,1)`$. We therefore have $`\mathrm{sign}(𝐡𝐠_1)=+1`$ while $`\mathrm{sign}(𝐡𝐠_2)=1`$ and again we see non-cancellation between the magneto-static and Higgs forces. However, now the potential is positive. The adjacent D-string and anti-D-string therefore repel.
This result is interesting in itself, but we can ask what it tells us about the original system of an adjacent $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair. As we observed already, these calculations do not apply directly to the system that we started out with, because the configuration they describe is related to our original one only after S-duality. However, we can make two relevant observations. One is that if we take a system of parallel D5-branes in type IIB and suspend a D3 brane in one interval and a $`\overline{\mathrm{D}}3`$ in the next, then we only need T-dualities, and no S-dualities, to relate that system to the $`SU(3)`$ monopole-antimonopole pair studied above. In this way we can remain at weak coupling throughout, hence the above computation reliably tells us that an adjacent brane-antibrane pair separated by a D5-brane does repel.
The second observation is that even for an adjacent brane-antibrane pair separated by an NS5-brane, the qualitative reason for repulsion still holds: charges of the same sign are deposited on the middle brane, while on the outer branes the product of deposited charges is zero. Thus, just by analysing charges, one sees that the pair should repel. We will confirm this by dualising to fractional branes in a subsequent section.
An interesting extension of the above model(s) arises as follows. Supposing we start with four parallel NS5-branes in type IIA theory. Then we can suspend a D4-brane between the first two and a $`\overline{\mathrm{D}}4`$-brane between the next two (Fig.(4.3)). In this situation the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair is not adjacent, but separated by an “empty” interval. The arguments of the previous section, based on charges, suggest that at tree level there is neither attraction nor repulsion between the $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair in this model. Hence we appear to have found a stable, non-BPS configuration from brane-antibrane constructions. Indeed, in the related D1-D3 system, there exists a corresponding non-BPS solution of the classical field equations, consisting of a monopole and anti-monopole placed in commuting $`SU(2)`$ subgroups.
Fig.(4.3): Non-adjacent $`\mathrm{D}4\overline{\mathrm{D}}4`$ pairs in type IIA.
Note that in this system, the locations of the first two NS5-branes in $`x^7,x^8,x^9`$ must be equal to the location of the D4-brane in these directions, while the locations of the next two NS5-branes in these directions must likewise be equal to that of the $`\overline{\mathrm{D}}4`$-brane ending on them. However, these two sets of locations need not be equal to each other.
This system can be shown to be T-dual to a $`Z_4`$ ALE singularity. The above D4 and $`\overline{\mathrm{D}}4`$ branes correspond to two of the eight different fractional branes that arise when a full $`\mathrm{D4}\overline{\mathrm{D}}4`$ pair is brought to a $`Z_4`$ singularity. The spectrum of light states for strings stretching between such a pair can in principle be analysed using quiver techniques, though we will not do so in this paper.
Note that if we replace the D4 and $`\overline{\mathrm{D}}4`$ branes with D2 and $`\overline{\mathrm{D}}2`$ branes respectively, we end up with a system that is similar to one that was investigated in Ref.\[21\]. In the latter, a well-separated pair of $`Z_2`$ ALE singularities arranged at opposite ends of a circle is shown to support a stable non-BPS state, consisting of a fractional D2-brane wrapped on the first singularity and a $`\overline{\mathrm{D}}2`$-brane wrapped on the second. As the singularities get closer, this system is unstable to decay into a single D-string of type IIA wrapped on the circle that connects them. As we will see, well-separated ALE singularities can be T-dualised to a set of NS5-branes on a very small circle, which is precisely the system being considered in this paper. The limit of a small circle ensures that the approximation of a common world-volume field theory is valid. While the system described in Fig.(4.3) is not dual to a pair of $`Z_2`$ singularites, for some purposes a $`Z_4`$ singularity with the two end cycles shrunk to zero size should provide similar physics. It would be interesting to explore whether a similar phase diagram to the one studied in Ref.\[22\] arises in the $`Z_4`$ model.
5. Borromean Branes
Consider now a configuration in type IIB string theory containing NS5-branes, D5-branes and D3-branes all together. The branes are aligned as follows:
$$\begin{array}{cc}\hfill NS5:& (x^1,x^2,x^3,x^4,x^5)\hfill \\ \hfill D5:& (x^1,x^2,x^7,x^8,x^9)\hfill \\ \hfill D3:& (x^1,x^2,x^6)\hfill \end{array}$$
This is the class of configurations studied by Hanany and Witten\[14\], and following them we also use the notation $`\stackrel{}{x}=(x^3,x^4,x^5)`$ and $`\stackrel{}{y}=(x^7,x^8,x^9)`$. Scalar fields representing translational modes in these directions will be denoted $`\stackrel{}{X}`$ and $`\stackrel{}{Y}`$.
To start with, consider two parallel NS5-branes at the same value of $`\stackrel{}{y}`$, say $`\stackrel{}{y}=\stackrel{}{0}`$, and separated by a finite distance in $`x^6`$. Suspend a D3-brane between them. This can be located at any value of $`\stackrel{}{x}`$, though it must have $`\stackrel{}{y}=0`$ to end on the NS5-branes. At some intermediate value of $`x^6`$, insert a D5-brane at a fixed $`\stackrel{}{x}`$, say $`\stackrel{}{x}=\stackrel{}{0}`$. The whole configuration is shown in Fig.(5.1).
Fig.(5.1): Hanany-Witten configuration.
Any two of the three branes in this construction break the supersymmetry to $`\frac{1}{4}`$ of the maximal value, so the field theory on the common intersection will have 8 supersymmetries. Adding the third brane does not further reduce the number of supersymmetries. To illustrate this, split the 32 supercharges of IIB string theory into four groups of 8, labelled $`Q_\lambda `$, $`\lambda =1,\mathrm{},4`$. Each brane breaks half of the supersymmetry as indicated in the following table,
$$\begin{array}{ccccccccc}& & Q_1& & Q_2& & Q_3& & Q_4\\ \\ NS5& & \times & & & & \times & & \\ \\ D5& & & & \times & & \times & & \\ \\ D3& & \times & & \times & & & & \end{array}$$
where a $`\times `$ denotes that the supersymmetry is broken by a given brane, while a $``$ means that it is preserved. We see that the 8 supercharges of $`Q_4`$ are preserved by all branes, and we therefore have a $`2+1`$-dimensional field theory with $`𝒩=4`$ supersymmetry describing the low energy dynamics of the D3-brane.
The massless spectrum of this $`2+1`$-dimensional field theory is as follows\[14\]. Firstly, the 3-brane suspended between two parallel NS5-branes has a pure $`𝒩=4`$ supersymmetric $`U(1)`$ gauge theory, containing a gauge field, three scalars associated to the transverse motion of the 3-brane in the $`\stackrel{}{x}`$ directions, and fermions. Inserting the D5-brane in the middle gives rise to a charged hypermultiplet from the open string connecting it to the D3-brane. This hypermultiplet becomes massless when the D3-brane moves to $`\stackrel{}{x}=0`$, where it touches the D5-brane.
Now let us ask what happens if we replace the D3-brane with a $`\overline{\mathrm{D}}3`$-brane, while keeping the NS5 and D5 branes unchanged (Fig.(5.2))<sup>5</sup> This configuration was described in Ref.\[23\].. The supersymmetries preserved by the system may be read off from Table (5.1), simply by exchanging each $`\times `$ and $``$ in the D3-brane row, to get
$$\begin{array}{ccccccccc}& & Q_1& & Q_2& & Q_3& & Q_4\\ \\ NS5& & \times & & & & \times & & \\ \\ D5& & & & \times & & \times & & \\ \\ \overline{\mathrm{D}}3& & & & & & \times & & \times \end{array}$$
The system now breaks all supersymmetries. Notice however that any pair of branes still preserves 8 supercharges. We call this configuration “Borromean” (in analogy with the famous topological configuration of three rings, where any two are unlinked but all three together are linked). We may restore supersymmetry by changing either of the 5-branes to an antibrane. If however we change both D5 and NS5-branes to $`\overline{\mathrm{D}}5`$ and $`\overline{\mathrm{NS}}5`$-branes respectively, so that all branes in Fig.(5.1) have been changed to antibranes, we once again have a situation preserving no supersymmetry.
Fig.(5.2): Borromean branes.
Let us now analyse the spectrum of the Borromean brane system. Since the $`\overline{\mathrm{D}}3`$-brane suspended between two NS5-branes is supersymmetric, it will give rise to an $`𝒩=4`$ $`U(1)`$ vector multiplet. On the other hand, if we consider the $`\overline{\mathrm{D}}3`$-brane along with the D5-brane, the pair is also supersymmetric and the open string joining them gives rise, as previously, to an $`𝒩=4`$ hypermultiplet. Since the light fields all come from considering two of the three branes at a time, we see that there is no tachyon in the one-particle Hilbert space of this theory. However, although pairwise supersymmetric, the coupling of the hypermultiplet to the vector multiplet does not preserve supersymmetry. We will now argue that the correct non-supersymmetric $`d=2+1`$-dimensional low-energy gauge theory is given by the usual $`𝒩=4`$ Lagrangian without the Yukawa terms coupling a hypermultiplet scalar and fermion to a vector multiplet fermion.
To see this we use the pairwise supersymmetry of the configuration to determine which fermions on the worldvolume of the D3-brane couple to the 3-5 strings, and which are projected out by the boundary conditions of the NS5-branes. Let us begin with the latter problem and consider an infinite 3-brane. The world-volume fields on this brane are all goldstone modes for broken symmetries: 6 scalars, $`\stackrel{}{X}`$ and $`\stackrel{}{Y}`$, and 16 goldstinoes (corresponding to the supercharges $`Q_1`$ and $`Q_2`$ for the D3-brane of Eqn.(5.1), and to $`Q_3`$ and $`Q_4`$ for the $`\overline{\mathrm{D}}3`$-brane of Eqn.(5.1)). The presence of the NS5-branes on which the 3-brane ends projects out half of these fields. Of the bosons, the surviving fields are $`\stackrel{}{X}`$, which are the space-time directions which do not arise as goldstone bosons on the NS5-brane, together with the $`A_0,A_1`$ and $`A_2`$ component of the gauge field. By supersymmetry, the surviving fermions are those goldstinoes which arise from supersymmetries broken by the 3-brane, but preserved by the NS5-branes. In the case of the D3-brane, this means $`Q_2`$, while for the $`\overline{\mathrm{D}}3`$-brane, it is $`Q_4`$.
Now let us examine the hypermultiplet fields arising from the 3-5 string. We are interested in determining to which fields on the 3-brane they couple. Once again, we consider first the case of an infinite 3-brane. The world-volume theory is $`d=3+1`$, $`𝒩=4`$ $`U(1)`$ gauge theory. The presence of the D5-brane breaks this supersymmetry by half and the $`𝒩=4`$ multiplet splits into an $`𝒩=2`$ vector multiplet and a neutral $`𝒩=2`$ hypermultiplet. The manner in which this split occurs is determined by the fact that the charged hypermultiplet arising from the 3-5 string is minimally coupled only to the former. Notice that as the D5-brane is not Lorentz invariant from the point of view of an infinite 3-brane world-volume, the gauge field will also be split, partly living in the vector multiplet and partly in the hypermultiplet. More specifically, the 3-5 string couples to the components of the gauge field $`A_0,A_1`$ and $`A_2`$, together with the scalars $`\stackrel{}{X}`$. To see the latter is particularly simple, as motion of the 3-brane in the $`\stackrel{}{x}`$ directions stretches the 3-5 string and so gives a mass to the charged hypermultiplet. Stated differently, the charged 3-5 hypermultiplet couples to those scalars which are goldstone modes for both the 3-brane and the D5-brane. By supersymmetry, it therefore also couples to the fermions on the 3-brane which are goldstinoes for both 3-brane and D5-brane. Hence, in the case of the D3-brane in (5.1), the 3-5 string couples to $`Q_2`$, while for the $`\overline{\mathrm{D}}3`$-brane of (5.1), it couples to $`Q_3`$.
Putting these two facts together allows us to determine the field theory for the Borromean branes of Fig.(5.2). The 3-5 string couples to the scalars $`\stackrel{}{X}`$, the gauge fields $`A_0,A_1,A_2`$ and the goldstinoes $`Q_3`$. However, the NS5-branes project out all fields except for the above bosons and the goldstinoes $`Q_4`$. In particular, the fermions to which the hypermultiplet couples are lost due to the boundary conditions imposed by the NS5-brane. Therefore, the gauge theory on the D3-brane is as in the supersymmetric case, except that the Yukawa terms coupling the 3-5 string hypermultiplet to the 3-3 string fermions is missing, thus breaking supersymmetry.
Having identified the non-supersymmetric gauge theory living on the D3-brane, we must now determine the correct vacuum state of this theory. Classically, the theory has flat directions given by the vacuum expectation values of the scalars $`\stackrel{}{X}`$ and the dual photon. As the theory is non-supersymmetric, the flat directions associated to the scalars $`\stackrel{}{X}`$ are not protected against quantum corrections. However, they do remain at one-loop as a consequence of bose-fermi degeneracy. Whether this cancellation continues at higher loops is an open question.
One may pass from the Borromean brane situation to that of adjacent brane-antibrane pairs, generalising those discussed in the previous section, by passing the D5-brane through one of the NS5-branes. The pairwise supersymmetry of these two 5-branes ensures that the Hanany-Witten transition\[14\] occurs as in the supersymmetric case, and we are left with a D3-brane suspended between the D5 and NS5-brane (Fig.(5.3)).
Fig.(5.3): Result of Hanany-Witten transition on Borromean branes.
The brane configuration is no longer Borromean, as the D3 and $`\overline{\mathrm{D}}3`$-branes are not pairwise supersymmetric. Applying the considerations of the previous section locally to the region near the middle 5-brane, we would expect that the $`\mathrm{D3}\overline{\mathrm{D}}3`$ pair repels. However, it is quite possible that in this model the Hanany-Witten transition is accompanied by a phase transition, so we are unable to relate the physics of Figs.(5.2) and (5.3). Note in particular that the dimensionless expansion parameter in the theory on the Borromean branes is $`e^2/r`$, where $`e`$ is the gauge coupling constant, and $`r`$ the distance between the D5-brane and $`\overline{\mathrm{D}}3`$-brane. The one-loop flatness of the vacuum moduli space of this theory ensures that any potential that is generated at large distance is of the form $`\pm 1/r^n`$, for $`n2`$. This is to be contrasted with the $`1/r`$ behaviour of adjacent brane-antibrane pairs separated by a D5-brane.
6. Dualising NS5-branes into ALE spaces
Consider type IIA string theory compactified on a 4-torus $`T^4`$ in the directions $`x^6,x^7,x^8,x^9`$ and orbifolded by the inversion $`I_{6789}`$. This gives rise to 16 fixed orbifold 5-planes. Each of these is locally a $`Z_2`$ ALE space. Or goal is to dualise this system and then, by going to a suitable region of moduli space, to isolate the dual of a single $`Z_2`$ ALE space. The material in this section is not new, having been discussed in slightly different language in Refs.\[6\] and in the context of conifold singularities in Refs.\[24,,25\]. The duality that we discuss below will be used in the subsequent section to map the suspended brane systems to brane-antibrane generalisations of quiver theories.
Under T-duality along $`x^6`$, the inversion $`I_{6789}`$ is mapped to $`(1)^{F_L}I_{6789}`$, as one can easily see by examining the action on various massless states. At the same time, the type IIA theory is mapped to the type IIB theory. Thus the dual to the original orbifold is the orbifold of type IIB string theory on a 4-torus quotiented by $`(1)^{F_L}I_{6789}`$. This too has 16 fixed points, but the nature of the orbifold 5-planes is quite different. To see this, perform an S-duality. The operator $`(1)^{F_L}`$ is mapped to the orientifolding operation $`\mathrm{\Omega }`$, and the resulting theory has 16 orientifold 5-planes (O5-planes) along with 16 (mirror pairs of) D5-branes. The D5-branes can sit on the orientifold planes to make charge-cancelled configurations, or they can move off. Hence it must be that, before performing S-duality, the 16 orbifold planes were made up of NS5-branes (which are S-dual to D5-branes in type IIB) and static planes that are S-dual to O5-planes. We call these “S5-planes”.
To summarise, the T-dual of type IIA on a $`T^4/Z_2`$ orbifold is type IIB on a 4-torus with 16 S5-planes and 16 pairs of NS5-branes. At the orbifold point of the type IIA theory, the flux of the 2-form $`B`$-field through each $`Z_2`$ singularity is $`\frac{1}{2}`$, while on the IIB side, each NS5-brane is on top of an S5-plane. Now since the T-duality was performed only on $`x^6`$, we can identify the $`x^7,x^8,x^9`$ positions on the two sides and decompactify them. As a result we have on the type IIA side, a pair of $`Z_2`$ ALE spaces located symmetrically around the $`x^6`$ circle, and at the origin of $`x^7,x^8,x^9`$, while on the type IIB side there is a pair of S5-planes similarly located symmetrically around the dual $`x^6`$ circle, and a pair of NS5-branes at points of this circle (Fig.(6.1)).
Fig.(6.1): Duality between $`Z_2`$ ALE spaces and S5-NS5 systems.
The matching of moduli is reasonably straightforward. In particular, if $`R`$ is the radius of the $`x^6`$ circle on the type IIA side then the dual radius is $`\frac{1}{R}`$. Let $`B,B^{}`$ be the fluxes of the $`B`$-field through the two $`Z_2`$ ALE spaces on the IIA side, and let $`X_6^{(1)},X_6^{(2)}`$ be the $`x_6`$ locations of the two NS5-branes on the IIB side. Then we have:
$$\begin{array}{cc}\hfill X_6^{(1)}& =\frac{1}{R}|BB^{}|\hfill \\ \hfill X_6^{(2)}& =\frac{1}{R}|B+B^{}|\hfill \end{array}$$
Note that with this matching, the type IIA theory has an enhanced gauge symmetry if $`B=0`$ or $`B^{}=0`$, coming from a BPS D2-brane wrapping the 2-cycle of the first or second ALE space, and this corresponds in type IIB to $`X_6^{(1)}=X_6^{(2)}`$ so the two NS5-branes meet. If both $`B=B^{}=0`$ then we have $`SU(2)\times SU(2)`$ enhanced gauge symmetry, and on the type IIB side, $`X_6^{(1)}=X_6^{(2)}=0`$, so the two NS5-branes meet on an S5-plane, where in fact we expect $`SO(4)SU(2)\times SU(2)`$.
Suppose we now fix $`B^{}=\frac{1}{2}`$ once and for all. Then no enhanced gauge symmetry will ever arise from the second ALE space, and henceforth we focus our attention on the region near the first ALE space. The dual picture now has the two NS5-branes at equal distances from their corresponding S5-planes. They can only meet far away from the S5-planes (in fact, at the midpoint of the strip depicted in Fig.(6.1)). Hence the S5 orbifold planes can be ignored. The result is a duality between a pair of NS5-branes arranged around a circle (in type IIB) and a solitary $`Z_2`$ ALE singularity (in type IIA). The roles of type IIA and type IIB can also be reversed.
Now we may wrap a D$`p`$-brane around the $`x^6`$ circle and perform the above duality. It turns into a D$`(p1)`$ brane at the $`Z_2`$ ALE singularity. Breaking the original $`p`$-brane on the NS5-branes corresponds to breaking the $`(p1)`$-brane into a pair of fractional branes, which are free to move around while remaining in the 5-plane of the ALE singularity.
In a similar fashion, one can start with a 2-torus slanted at a particular angle and carry out a $`Z_k`$ quotient $`z\omega z`$ where $`\omega ^k=1`$. Near the origin, the system can be described in T-dual language by a set of $`k`$ NS5-branes located around the $`x^6`$ circle. The separations of the NS5-branes along $`x^6`$ are given by B-fields through the various 2-cycles of the $`Z_k`$ singularity, while their locations in $`x^7,x^8,x^9`$ are determined by the geometrical size parameters of these 2-cycles.
7. Brane-Antibrane Pairs at an ALE Singularity
In this section we obtain the spectrum of light fields (including possible tachyons) on brane-antibrane pairs placed at ALE quotient singularites. For simplicity we will work mainly with the case of $`R^4/Z_2`$. We start with this singularity, where $`R^4`$ corresponds to the directions $`x^6,x^7,x^8,x^9`$, and bring a $`\mathrm{D}p\overline{\mathrm{D}}p`$ pair close to it ($`p`$ is even in type IIA and odd in type IIB). Our analysis will closely parallel that of Ref.\[7\] where BPS D-branes at ALE spaces were studied. Some aspects of brane-antibrane systems at ALE spaces have been investigated previously in Ref.\[26\].
In order to work out the spectrum of the resulting theory, we first introduce the “mirror” of the $`\mathrm{D}p\overline{\mathrm{D}}p`$ pair in the $`Z_2`$ singularity, thus making four branes altogether. Let these four branes be labelled $`1,1^{},\overline{1},\overline{1}^{}`$, where the prime denotes a mirror brane and the bar denotes an antibrane. Thus there are $`4\times 4`$ Chan-Paton factors, which we classify as follows. We use the $`2\times 2`$ Chan-Paton labels $`A,B,C,D`$ to distinguish between strings on a brane, on a mirror brane, or between the two, while the labels $`a,b,c,d`$ are used (as before) to distinguish between strings on a brane, on an antibrane, or between the two. The correspondence between $`a,b,c,d`$ and Chan-Paton matrices was given in Eqn.(2.1), and the analogous correspondence holds also for $`A,B,C,D`$.
With these definitions, we have the correspondence:
$$\begin{array}{c}11:Aa\\ 11^{}:Ca\\ 1^{}1:Da\\ 1^{}1^{}:Ba\end{array}\begin{array}{c}\overline{1}\overline{1}:Ab\\ \overline{1}\overline{1}^{}:Cb\\ \overline{1}^{}\overline{1}:Db\\ \overline{1}^{}\overline{1}^{}:Bb\end{array}\begin{array}{c}1\overline{1}:Ac\\ 1\overline{1}^{}:Cc\\ 1^{}\overline{1}:Dc\\ 1^{}\overline{1}^{}:Bc\end{array}\begin{array}{c}\overline{1}1:Ad\\ \overline{1}1^{}:Cd\\ \overline{1}^{}1:Dd\\ \overline{1}^{}1^{}:Bd\end{array}$$
The spectrum of this system, before projection by the $`Z_2`$, contains gauge fields $`A_\mu ^{(I)}`$, $`\mu =0,\mathrm{},p`$, along with scalars $`X_i^{(I)}`$, $`i=p+1,\mathrm{},5`$, and $`X_m^{(I)}`$, $`m=6,7,8,9`$. These fields all lie in the adjoint of $`U(2)\times U(2)`$, with the superscript $`(I)`$ labelling the relevant factor. Along with the obvious fermionic counterparts, these make up a set of massless fields arising in the Chan-Paton sectors that appear in the first two columns of Eqn.(7.1). Each Chan-Paton sector in the remaining two columns gives rise to a tachyon $`T`$, hence there are 8 real tachyons altogether. These are accompanied by a set of massless Ramond fermions which are oppositely GSO-projected as compared to the ones that accompanied the gauge field and massless scalars.
There are two interesting $`Z_2`$ involutions: the inversion that creates the ALE space, which we denote by $`I_{6789}`$, and the symmetry $`(1)^{F_L}`$ which we have used before. In our conventions, these involutions act as conjugation by:
$$\begin{array}{cc}\hfill I_{6789}& :\sigma _11\hfill \\ \hfill (1)^{F_L}& :1\sigma _1\hfill \end{array}$$
This follows from the fact that the first $`\sigma _1`$ factor exchanges a brane with its mirror, while the second $`\sigma _1`$ exhanges a brane with an antibrane.
It is now straightforward to write down the transformations of the various fields under $`I_{6789}`$:
$$\begin{array}{ccccc}\underset{¯}{\mathrm{Field}}& & \underset{¯}{\mathrm{CP}\mathrm{factor}}& & \underset{¯}{I_{6789}}\\ & & & \\ A_\mu ^{(I)},X_i^{(I)}& & \{1,\sigma _1\}\{1,\sigma _3\}& & +\\ X_m^{(I)}& & \{i\sigma _2,\sigma _3\}\{1,\sigma _3\}& & +\\ A_\mu ^{(I)},X_i^{(I)}& & \{i\sigma _2,\sigma _3\}\{1,\sigma _3\}& & \\ X_m^{(I)}& & \{1,\sigma _1\}\{1,\sigma _3\}& & \\ T& & \{1,\sigma _1\}\{\sigma _1,i\sigma _2\}& & +\\ T& & \{i\sigma _2,\sigma _3\}\{\sigma _1,i\sigma _2\}& & \end{array}$$
Thus under orbifolding by $`I_{6789}`$, the fields which are invariant survive, while the others are projected out. The gauge group that survives (for a brane-antibrane pair on top of the orbifold singularity) is $`U(1)^4`$. This statement really applies when all vev’s are zero, which is the common origin of the Higgs branch and Coulomb branch. We see from the table that 4 real tachyons survive in the orbifolded theory, these can be grouped into two complex tachyons that transform as the bi-fundamental of the first $`U(1)\times U(1)`$ and the second $`U(1)\times U(1)`$ respectively.
On the Higgs branch, the brane-antibrane pair leaves the orbifold plane and the gauge group is Higgsed to $`U(1)\times U(1)`$. This corresponds to retaining the following fields from the above table:
$$\begin{array}{cc}\hfill A_\mu ^{(I)},X_i^{(I)}& :1\{1,\sigma _3\}\hfill \\ \hfill X_m^{(I)}& :\sigma _3\{1,\sigma _3\}\hfill \\ \hfill T& :1\{\sigma _1,i\sigma _2\}\hfill \end{array}$$
(For $`X_m`$ this involves a gauge choice, see Ref.\[27\].) Thus on this branch there is a single complex tachyon, transforming as the bi-fundamental of $`U(1)\times U(1)`$.
We can relate all this to the brane constructions that we discussed above. Compactify $`x^6`$ on a circle and T-dualize, then we have a pair of NS5-branes on the circle and a $`\mathrm{D}(p+1)\overline{\mathrm{D}}(p+1)`$ pair running all the way around the circle (Fig.(6.1)).
Fig.(7.1): A pair of NS5-branes and a $`\mathrm{D}(p+1)\overline{\mathrm{D}}(p+1)`$ pair wrapped on $`x^6`$.
The figure describes the point where the Higgs branch and Coulomb branch meet, and it is clear that the gauge group should be $`U(1)^4`$ as we have shown. Moreover, in this brane construction one expects 2 complex tachyons, corresponding to open strings connecting a brane and antibrane segment suspended between the same NS5-branes. They manifestly carry the right charges. Moving onto the Higgs branch is accomplished by taking the brane-antibrane pair away from the NS5-branes, this breaks the gauge group to $`U(1)\times U(1)`$ and there is just one complex tachyon.
To understand the Coulomb branch, observe that both the brane and the antibrane which wrap all the way around $`x^6`$ in the dual brane construction can break on the NS5-branes. (In the T-dual picture this corresponds to the fact that branes at an ALE singularity can break into fractional branes\[27,,7,,28,,29\]. Hence, altogether we have four brane segments in the problem (or rather, two brane segments and two antibrane segments). We can separate all these and take some of the segments off to infinity. For example, suppose we take a brane and an antibrane stretching along the same segment to infinity. The result (after taking the radius very large) will be the diagram of Fig.(2.1). On the other hand, we can take away an antibrane from one segment and a brane from the adjacent segment. The result will be an adjacent brane-antibrane pair, similar to the construction in Fig.(4.1). We will use these facts in the T-dual picture to understand more about the brane constructions of Figs.(2.1),(4.1).
Returning to the $`\mathrm{D}p\overline{\mathrm{D}}p`$ pair at an ALE singularity, we should now identify the Chan-Paton factors corresponding to strings connecting fractional branes. Recall that for a single D$`p`$-brane at a $`Z_2`$ singularity, the Coulomb branch arises when $`X_m=0`$ and $`X_i\{1,\sigma _1\}`$. It is on this branch that the single D$`p`$-brane is said to split into a pair of fractional branes and the gauge group becomes $`U(1)\times U(1)`$. Note that the individual gauge fields associated to the two fractional branes are associated to the Chan-Paton factors $`1+\sigma _1`$ and $`1\sigma _1`$ respectively.
For a $`\mathrm{D}p\overline{\mathrm{D}}p`$ system, the Higgs branch is described by $`X_m^{(I)}0`$, $`X_i^{(I)}1\{1,\sigma _3\}`$ while the Coulomb branch is $`X_m^{(I)}=0`$, $`X_i^{(I)}\{1,\sigma _1\}\{1,\sigma _3\}`$. This time, as expected, the latter branch describes four fractional branes altogether. Let us denote these fractional branes as $`1_f,1_f^{},\overline{1}_f`$ and $`\overline{1}_f^{}`$ (note that in this case, the prime does not denote the mirror image, though the bar still denotes the antibrane). The interpretations of these branes are as follows: $`1_f`$ is a (p+2)-brane wrapped on the vanishing 2-cycle $`\mathrm{\Sigma }`$ at the ALE singularity<sup>6</sup> At the orbifold point, we have an unresolved orbifold singularity through which the $`B`$-field has a flux of $`\frac{1}{2}`$. This also means that the NS5-branes are symmetrically placed along the $`x^6`$ circle in the T-dual construction. However, we can of course allow $`B`$ to vary, and in the following discussion we have in mind this more general case.. $`1_f^{}`$ is an anti-(p+2)-brane wrapped on the same 2-cycle. On the other hand, $`\overline{1}_f`$ is an anti-(p+2)-brane, while $`\overline{1}_f^{}`$ is a (p+2)-brane.
Following the discussion in Ref.\[29\], we can assign the following world-volume couplings to the four fractional branes:
$$\begin{array}{cc}& 1_f:(C_{RR}^{(p+3)}+(B_{NS}F_{1_f}))C_{RR}^{(p+1)}\hfill \\ & 1_f^{}:(C_{RR}^{(p+3)}(B_{NS}F_{1_f^{}}))C_{RR}^{(p+1)}\hfill \\ & \overline{1}_f:(C_{RR}^{(p+3)}(B_{NS}F_{\overline{1}_f}))C_{RR}^{(p+1)}\hfill \\ & \overline{1}_f^{}:(C_{RR}^{(p+3)}+(B_{NS}F_{\overline{1}_f^{}}))C_{RR}^{(p+1)}\hfill \end{array}$$
Moreover, the world-volume field strengths are given by $`_\mathrm{\Sigma }F_{1_f}=_\mathrm{\Sigma }F_{\overline{1}_f}=0`$ and $`_\mathrm{\Sigma }F_{1_f^{}}=_\mathrm{\Sigma }F_{\overline{1}_f^{}}=1`$.
We would like to find the spectrum arising from all possible configurations of open strings on this set of four fractional branes. This can be read off directly from (7.1). First consider open strings starting and ending on the same brane. These each carry a gauge field $`A_\mu `$ and a pair of scalars $`X_i`$ $`(i=4,5)`$ having Chan-Paton factors:
$$\begin{array}{cc}\hfill 1_f1_f& :\frac{1}{2}(1+\sigma _1)\frac{1}{2}(1+\sigma _3)\hfill \\ \hfill 1_f^{}1_f^{}& :\frac{1}{2}(1\sigma _1)\frac{1}{2}(1+\sigma _3)\hfill \\ \hfill \overline{1}_f\overline{1}_f& :\frac{1}{2}(1+\sigma _1)\frac{1}{2}(1\sigma _3)\hfill \\ \hfill \overline{1}_f^{}\overline{1}_f^{}& :\frac{1}{2}(1\sigma _1)\frac{1}{2}(1\sigma _3)\hfill \end{array}$$
Next, we have in principle 6 oriented strings going from each one of the four fractional branes to another, and each string has two orientations, for a total of 12. However, some of these are projected out. In fact, we get massless hypermultiplet scalars $`X_m`$ as follows:
$$\begin{array}{cc}\hfill 1_f1_f^{}& :\frac{1}{2}(\sigma _3+i\sigma _2)\frac{1}{2}(1+\sigma _3)\hfill \\ \hfill 1_f^{}1_f& :\frac{1}{2}(\sigma _3i\sigma _2)\frac{1}{2}(1+\sigma _3)\hfill \\ \hfill \overline{1}_f\overline{1}_f^{}& :\frac{1}{2}(\sigma _3+i\sigma _2)\frac{1}{2}(1\sigma _3)\hfill \\ \hfill \overline{1}_f^{}\overline{1}_f& :\frac{1}{2}(\sigma _3i\sigma _2)\frac{1}{2}(1\sigma _3)\hfill \end{array}$$
Tachyons arise from the following open strings:
$$\begin{array}{cc}\hfill 1_f\overline{1}_f& :\frac{1}{2}(1+\sigma _1)\frac{1}{2}(\sigma _1+i\sigma _2)\hfill \\ \hfill \overline{1}_f1_f& :\frac{1}{2}(1+\sigma _1)\frac{1}{2}(\sigma _1i\sigma _2)\hfill \\ \hfill 1_f^{}\overline{1}_f^{}& :\frac{1}{2}(1\sigma _1)\frac{1}{2}(\sigma _1+i\sigma _2)\hfill \\ \hfill \overline{1}_f^{}1_f^{}& :\frac{1}{2}(1\sigma _1)\frac{1}{2}(\sigma _1i\sigma _2)\hfill \end{array}$$
Thus there are four real tachyons in the Coulomb branch. The remaining open strings correspond in principle to the following Chan-Paton factors:
$$\begin{array}{cc}\hfill 1_f\overline{1}_f^{}& :\frac{1}{2}(\sigma _3+i\sigma _2)(\sigma _1+i\sigma _2)\hfill \\ \hfill \overline{1}_f^{}1_f& :\frac{1}{2}(\sigma _3i\sigma _2)(\sigma _1i\sigma _2)\hfill \\ \hfill 1_f^{}\overline{1}_f& :\frac{1}{2}(\sigma _3i\sigma _2)(\sigma _1+i\sigma _2)\hfill \\ \hfill \overline{1}_f1_f^{}& :\frac{1}{2}(\sigma _3+i\sigma _2)(\sigma _1i\sigma _2)\hfill \end{array}$$
However, as one can see from (7.1), these strings are projected out: the tachyon is removed by $`I_{6789}`$ and the massless states are removed by the anti-GSO projection. Hence we find the novel result that strings connecting $`1_f`$ to $`\overline{1}_f^{}`$ have no tachyonic or massless bosonic states.
Now let us consider first the system of a suspended $`\mathrm{D}(p+1)\overline{\mathrm{D}}(p+1)`$ pair as in Fig.(2.1). This is described in the T-dual version by keeping the fractional branes $`1_f`$ and $`\overline{1}_f`$. We see that for all values of $`_\mathrm{\Sigma }B_{NS}`$, the total (p+2)-brane and p-brane charges add up to zero. The analysis above tells us that the spectrum on this pair consists of a $`U(1)\times U(1)`$ gauge field, massless scalars $`X_i`$ and a complex tachyon $`T`$. All this is consistent with the fact that we expect the original pair to be able to annihilate completely. Thus we confirm the heuristic picture of this system developed in Section 2.
Turning next to the system of an adjacent $`\mathrm{D}(p+1)\overline{\mathrm{D}}(p+1)`$ pair as in Fig.(4.1), we see that in this case the dual description is obtained by keeping the fractional branes $`1_f`$ and $`\overline{1}_f^{}`$ (or equivalently $`\overline{1}_f`$ and $`1_f^{}`$, with some sign changes in the following). In this case the picture is very different. The system has a net (p+2)-brane charge of $`+2`$, and a net p-brane charge of $`2B_{NS}1`$. The latter vanishes only at the symmetric point $`B_{NS}=\frac{1}{2}`$, for which the NS5-branes in the original construction are equally spaced around the circle. As far as the spectrum is concerned, there are neither any tachyons nor any massless scalars coming from open strings between the pair.
Finally, one can consider non-adjacent brane-antibrane pairs from this point of view by taking a brane-antibrane pair to a $`Z_4`$ ALE singularity, where it can split into a total of 8 fractional branes. We will not describe this explicitly here.
8. Boundary State Computation of Brane-Antibrane Forces
In this section our aim is to construct the boundary states corresponding to the fractional D3-branes and use them to compute forces between pairs of fractional branes of the various types discussed in the previous section. We will use the conventions of Ref.\[30\].
For this we have to construct consistent boundary states which describe D3-branes in the $`Z_2`$ orbifold of type IIB string theory where the orbifolding group $`Z_2`$ is generated by $`I_{6789}`$ (which we refer to as $`R`$). Boundary states at orbifolds have been classified and constructed in great generality in Ref.\[31\]. Here we will only describe the formulae relevant for our purpose.
Let us consider a D3-brane in type IIB string theory on $`R^{9,1}`$. Take the D3-brane to be along the $`x^3,x^4,x^5`$ directions, and situated at the origin in $`x^6,x^7,x^8,x^9`$ directions. We make a double Wick rotation $`x^0ix^0`$, $`x^1ix^1`$, so that the light-cone directions are $`x^1\pm x^2`$ and the D3-brane world-volume directions are all space-like. We first define the basic boundary states in terms of which the D3-brane state will be defined.
In the untwisted sector these states are:
$$\begin{array}{cc}\hfill |k,\eta _{\genfrac{}{}{0pt}{}{NSNS;U}{RR;U}}& =\mathrm{exp}(\underset{n=1}{\overset{\mathrm{}}{}}[\frac{1}{n}\underset{\mu =0,3,4,5}{}\alpha _n^\mu \stackrel{~}{\alpha }_n^\mu +\frac{1}{n}\underset{\mu =6}{\overset{9}{}}\alpha _n^\mu \stackrel{~}{\alpha }_n^\mu ]\hfill \\ & +i\eta \underset{r>0}{}[\underset{\mu =0,3,4,5}{}\psi _r^\mu \stackrel{~}{\psi }_r^\mu +\underset{\mu =6}{\overset{9}{}}\psi _r^\mu \stackrel{~}{\psi }_r^\mu ])|k,\eta _{\genfrac{}{}{0pt}{}{NSNS;U}{RR;U}}^{(0)}\hfill \end{array}$$
where $`k(k^1,k^2,k^6,k^7,k^8,k^9)`$ and $`\eta =\pm `$. In the NSNS sector, the labels $`n`$, $`r`$ are
$$\begin{array}{cc}\hfill n& Z_+\text{for}\mu =0,3,4,\mathrm{},9\hfill \\ \hfill r& Z_+\frac{1}{2}\text{for}\mu =0,3,4,\mathrm{},9\hfill \end{array}$$
while in the RR sector they are:
$$n,rZ_+\text{for}\mu =0,3,4,\mathrm{},9$$
The NSNS vacuum state is independent of $`\eta `$, while the RR vacuum states are defined as:
$$\begin{array}{cc}\hfill \psi _{}^\mu |k,\eta _{RR;U}^{(0)}& =0\text{for}\mu =0,3,4,5.\hfill \\ \hfill \psi _+^\mu |k,\eta _{RR;U}^{(0)}& =0\text{for}\mu =6,7,8,9.\hfill \\ \hfill |k,+_{RR;U}^{(0)}& =\underset{\mu =6}{\overset{9}{}}\psi _{}^\mu \underset{\mu =0,3}{\overset{5}{}}\psi _+^\mu |k,_{RR;U}^{(0)}\hfill \end{array}$$
Next let us turn to the twisted sector of the $`Z_2`$ orbifold. Here the relevant boundary states are given by:
$$\begin{array}{cc}\hfill |k,\eta _{\genfrac{}{}{0pt}{}{NSNS;T}{RR;,T}}& =\mathrm{exp}(\underset{n>0}{\overset{\mathrm{}}{}}[\frac{1}{n}\underset{\mu =0,3,4,5}{}\alpha _n^\mu \stackrel{~}{\alpha }_n^\mu +\frac{1}{n}\underset{\mu =6}{\overset{9}{}}\alpha _n^\mu \stackrel{~}{\alpha }_n^\mu ]\hfill \\ & +i\eta \underset{r>0}{}[\underset{\mu =0,3,4,5}{}\psi _r^\mu \stackrel{~}{\psi }_r^\mu +\underset{\mu =6}{\overset{9}{}}\psi _r^\mu \stackrel{~}{\psi }_r^\mu ])|k,\eta _{\genfrac{}{}{0pt}{}{NSNS;T}{RR;T}}^{(0)}\hfill \end{array}$$
where $`k(k^1,k^2)`$. In the twisted NSNS sector, the labels $`n`$, $`r`$ are
$$\begin{array}{cc}\hfill n& Z_+\text{for}\mu =0,3,4,5,\hfill \\ & Z_+\frac{1}{2}\text{for}\mu =6,7,8,9,\hfill \\ \hfill r& Z_+\frac{1}{2}\text{for}\mu =0,3,4,5,\hfill \\ & Z_+\text{for}\mu =6,7,8,9.\hfill \end{array}$$
while in the twisted RR sector they are:
$$\begin{array}{cc}\hfill n,r& Z_+\text{for}\mu =0,3,4,5,\hfill \\ & Z_+\frac{1}{2}\text{for}\mu =6,7,8,9.\hfill \end{array}$$
Since in the twisted sector both NSNS and RR sectors have zero modes we have to define the vacua carefully (as we did for the RR sector of the untwisted sector). Define
$$\begin{array}{cc}\hfill \psi _+^\mu |k,_{NSNS;T}^{(0)}& =0\text{for}\mu =6,7,8,9.\hfill \\ \hfill |k,+_{NSNS;T}^{(0)}& =\underset{\mu =6}{\overset{9}{}}\psi _{}^\mu |k,_{NSNS;T}^{(0)}\hfill \end{array}$$
Similarly for the RR vacuum state, define
$$\begin{array}{cc}\hfill \psi _{}^\mu |k,_{RR;T}^{(0)}& =0\text{for}\mu =0,3,4,5,.\hfill \\ \hfill |k,+_{RR;T}^{(0)}& =\underset{\mu =0,3}{\overset{5}{}}\psi _+^\mu |k,_{RR;T}^{(0)}\hfill \end{array}$$
Now we integrate these states over momentum, to define states corresponding to branes at fixed positions. In the untwisted sector we get:
$$\begin{array}{cc}\hfill |\eta _{NSNS;U}& =𝒩\underset{\mu =1,2,6}{\overset{9}{}}dk^\mu |k,\eta _{NSNS;U}\hfill \\ \hfill |\eta _{RR;U}& =4i𝒩\underset{\mu =1,2,6}{\overset{9}{}}dk^\mu |k,\eta _{RR;U}\hfill \end{array}$$
while the corresponding states in the twisted sector are:
$$\begin{array}{cc}\hfill |\eta _{NSNS;T}& =2\stackrel{~}{𝒩}\underset{\mu =1}{\overset{2}{}}dk^\mu |k,\eta _{NSNS;T}\hfill \\ \hfill |\eta _{RR}& =2i\stackrel{~}{𝒩}\underset{\mu =1}{\overset{2}{}}dk^\mu |k,\eta _{RR;T}\hfill \end{array}$$
Next we combine the above states into the appropriate GSO-invariant linear combinations to describe D3 and $`\overline{\mathrm{D}}3`$ branes at the ALE space. In the untwisted sector this gives the states:
$$\begin{array}{cc}\hfill |U_{NSNS}& =\frac{1}{\sqrt{2}}\left(|+_{NSNS;U}|_{NSNS;U}\right),\hfill \\ \hfill |U_{RR}& =\frac{1}{\sqrt{2}}\left(|+_{RR;U}+|_{RR;U}\right),\hfill \end{array}$$
while in the twisted sector the GSO-invariant combinations are:
$$\begin{array}{cc}\hfill |T_{NSNS}& =\frac{1}{\sqrt{2}}\left(|+_{NSNS;T}+|_{NSNS;T}\right)\hfill \\ \hfill |T_{RR}& =\frac{1}{\sqrt{2}}\left(|+_{RR;T}+|_{RR;T}\right)\hfill \end{array}$$
Finally we can combine the untwisted and twisted sector boundary states to produce states that describe branes in the full theory. We find four independent consistent boundary states for D3, $`\overline{\mathrm{D}}3`$, which can be identified with the four fractional branes $`1_f,1_f^{},\overline{1}_f,\overline{1}_f^{}`$. The states along with their identifications are as follows:
$$\begin{array}{cc}\hfill |\mathrm{D3},+& =\frac{1}{2}(|U_{NSNS}+|U_{RR}+|T_{NSNS}+|T_{RR}):1_f\hfill \\ \hfill |\mathrm{D3},& =\frac{1}{2}(|U_{NSNS}+|U_{RR}|T_{NSNS}|T_{RR}):1^{}_f\hfill \\ \hfill |\overline{\mathrm{D}}3,+& =\frac{1}{2}(|U_{NSNS}|U_{RR}|T_{NSNS}+|T_{RR}):\overline{1}^{}_f\hfill \\ \hfill |\overline{\mathrm{D}}3,& =\frac{1}{2}(|U_{NSNS}|U_{RR}+|T_{NSNS}|T_{RR}):\overline{1}_f\hfill \end{array}$$
Now we can explicitly write down the closed string tree amplitudes between the above boundary states in terms of open string one loop amplitudes. First, we choose $`32𝒩^2=\frac{v^{(4)}}{(2\pi )^4}`$, where $`v^{(4)}`$ is the infinite volume of the brane along the $`x^6,x^7,x^8,x^9`$ directions. This normalization is determined by the requirement that amplitudes between boundary states can be interpreted as open-string traces. Next, define $`\stackrel{~}{q}=\mathrm{exp}(\pi t)`$ and
$$\begin{array}{cc}\hfill f_1(q)& =q^{\frac{1}{12}}\underset{n=1}{\overset{\mathrm{}}{}}(1q^{2n})=\eta (q^2)\hfill \\ \hfill f_2(q)& =\sqrt{2}q^{\frac{1}{12}}\underset{n=1}{\overset{\mathrm{}}{}}(1+q^{2n})=\sqrt{\frac{\theta _2(q^2)}{\eta (q^2)}}\hfill \\ \hfill f_3(q)& =q^{\frac{1}{24}}\underset{n=1}{\overset{\mathrm{}}{}}(1+q^{2n+1})=\sqrt{\frac{\theta _3(q^2)}{\eta (q^2)}}\hfill \\ \hfill f_4(q)& =q^{\frac{1}{24}}\underset{n=1}{\overset{\mathrm{}}{}}(1q^{2n+1})=\sqrt{\frac{\theta _4(q^2)}{\eta (q^2)}}\hfill \end{array}$$
Then we find:
$$\begin{array}{cc}\hfill _0^{\mathrm{}}𝑑l\mathrm{D3},+|e^{lH_c}|\mathrm{D3},+=& _0^{\mathrm{}}\frac{dt}{2t}\mathrm{tr}_{NSR}\left(\frac{1+(1)^F}{2}\frac{1+R}{2}e^{2tH_0}\right)\hfill \\ \hfill =& \frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\{\frac{f_3(\stackrel{~}{q})^8f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\hfill \\ & +4\frac{f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\}\hfill \end{array}$$
The amplitude between any other fractional brane and itself is given by the same expression. Using the abstruse identity $`f_3(\stackrel{~}{q})^8f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8=0`$, we see that this amplitude vanishes, as expected.
Next, consider the amplitude:
$$\begin{array}{cc}\hfill _0^{\mathrm{}}𝑑l\mathrm{D3},+|e^{lH_c}|\mathrm{D3},=& _0^{\mathrm{}}\frac{dt}{2t}\mathrm{tr}_{NSR}\left(\frac{1+(1)^F}{2}\frac{1R}{2}e^{2tH_0}\right)\hfill \\ \hfill =& \frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\{\frac{f_3(\stackrel{~}{q})^8f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\hfill \\ & 4\frac{f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\}\hfill \end{array}$$
This gives the force between $`1_f`$ and $`1_f^{}`$, and is also equal to zero.
The next amplitude of interest is:
$$\begin{array}{cc}\hfill _0^{\mathrm{}}𝑑l\mathrm{D3},+|e^{lH_c}|\overline{\mathrm{D}}3,=& _0^{\mathrm{}}\frac{dt}{2t}\mathrm{tr}_{NSR}\left(\frac{1(1)^F}{2}\frac{1+R}{2}e^{2tH_0}\right)\hfill \\ \hfill =& \frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\{\frac{f_3(\stackrel{~}{q})^8+f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\hfill \\ & +4\frac{f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4+f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\}\hfill \end{array}$$
This can be simplified to:
$$\frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\left\{2\frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}+8\frac{f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\right\}$$
showing that it is strictly positive. Thus the force is attractive. This is as expected, since this amplitude describes the force between $`1_f`$ and $`\overline{1}_f`$.
Finally we evaluate the last independent amplitude:
$$\begin{array}{cc}\hfill _0^{\mathrm{}}𝑑l\overline{\mathrm{D}}3,+|e^{lH_c}|\mathrm{D3},+=& _0^{\mathrm{}}\frac{dt}{2t}\mathrm{tr}_{NSR}\left(\frac{1(1)^F}{2}\frac{1R}{2}e^{2tH_0}\right)\hfill \\ \hfill =& \frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\{\frac{f_3(\stackrel{~}{q})^8+f_4(\stackrel{~}{q})^8f_2(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\hfill \\ & 4\frac{f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4+f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\}\hfill \end{array}$$
which simplifies to:
$$\begin{array}{cc}& \frac{v^{(4)}}{32(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\left\{2\frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}8\frac{2f_4(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_1(\stackrel{~}{q})^4f_2(\stackrel{~}{q})^4}\right\}\hfill \\ & =\frac{v^{(4)}}{16(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\left[14\frac{f_1(\stackrel{~}{q})^4f_3(\stackrel{~}{q})^4}{f_2(\stackrel{~}{q})^4f_4(\stackrel{~}{q})^4}\right]\hfill \end{array}$$
Inserting the identity $`f_2(\stackrel{~}{q})^2f_3(\stackrel{~}{q})^2f_4(\stackrel{~}{q})^2=2`$, this amplitude can be rewritten:
$$\frac{v^{(4)}}{16(2\pi )^4}_0^{\mathrm{}}\frac{dt}{t^3}\frac{f_4(\stackrel{~}{q})^8}{f_1(\stackrel{~}{q})^8}\left[1\theta _3(\stackrel{~}{q}^2)^4\right]$$
and since $`\theta _3>1`$ for real arguments, it follows that the integrand is strictly negative, implying that the force between the $`1_f`$ and $`\overline{1^{}}_f`$ is repulsive. This confirms our claim in a previous section that the force between an adjacent suspended brane-antibrane pair is repulsive.
To summarise, we have shown that the force between two fractional $`1_f`$ branes, or between $`1_f`$ and $`1_f^{}`$, is zero, as expected. The force between $`1_f`$ and $`\overline{1}_f`$ does not vanish, but is instead attractive, as one expects for a brane-antibrane pair. Finally, the force between a $`1_f`$ brane and a $`\overline{1}_f^{}`$ is repulsive.
9. Concluding Remarks
From the variety of models considered in this paper, a number of general physical observations and open questions emerge. The fascinating structure of brane-antibrane pairs and unstable D-branes in type II string theories, and their various decay modes, are likely to be a pointer towards more fundamental structures underlying string theory. The systems we consider, with intersecting branes, lead to an even more complex picture of decay modes and interactions.
Combining some of the ingredients we have discussed could lead to dynamically stable non-BPS brane configurations and associated field theories. For example, one could imagine putting together parallel branes, which attract, and adjacent branes, which repel, in a variety of ways.
Various dual configurations exist with the Borromean property, though we have not described them here. They presumably merit careful investigation. The model that we have considered here has Bose-Fermi degeneracy in its spectrum through one-loop order. A curious occurrence of Bose-Fermi degeneracy in a very different class of non-BPS systems was described in Ref.\[32,,33\], and it would be interesting to know if there is any common thread that links these brane configurations and field theories.
Acknowledgements:
We would like to thank Atish Dabholkar, Keshav Dasgupta, Nick Dorey, Kimyeong Lee, Soo-Jong Rey, Ashoke Sen, Sandip Trivedi, Gerard Watts and Piljin Yi for helpful discussions. D.T. would like to express his gratitude to the Tata Institute of Fundamental Research, the Mehta Research Institute and the Korean Institute for Advanced Study for their kind hospitality. D.T. is supported by an EPSRC fellowship.
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warning/0001/gr-qc0001001.html | ar5iv | text | # Can Quantum Cosmology Give Observational Consequences of Many-Worlds Quantum Theory?
## Introduction
Ever since Hugh Everett III formulated his many-worlds alternative everett57 ; dewitt73 to the Copenhagen version of quantum theory, there has been considerable discussion of its merits. Many people, including some of the original supporters of the many-worlds version, have expressed the opinion that the many-worlds version is empirically indistinguishable from the Copenhagen version, so that the difference is merely metaphysical.
For example, in the first wide popularization of the many-worlds or Everett-Wheeler-Graham (EWG) version of quantum theory, Everett’s bulldog Bryce DeWitt stated dewitt70 , “Clearly the EWG view of quantum mechanics leads to experimental predictions identical with those of the Copenhagen view.”
Everett’s Ph.D. supervisor John Wheeler, who initially supported the many-worlds version wheeler57 , has recently summarized it as follows wheeler98 : “Does it offer any new insights? Does it predict outcomes of experiments that differ from outcomes predicted in conventional quantum theory? The answer to the first question is emphatically yes. The answer to the second question is emphatically no.”
Roland Omnès, though never a supporter of the many-worlds version to my knowledge, has thought about it deeply and concluded omnes94 , “If quantum mechanics were absolutely true and Everett were right, no experiment would be able to confirm or reject it. …It is not science because no experiment can show it to be wrong.”
However, David Deutsch has argued deutsch85 that the many-worlds version of quantum theory would be confirmed if an observer could “split” into two copies which make different observations, remember that they observed but not what they observed, and then are rejoined coherently. Although there are conceptual loopholes (such as claiming after the experiment that the observer’s memory of having made a definite observation is merely a false memory), I believe this argument is fairly strong evidence that the difference between the Copenhagen and many-worlds versions of quantum theory is, in principle at least, a matter that could be experimentally tested. Nevertheless, this proposed test appears to be technically very difficult.
Because of the difficulties of Deutsch’s proposed experiment, here I wish to raise the possibility that quantum cosmology might in principle lead to empirical distinctions between the Copenhagen and the many-worlds versions of quantum theory.
By the Copenhagen version, I essentially mean what I might more accurately call a single-history version, in which quantum theory gives probabilities for various alternative sequences of events, but only one sequence actually occurs. Each such alternative sequence might be called a “history” or a “world.”
In the many-worlds version, in contrast, all of the possible histories or worlds with nonzero quantum probabilities actually occur, with the quantum probabilities being not probabilities for the histories to be actualized (since all are), but instead essentially measures for the magnitude of the existence of the various histories.
## Consequences of Different Numbers of Observers
There can be significant differences in typical observations if the number of observers varies greatly from “world” to “world.” Consider the following toy models:
Quantum Cosmology Model I
World 1: Observers; measure or probability $`10^{100}`$
World 2: No observers; measure or probability $`110^{100}`$
In a single-history version of this Model I, World 1 is very improbable to occur at all, so any observation would be strong evidence against the single-history version. In a many-worlds version, World 1 does occur, so observations are not evidence against that theory.
Quantum Cosmology Model II
World A: $`10^{10}`$ observers during collapse; measure $`110^{30}`$
World B: $`10^{90}`$ observers during expansion; measure $`10^{30}`$
In a single-history version of Model II, World B is very improbable, so a random observation should expect to see a collapsing universe, Hubble constant $`H<0`$, and the probability that $`H>0`$ is observed is only $`10^{30}`$.
In contrast, in a many-worlds version of Model II, all of the observations occur, with measures presumably given by something like the expectation values of positive operators each associated with a corresponding observation page95 ; page96 . I shall assume that the observers are sufficiently similar that the total measure of a certain set of observations in a certain world (e.g., of whether the universe is expanding) is roughly proportional to the total number of observers in that world who make the observation, multiplied by the quantum measure of that world. I shall also assume that the fraction of observers who do observe whether the universe is expanding or contracting is the same in both World A and World B.
Then the total measure for World A observations of a collapsing universe is roughly proportional to the $`10^{10}`$ observers times the quantum measure of nearly unity for that world, or $`10^{10}`$, whereas the total measure for World B observations of an expanding universe is roughly proportional to $`10^{90}`$ observers of that world times the quantum measure of $`10^{30}`$ for that world, or $`10^{60}`$. Thus a random observation chosen from the sample of all existing observations in the many-worlds version is about $`10^{50}`$ times more likely to be from World B, seeing $`H>0`$, than it is to be from World A, seeing $`H<0`$, a situation qualitatively the reverse of the relative probabilities in a single-history version, such as the Copenhagen version of quantum theory.
Now if one accepted the basic quantum measures of the two worlds in Quantum Cosmology Model II but was not very certain whether a single-history or a many-worlds version of quantum theory were correct, then if one made an observation of whether the universe were expanding or contracting, it would give strong evidence as to which version is correct.
One way to explain the difference between sampling a random observation in single-history versus many-worlds quantum theories is with lottery tickets. Suppose that we have a quantum cosmological model with the following two worlds:
World 1: $`N_1`$ observers; quantum measure or probability $`p_1`$
World 2: $`N_2`$ observers; quantum measure or probability $`p_2`$
The single-history version of quantum theory is like assigning lottery tickets to World 1 and World 2 in the ratio $`p_1:p_2`$. Then a lottery ticket is chosen at random to select which world, and its observers, exist.
The many-worlds version of quantum theory is like assigning lottery tickets to each observer in World 1 and 2 with ratio $`p_1:p_2`$, so that the ratio of the total number of lottery tickets in world 1 to that in world 2 is $`N_1p_1:N_2p_2`$. All the observers exist, but with different measures for their reality, analogous to holding different numbers of lottery tickets. Choosing a measure-weighted observer (or, better, observation) at random is analogous to choosing a lottery ticket at random. The choice really is not made (since all observations really exist in the many-worlds version), but for saying which observations are typical, is is helpful to imagine their being chosen randomly.
## Preliminary Evidence from Hartle-Hawking
We cannot yet calculate probabilities for our observations from an accepted model of the quantum cosmology quantum measures, so we cannot yet perform a definitive test of whether the single-history or the many-worlds version of quantum theory is correct. However, we can examine some highly speculative preliminary suggestions from the Hartle-Hawking ‘no-boundary’ proposal hawking82 ; hartle83 ; hawking84 applied to a $`k=+1`$ Friedmann-Robertson-Walker model with a minimally coupled massive scalar field (potential $`\frac{1}{2}m^2\varphi ^2`$).
In this minisuperspace model, an approximation of the stationary phase approximation for the path integral in which the scalar field starts at a value $`\varphi _i`$ large compared with the Planck value (unity here) leads to the universe nucleating with initial size
$$a_i^2=\frac{3}{4\pi m^2\varphi _i^2}=\frac{p}{\pi }$$
(1)
and quantum measure roughly proportional to
$$e^{2S_E}e^{\pi a_i^2}=e^p$$
(2)
with $`p\pi a_i^2`$. Observations suggest $`m10^6`$ linde90 .
This is actually a measure density, and it is not clear what the prefactor should be. One simple choice is $`dp=2\pi a_ida_1`$. The resulting measure would diverge if integrated to $`p=\mathrm{}`$ or $`a_i=\mathrm{}`$, but this would correspond to $`\varphi _i=0`$, where the approximation is invalid. To get an inflationary solution, one needs $`\varphi _i>\varphi _{\mathrm{min}}1`$, so $`a_i<a_m=\sqrt{3/4\pi }/(m\varphi _{\mathrm{min}})1/m`$ or $`p<p_m=3/(4m^2\varphi _i^2)1/m^2`$. Cut off the measure density there and normalize it, so we get the simple idealization
$$P(p<p^{})\frac{e^p^{}1}{e^{p_m}1}e^{p_m}(e^p^{}1)$$
(3)
for $`p^{}\pi a_i^2<p_m\pi a_m^21/m^210^{12}1`$.
After the universe nucleates, it undergoes slow-roll inflation with $`\varphi `$ decreasing from $`\varphi _i`$ to $`\varphi _e1`$ and the volume increasing to
$$V_e=V_i\left(\frac{a_e}{a_i}\right)^3\frac{\sqrt{27\pi }}{4m^3\varphi _i^3}e^{6\pi (\varphi _i^2\varphi _e^2)}=2\sqrt{\pi }e^{6\pi \varphi _e^2}p^{3/2}\mathrm{exp}\frac{4.5\pi }{m^2p}p^{3/2}\mathrm{exp}\frac{4.5\pi }{m^2p},$$
(4)
which implies that for $`m^3V_e1`$,
$$p\frac{4.5\pi /m^2}{\mathrm{ln}(m^3V_e)+1.5\mathrm{ln}\mathrm{ln}(m^3V_e)}.$$
(5)
The entropy density after reheating is
$$s_eT_e^3\rho _e^{3/4}(m^2\varphi _e^2)^{3/4}m^{3/2}10^9.$$
(6)
By comparison, the entropy density of radiation today is
$$s_0\frac{86\pi ^2}{165}T_0^31.22\times 10^{95}.$$
(7)
Assuming essentially adiabatic expansion after reheating, one gets that the volume of the universe today is
$$V_0\frac{s_e}{s_0}V_e10^{95}m^{3/2}V_e10^{95}m^{3/2}p^{3/2}\mathrm{exp}\frac{4.5\pi }{m^2p}10^{86}p^{3/2}e^{1.4\times 10^{13}/p}.$$
(8)
Now to get something analogous to Quantum Cosmology Model II above, we need to consider what values of $`p`$ give observers mainly seeing the universe either contracting or expanding, and how many observers are produced as a function of $`p`$.
Let us make the crude assumption that observers require a universe of an age at least of the order of $`10^{60}`$, a tenth of the age of our actual universe, and hence a volume of the order of $`10^{181}`$, in order for suitable habitats to have evolved (e.g., planets around stars). This would give a lower limit on the volume at the end of inflation of about
Ve>
1086m3/21095.similar-tosubscript𝑉𝑒>
superscript1086superscript𝑚32superscript1095V_{e}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$>$\hss}
\lower 5.0pt\vbox{\hbox{$\sim$}}}}\ }10^{86}m^{-3/2}\sim 10^{95}. (9)
Inserting this back into the approximate relation between $`V_e`$ and $`p`$ gives
$$p<p_{\mathrm{max}}\frac{4.5\pi /m^2}{\mathrm{ln}(10^{77})+1.5\mathrm{ln}\mathrm{ln}(10^{77})}\frac{4.5\pi /m^2}{185}7.6\times 10^{10}$$
(10)
as the crude condition for the existence of observers.
However, if $`p`$ is sufficiently near this upper limit $`p_{\mathrm{max}}`$ for the existence of observers, then the universe will just barely last long enough for them, and they will mostly exist near the end of the lifetime of the universe, when it is collapsing. For most observers to see the universe expanding, $`V_e`$ must be sufficiently larger that the lifetime of the universe is long enough for most observers to exist while the universe is still expanding. If the present age of the universe is a typical time for observers, then one might estimate that the universe must still be expanding at an age of roughly $`10^{61}`$ for most observers to see the universe expanding, and hence for it to have a volume of at least of the order of $`10^{184}`$ then. This leads us to Ve>
1089m3/21098similar-tosubscript𝑉𝑒>
superscript1089superscript𝑚32superscript1098V_{e}{\ \lower-1.2pt\vbox{\hbox{\hbox to0.0pt{$>$\hss}
\lower 5.0pt\vbox{\hbox{$\sim$}}}}\ }10^{89}m^{-3/2}\sim 10^{98} and
$$p<p_{\mathrm{exp}}\frac{4.5\pi /m^2}{\mathrm{ln}(10^{80})+1.5\mathrm{ln}\mathrm{ln}(10^{80})}\frac{4.5\pi /m^2}{192}7.4\times 10^{10}$$
(11)
as the crude condition for most observers to see the universe expanding.
In other words, in the Hartle-Hawking minisuperspace model under consideration, if $`0<p<p_{\mathrm{exp}}7.4\times 10^{10}`$, observers will exist and will mostly see the universe expanding; if $`p_{\mathrm{exp}}<p<p_{\mathrm{max}}7.6\times 10^{10}`$, observers will exist but will mostly see the universe contracting; and if $`p_{\mathrm{max}}<p`$, essentially no observers will exist.
First, consider a Copenhagen or other single-history version of this quantum minisuperspace model, in which the wavefunction collapses to give a single macroscopic history or world, a classical Friedmann-Robertson-Walker universe characterized by $`\varphi _i`$, $`a_i`$, or $`p`$.
Using the results above, the probability for the wavefunction to collapse to classical universe that lasts long enough for observers is
$$P(\mathrm{observers})\frac{e^{p_{\mathrm{max}}}1}{e^{p_m}1}e^{p_{\mathrm{max}}p_m}10^{401000000000},$$
(12)
and the probability for it to have observers that mostly see the universe expanding is
$$P(\mathrm{observers}\mathrm{seeing}\mathrm{expansion})e^{p_{\mathrm{exp}}p_m}10^{402000000000},$$
(13)
both of which are utterly tiny.
Therefore, unless one had an uncertainty less than roughly $`e^{0.924\times 10^{12}}10^{40100000000}`$ that this single-history model was correct, the evidence that observers exist would be overwhelming evidence against it.
Even if one somehow claimed that observers were necessary (i.e., that the wavefunction could not collapse to a world with no observers), the conditional probability of a world with observers mostly seeing the universe expand, given the condition that observers exist, is only
$$P(\mathrm{observers}\mathrm{seeing}\mathrm{expansion}|\mathrm{observers}\mathrm{exist})e^{p_{\mathrm{exp}}p_{\mathrm{max}}}10^{1000000000}.$$
(14)
Thus the observation that the universe is expanding would be strong evidence against the single-history version of this model.
On the other hand, if one takes a many-worlds version of this Hartle-Hawking minisuperspace model, all components of the wavefunction exist that have positive measure, no matter how small, so observers will exist in a generalization of the model that allows sufficient structure for observers. Therefore, the existence of observers would not be evidence against a many-worlds version of a model sufficiently general to allow observers within at least some components of the wavefunction.
However, we can still ask for the relative probabilities of observations that the universe is contracting or expanding. In a many-worlds version, this will be roughly proportional to the bare quantum probability for each world multiplied by the the number of observers for that world. The unnormalized bare quantum probability was given above as $`dPe^pdp`$ for $`0<p<p_m1/m^210^{12}`$. For $`0<p<p_{\mathrm{max}}`$, observers exist, and it is reasonable to assume that the number of them is very roughly proportional to the volume $`V_0`$ of the universe when the entropy density is roughly the value of $`10^{95}`$ that we observe, assuming that our observations of this quantity are typical. Then in this range of $`p`$, we get an unnormalized observational probability density roughly proportional to
$$dP_{\mathrm{obs}}V_0dPV_0e^pdp10^{95}m^{3/2}p^{3/2}\mathrm{exp}\left(\frac{4.5\pi }{m^2p}+p\right)dp.$$
(15)
The integral of this over $`0<p<p_{\mathrm{max}}`$ diverges for $`p0`$ ($`V_0\mathrm{}`$), because of the divergence in the number of observers there, so effectively all of the observational probability occurs at that limit (infinitely large universes with infinitely many observers, presumably almost all seeing the universe expanding, since the stars would have all burned out in the infinite time it takes the infinitely large universe to recollapse.)
Therefore, in a many-worlds version of this Hartle-Hawking quantum cosmological model, one would presumably expect with very nearly unit probability that a random observation would see the universe expanding, the opposite of one’s expectation for a single-history version of the same model. Thus if one accepted the basic model and allowed at least some reasonable uncertainty as to whether a single-history or a many-worlds version of the model is correct (before considering the evidence of the sign of the Hubble constant), then one’s observation of whether the universe is expanding or contracting would give very strong evidence in support of either the many-worlds or the single-history version respectively. This is very similar qualitatively to the toy Quantum Cosmology Model II discussed above, except here the quantitative differences are even grossly more severe.
Of course, this preliminary evidence from a particular implementation of the Hartle-Hawking no-boundary proposal is highly speculative and is meant to be mainly illustrative, because of the many uncertainties of the model.
## Conclusions
If the amount of observations (roughly, the number of observers) varies for different wavefunction components, then observation probabilities depend on whether only one component occurs in actuality (a single-history version of quantum theory, where observations truly are made only in one history, world, or component of the wavefunction), or whether many do (a many-worlds version, where observations truly are made in many histories, worlds, or components of the wavefunction). In particular, if components with relatively few observations dominate the quantum amplitude, but other components with testably different observations dominate the expectation value of the number of observations, which observations are most probable varies between single-history and many-worlds quantum theories.
The Hartle-Hawking wavefunction might allow a test from the observed expansion of the universe, but as of now it is highly speculative whether it is correct and what relative probabilities it would give for observing the universe expanding.
I acknowledge helpful discussions with Meher Antia, Jerry Finkelstein, Valeri Frolov, Jim Hartle, Jacques Mallah, and William Unruh. This research was supported in part by the Natural Sciences and Engineering Research Council of Canada.
A shorter version of this paper has been circulated page99 and has been reported on in the lay literature antia99 . |
warning/0001/gr-qc0001067.html | ar5iv | text | # The Interaction of Dirac Particles with Non-Abelian Gauge Fields and Gravity – Bound States
## 1 Introduction
The coupling of gravity to other classical force fields and to quantum mechanical particles has led to many interesting solutions of Einstein’s equations and has given some insight into the nature of the nonlinear interactions. The first such examples are the Bartnik-McKinnon (BM) solutions of the Einstein-Yang/Mills (EYM) equations . For these solutions, the repulsive Yang-Mills force compensates the attractive gravitational force; unfortunately, these solutions are unstable . If, on the other hand, one considers quantum mechanical Dirac particles, the gravitational attraction is balanced by the repulsion due to the Heisenberg Uncertainty Principle, and this leads to stable bound states of the resulting Einstein-Dirac (ED) system . However, a pure ED system is somewhat artificial, because physical Dirac particles also exhibit electroweak and strong interactions, which in all realistic situations are much stronger than gravity. Thus the question arises if Dirac particles in a gravitational field still form bound states if an additional strong coupling to a non-abelian Yang-Mills (YM) field is taken into account (the case of an abelian gauge field was considered in ). Related questions are, do the BM solutions become stable if one includes Dirac particles, and which physical effects does the nonlinear coupling in the Einstein-Dirac-Yang/Mills (EDYM) equations lead to?
In order to address these questions, we consider here a static, spherically symmetric EDYM system of one Dirac particle in a gravitational field and an $`SU(2)`$ Yang-Mills field. In this system, the spinors couple only to the magnetic component of the YM field, and we thus obtain a consistent ansatz by setting the electric component identically equal to zero. In the limit of weak coupling of the spinors, our system goes over to the EYM system . In contrast to the two-particle singlet state studied in , we consider here only one Dirac particle (this becomes possible because the inclusion of the $`SU(2)`$ Yang-Mills field changes the representation of the rotation group on the spinors; see Section 2). Thus one cannot recover exactly the ED system , but the limit of a weak Yang-Mills field yields equations which are closely related to the ED equations of the two-particle singlet.
By numerically seeking bound states of the EDYM system, we find a surprisingly rich solution structure. First of all, we find solutions where the Dirac particle is bound by the gravitational attraction, and where the Dirac particle also generates a YM field. Stable solutions of this type exist also in the physically realistic situation of weak gravitational and strong YM coupling. This result shows that the magnetic component of the YM field, which usually has a repulsive effect (like e.g. for the BM solutions) cannot prevent Dirac particles from forming stable bound states. We also find other types of solutions where the binding comes about through the nonlinear interaction of all three fields. These solutions have stable and unstable branches, whereby the stable solutions are found for weak gravitational coupling, provided that the YM coupling is sufficiently strong, but not too strong. Finally, we study the relation between these solutions and the BM solutions. We find one-parameter families of solutions which join the BM ground state with our new solutions. This shows that the BM ground state can be made stable by the inclusion of a Dirac particle, but only if the coupling to the spinors is sufficiently strong. The first excited BM state, on the other hand, cannot be joined with our new solutions. Namely, perturbing this state by an additional Dirac particle yields a separate unstable branch of EDYM solutions.
The plan of the paper is as follows. In Section 2 we derive the $`SU(2)`$-EDYM equations. In Section 3 we obtain a limiting system constructed by letting the gravitational constant tend to zero and letting the YM coupling constant tend to infinity. We find numerical solutions of this system and discuss their properties. In the last section, we consider solutions of the full EDYM equations, obtained by tracing the solutions of our limiting system and the BM solutions while continuously varying the coupling constants.
## 2 Derivation of the Equations
The general EDYM equations are obtained by variation over Lorentzian metrics $`g_{ij}`$, YM connections $`𝒜`$, and Dirac wave functions $`\mathrm{\Psi }`$, of the action
$$S=\left(\frac{1}{16\pi \kappa }R+\overline{\mathrm{\Psi }}(Gm)\mathrm{\Psi }\frac{1}{16\pi e^2}\text{Tr}(F_{ij}F^{ij})\right)\sqrt{detg}d^4x,$$
(2.1)
where $`R`$ is scalar curvature, $`G`$ is the Dirac operator (which depends on the gravitational and YM fields), $`F_{ij}`$ is the YM field tensor, and where the trace is taken over the YM index. The gravitational and YM coupling constants are denoted by $`\kappa `$ and $`e`$, respectively. The appearance of the factor $`e^2`$ in (2.1) requires a brief explanation. In contrast to the usual form of the gauge-covariant derivative $`D_j=_jieA_j`$, we use here the convention $`D_j=_jiA_j`$ (this makes it possible to work with the particularly convenient form of the gauge potentials used in ). Our convention is obtained from the usual one by the transformation $`A_je^1A_j`$. Under this transformation, the field strength tensor behaves like $`F_{ij}e^1F_{ij}`$, and this gives rise to the factor $`e^2`$ in (2.1).
In this paper, we shall study a particular EDYM system, which is obtained as follows. First of all, we consider a spherically symmetric, static metric in polar coordinates,
$$ds^2=\frac{1}{T(r)^2}dt^2\frac{1}{A(r)}dr^2r^2d\vartheta ^2r^2\mathrm{sin}^2\vartheta d\phi ^2,$$
(2.2)
with positive functions $`A`$ and $`T`$. The Einstein tensor corresponding to this metric is given in . Moreover, as in , we restrict attention to the magnetic component of an $`SU(2)`$ Yang-Mills field and choose for the YM potential the ansatz
$$𝒜=w(r)\tau ^1d\vartheta +(\mathrm{cos}\vartheta \tau ^3+w(r)\mathrm{sin}\vartheta \tau ^2)d\phi $$
(2.3)
with a real function $`w`$, where $`\stackrel{}{\tau }=\frac{1}{2}\stackrel{}{\sigma }`$ is the standard basis of $`su(2)`$ ($`\stackrel{}{\sigma }`$ are the Pauli matrices). The YM current $`j`$ and energy-momentum tensor $`T_j^i=\frac{1}{4\pi e^2}\text{Tr }(F^{ia}F_{ja}\frac{1}{4}F^{ab}F_{ab}\delta _j^i)`$ are computed to be
$`j`$ $`=`$ $`{\displaystyle \frac{1}{4\pi e^2}}\left({\displaystyle \frac{A}{2r^2}}w^{\prime \prime }{\displaystyle \frac{A^{}T2AT^{}}{4r^2T}}w^{}{\displaystyle \frac{w(1w^2)}{2r^4}}\right)\left(\sigma ^1{\displaystyle \frac{}{\vartheta }}+\sigma ^2\mathrm{csc}(\vartheta ){\displaystyle \frac{}{\phi }}\right)`$
$`T_0^0`$ $`=`$ $`{\displaystyle \frac{1}{4\pi e^2}}\left({\displaystyle \frac{2}{r^4}}(1w^2)^2{\displaystyle \frac{4}{r^2}}Aw^2\right)`$ (2.4)
$`T_1^1`$ $`=`$ $`{\displaystyle \frac{1}{4\pi e^2}}\left({\displaystyle \frac{2}{r^4}}(1w^2)^2+{\displaystyle \frac{4}{r^2}}Aw^2\right)`$ (2.5)
$`T_2^2`$ $`=`$ $`T_3^3={\displaystyle \frac{1}{4\pi e^2}}\left({\displaystyle \frac{2}{r^4}}(1w^2)^2\right),`$ (2.6)
and all other components vanish. When coupled to the Dirac spinors, the YM potential (2.3) has the disadvantage that it depends on $`\vartheta `$ and $`\phi `$, in a way which makes it difficult to separate variables in the Dirac equation. To remedy this, we perform the $`SU(2)`$ gauge transformation $`𝒜_jU𝒜_jU^1+iU(_jU^1)`$ with
$$U(\vartheta ,\phi )=\mathrm{exp}\left(i\phi \tau ^1\right)\mathrm{exp}\left(i(\vartheta +\pi )\tau ^3\right)\mathrm{exp}\left(\frac{i\pi }{2}\tau ^2\right).$$
The resulting YM potential is
$$𝒜=(w1)r\mathrm{sin}\vartheta (\tau ^\phi d\vartheta \tau ^\vartheta d\phi ),$$
(2.7)
where we use the following “polar” linear combinations of the $`\tau `$ matrices,
$`\tau ^r`$ $`=`$ $`\tau ^1\mathrm{cos}\vartheta +\tau ^2\mathrm{sin}\vartheta \mathrm{cos}\phi +\tau ^3\mathrm{sin}\vartheta \mathrm{sin}\phi `$
$`\tau ^\vartheta `$ $`=`$ $`{\displaystyle \frac{1}{r}}\left(\tau ^1\mathrm{sin}\vartheta +\tau ^2\mathrm{cos}\vartheta \mathrm{cos}\phi +\tau ^3\mathrm{cos}\vartheta \mathrm{sin}\phi \right)`$
$`\tau ^\phi `$ $`=`$ $`{\displaystyle \frac{1}{r\mathrm{sin}\vartheta }}\left(\tau ^2\mathrm{sin}\phi +\tau ^3\mathrm{cos}\phi \right).`$ (2.8)
By minimally coupling the $`SU(2)`$ potential (2.7) to the Dirac operator in the gravitational field \[3, Eq. (2.23)\], we obtain the Dirac operator
$`G`$ $`=`$ $`iT\gamma ^t_t+\gamma ^r\left(i\sqrt{A}_r+{\displaystyle \frac{i}{r}}(\sqrt{A}1){\displaystyle \frac{i}{2}}\sqrt{A}{\displaystyle \frac{T^{}}{T}}\right)+i\gamma ^\vartheta _\vartheta +i\gamma ^\phi _\phi `$ (2.9)
$`+{\displaystyle \frac{2i}{r}}(w1)(\stackrel{}{\gamma }\stackrel{}{\tau }\gamma ^r\tau ^r)\tau ^r,`$
where $`(\gamma ^j)_{j=t,r,\vartheta ,\phi }`$ are, in analogy to (2.8), the Dirac matrices of Minkowski space in polar coordinates, where we work in the Dirac representation,
$$\gamma ^t=\left(\begin{array}{cc}\text{1 1}& 0\\ 0& \text{1 1}\end{array}\right),\stackrel{}{\gamma }=\left(\begin{array}{cc}0& \stackrel{}{\sigma }\\ \stackrel{}{\sigma }& 0\end{array}\right).$$
(2.10)
Notice that the Dirac operator (2.9) acts on 8-component wave functions; this is because the additional YM index doubles the number of components compared to usual Dirac spinors. More precisely, it is convenient to regard the wave functions as sections of
$$\text{}\text{C}^8=\text{}\text{C}_{\text{up/down}}^2\text{}\text{C}_{\text{large/small}}^2\text{}\text{C}_{\text{YM}}^2,$$
(2.11)
where $`\text{}\text{C}_{\text{up/down}}^2`$ describes the two spin orientations, $`\text{}\text{C}_{\text{large/small}}^2`$ corresponds to the upper and lower components of the Dirac spinor (i.e., usual Dirac spinors are sections of $`\text{}\text{C}^4=\text{}\text{C}_{\text{up/down}}^2\text{}\text{C}_{\text{large/small}}^2`$), and $`\text{}\text{C}_{\text{YM}}^2`$ is acted upon by the $`SU(2)`$ gauge group. For clarity, we shall refer to the factors in (2.11) by separate indices, i.e. we write a wave function $`\mathrm{\Psi }`$ as $`(\mathrm{\Psi }^{\alpha ua})_{\alpha ,u,a=1,2}`$, where $`\alpha `$, $`u`$, and $`a`$ label the components of $`\text{}\text{C}_{\text{up/down}}^2`$, $`\text{}\text{C}_{\text{large/small}}^2`$, and $`\text{}\text{C}_{\text{YM}}^2`$, respectively. Thus the operators $`\stackrel{}{\tau }`$ act on the index $`a`$, the spin operators $`\stackrel{}{S}`$ are given by $`\stackrel{}{S}=\frac{1}{2}\stackrel{}{\sigma }`$ acting on the Greek indices, and $`\gamma ^t`$ coincides with the matrix $`\gamma ^t=\text{diag}(1,1)`$ acting on the index $`u`$, i.e.
$$\gamma ^t\mathrm{\Psi }^{\alpha ua}=\{\begin{array}{cc}\mathrm{\Psi }^{\alpha ua}& \text{if }u=1\hfill \\ \mathrm{\Psi }^{\alpha ua}& \text{if }u=2\hfill \end{array}.$$
It is apparent in (2.9) that the Dirac operator commutes with the three operators
$$\stackrel{}{J}=\stackrel{}{L}+\stackrel{}{S}+\stackrel{}{\tau },$$
(2.12)
where $`\stackrel{}{L}`$ is angular momentum. It is very convenient to regard the operators $`\stackrel{}{J}`$ as the total angular momentum operators of the system. Since the total angular momentum operators are the infinitesimal generators of rotations (as explained for angular momentum in \[5, par. 26\]), we can then interpret (2.12) as saying that the inclusion of the YM field influences the representation of the rotation group on the spinors. The Dirac operator is invariant under this group representation, because the operators $`\stackrel{}{J}`$ commute with $`G`$; this makes spherical symmetry of the Dirac operator manifest.
Since (2.12) coincides with the formula for the addition of angular momentum $`\stackrel{}{L}`$ to two spin-$`\frac{1}{2}`$-operators $`\stackrel{}{S}`$ and $`\stackrel{}{\tau }`$, it is clear that the operators $`\stackrel{}{J}`$ have integer eigenvalues. Thus we can expect that the operator $`J^2`$ has a nontrivial kernel. In this case, the simplest spherically symmetric ansatz for the Dirac particles would be to take one Dirac particle whose wave function is in the kernel of $`J^2`$. We now work out this ansatz in detail, whereby we consider $`\stackrel{}{J}`$ as operators on the spinors $`\mathrm{\Phi }^{\alpha a}(\vartheta ,\phi )`$ on $`S^2`$ (i.e. the $`\mathrm{\Phi }^{\alpha a}`$ are sections of $`\text{}\text{C}^4=\text{}\text{C}_{\text{up/down}}^2\text{}\text{C}_{\text{YM}}^2`$). Adding the two spin operators $`\stackrel{}{S}`$ and $`\stackrel{}{\tau }`$, we can decompose $`\text{}\text{C}_{\text{up/down}}^2\text{}\text{C}_{\text{YM}}^2`$ into the direct sum of one state of total spin zero and three states of total spin one (see \[5, par. 31\]; these states are usually called the singlet and triplet states, respectively). By subsequently adding the angular momentum $`\stackrel{}{L}`$ according to the standard rules for the addition of angular momentum \[5, par. 31\], one sees that the operator $`J^2`$ has indeed a nontrivial kernel. More precisely, the kernel of $`J^2`$ is two-dimensional, spanned by two vectors $`\mathrm{\Phi }_0`$ and $`\mathrm{\Phi }_1`$ with angular momentum zero and one, respectively. The state $`\mathrm{\Phi }_0`$ is (up to a phase) uniquely characterized by the conditions
$$\stackrel{}{L}\mathrm{\Phi }_0=\mathrm{\hspace{0.33em}0}=(\stackrel{}{S}+\stackrel{}{\tau })\mathrm{\Phi }_0\text{and}\mathrm{\Phi }_0_{S^2}=\mathrm{\hspace{0.33em}1}.$$
(2.13)
Using (2.13), we can write $`\mathrm{\Phi }_1`$ as
$$\mathrm{\Phi }_1=\mathrm{\hspace{0.33em}2}S^r\mathrm{\Phi }_0=2\tau ^r\mathrm{\Phi }_0.$$
(2.14)
Namely, representing $`S^r`$ and $`\tau ^r`$ in the form
$$S^r=\stackrel{}{x}\stackrel{}{s}\text{and}\tau ^r=\stackrel{}{x}\stackrel{}{\tau },$$
and using the standard commutation relations between the components of $`\stackrel{}{L}`$, $`\stackrel{}{x}`$, and $`\stackrel{}{S}`$ (see \[5, pars. 26 and 54\]), we obtain that
$`\stackrel{}{J}\mathrm{\Phi }_1`$ $`=`$ $`2[\stackrel{}{J},S^r]\mathrm{\Phi }_0=\mathrm{\hspace{0.33em}2}[\stackrel{}{L},(\stackrel{}{x}\stackrel{}{S})]\mathrm{\Phi }_0+\mathrm{\hspace{0.25em}2}[\stackrel{}{S}+\stackrel{}{\tau },(\stackrel{}{x}\stackrel{}{S})]\mathrm{\Phi }_0`$ (2.15)
$`=`$ $`2i\stackrel{}{x}\stackrel{}{S}\mathrm{\Phi }_0+\mathrm{\hspace{0.25em}2}i\stackrel{}{x}\stackrel{}{S}\mathrm{\Phi }_0=\mathrm{\hspace{0.33em}0}`$
(where $``$ is the wedge product in $`\text{I R}^3`$), and
$$\mathrm{\Phi }_1_{S^2}^2=_{S^2}<\text{ }2S^r\mathrm{\Phi }_0,\mathrm{\hspace{0.25em}2}S^r\mathrm{\Phi }_0>d\omega =\mathrm{\Phi }_0_{S^2}^2=\mathrm{\hspace{0.33em}1}.$$
One can verify directly that $`\mathrm{\Phi }_1`$ has angular momentum one; namely
$`L^2\mathrm{\Phi }_1`$ $`=`$ $`2L^2S^r\mathrm{\Phi }_0=\mathrm{\hspace{0.33em}2}\stackrel{}{L}[\stackrel{}{L},S^r]\mathrm{\Phi }_0=2i\stackrel{}{L}(\stackrel{}{x}\stackrel{}{S})\mathrm{\Phi }_0`$
$`=`$ $`2i[\stackrel{}{L},(\stackrel{}{x}\stackrel{}{S})]\mathrm{\Phi }_0=\mathrm{\hspace{0.33em}4}(\stackrel{}{x}\stackrel{}{S})\mathrm{\Phi }_0=l(l+1)\mathrm{\Phi }_1`$
with $`l=1`$. Furthermore, using the fact that $`(2S^r)^2=1=(2\tau ^r)^2`$ and $`S^2=\frac{3}{4}=\tau ^2`$, we obtain that
$`\mathrm{\Phi }_0`$ $`=`$ $`2S^r\mathrm{\Phi }_1=2\tau ^r\mathrm{\Phi }_1`$ (2.16)
$`\stackrel{}{S}\stackrel{}{\tau }\mathrm{\Phi }_0`$ $`\stackrel{(\text{2.13})}{=}`$ $`\stackrel{}{S}\stackrel{}{S}\mathrm{\Phi }_0={\displaystyle \frac{3}{4}}\mathrm{\Phi }_0`$ (2.17)
$`\stackrel{}{S}\stackrel{}{\tau }\mathrm{\Phi }_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}((\stackrel{}{S}+\stackrel{}{\tau })^2S^2\tau ^2)\mathrm{\Phi }_1\stackrel{(\text{2.15})}{=}{\displaystyle \frac{1}{2}}(L^2S^2\tau ^2)\mathrm{\Phi }_1={\displaystyle \frac{1}{4}}\mathrm{\Phi }_1`$ (2.18)
$`\stackrel{}{S}\stackrel{}{L}\mathrm{\Phi }_0`$ $`=`$ $`0`$ (2.19)
$`\stackrel{}{S}\stackrel{}{L}\mathrm{\Phi }_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}((\stackrel{}{S}+\stackrel{}{L})^2S^2L^2)\mathrm{\Phi }_1={\displaystyle \frac{1}{2}}(\tau ^2S^2L^2)\mathrm{\Phi }_1=\mathrm{\Phi }_1.`$ (2.20)
Finally, it is useful to observe that (cf. \[6, equation (3.3)\])
$$S^\vartheta _\vartheta +S^\phi _\phi =\frac{2}{r}S^r(\stackrel{}{S}\stackrel{}{L}).$$
(2.21)
Using the relations (2.13)–(2.21), one can easily compute the Dirac operator (2.9) on the invariant subspace $`J^2=0`$. It turns out that we obtain a consistent ansatz for the Dirac wave function by setting
$$\mathrm{\Psi }^{\alpha ua}(t,r,\vartheta ,\phi )=e^{i\omega t}\frac{\sqrt{T(r)}}{r}\left(\alpha (r)\mathrm{\Phi }_0^{\alpha a}(\vartheta ,\phi )\delta _{u,1}+i\beta (r)\mathrm{\Phi }_1^{\alpha a}(\vartheta ,\phi )\delta _{u,2}\right)$$
(2.22)
with real functions $`\alpha `$ and $`\beta `$, where $`\omega >0`$ is the energy of the Dirac particle, and $`\delta _{.,.}`$ is the Kronecker delta. For this ansatz, the Dirac equation reduces to the system of ODEs
$`\sqrt{A}\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{w}{r}}\alpha (m+\omega T)\beta `$ (2.23)
$`\sqrt{A}\beta ^{}`$ $`=`$ $`(m+\omega T)\alpha {\displaystyle \frac{w}{r}}\beta .`$ (2.24)
The Dirac current $`j`$ and Dirac energy-momentum tensor $`T_{jk}=\text{Re }(\overline{\mathrm{\Psi }}G_{(j}D_{k)}\mathrm{\Psi })`$ corresponding to the ansatz (2.22) are obtained by a straightforward computation similar to that in . The result is
$`j`$ $`=`$ $`{\displaystyle \frac{2T}{r^3}}\alpha \beta \left(\sigma ^1{\displaystyle \frac{}{\vartheta }}+\sigma ^2\mathrm{csc}(\vartheta ){\displaystyle \frac{}{\phi }}\right)`$
$`T_0^0`$ $`=`$ $`{\displaystyle \frac{\omega T^2}{r^2}}(\alpha ^2+\beta ^2)`$
$`T_1^1`$ $`=`$ $`{\displaystyle \frac{\omega T^2}{r^2}}(\alpha ^2+\beta ^2)+\mathrm{\hspace{0.25em}2}{\displaystyle \frac{T}{r^3}}w\alpha \beta +{\displaystyle \frac{mT}{r^2}}(\alpha ^2\beta ^2)`$
$`T_2^2`$ $`=`$ $`T_3^3={\displaystyle \frac{T}{r^3}}w\alpha \beta ,`$
and all other components vanish. The normalization condition for the spinors is (as in ),
$$_0^{\mathrm{}}(\alpha ^2+\beta ^2)\frac{T}{\sqrt{A}}=\mathrm{\hspace{0.33em}1}.$$
(2.25)
By substituting the formulas for the YM current and energy-momentum tensor into the Einstein and YM equations, we obtain the following system of ODEs,
$`rA^{}`$ $`=`$ $`1A{\displaystyle \frac{\kappa }{e^2}}{\displaystyle \frac{(1w^2)^2}{r^2}}\mathrm{\hspace{0.25em}2}\kappa \omega T^2\left(\alpha ^2+\beta ^2\right){\displaystyle \frac{2\kappa }{e^2}}Aw^2`$ (2.26)
$`2rA{\displaystyle \frac{T^{}}{T}}`$ $`=`$ $`1+A+{\displaystyle \frac{\kappa }{e^2}}{\displaystyle \frac{(1w^2)^2}{r^2}}+\mathrm{\hspace{0.25em}2}\kappa mT\left(\alpha ^2\beta ^2\right)\mathrm{\hspace{0.25em}2}\kappa \omega T^2\left(\alpha ^2+\beta ^2\right)`$ (2.27)
$`+4\kappa {\displaystyle \frac{T}{r}}w\alpha \beta {\displaystyle \frac{2\kappa }{e^2}}Aw^2`$
$`r^2Aw^{\prime \prime }`$ $`=`$ $`\left(1w^2\right)w+e^2rT\alpha \beta r^2{\displaystyle \frac{A^{}T2AT^{}}{2T}}w^{}.`$ (2.28)
The Einstein equations are (2.26) and (2.27), whereas (2.28) is the YM equation. Our EDYM system is given by the five ODEs (2.23), (2.24), (2.26)–(2.28), together with the normalization condition (2.25).
We are interested here in bound states of the Dirac particles. Thus we want to find particle-like solutions of our EDYM system, i.e. solutions which are smooth and tend to the vacuum solution as $`r\mathrm{}`$. According to the explicit formulas (2.4)–(2.6), the energy-momentum tensor of the YM field is regular at $`r=0`$ only when $`|w(0)|=1`$ and $`w^{}(0)=0`$. Using the remaining gauge freedom, we can assume that $`w(0)=1`$, and thus
$$w(r)=\mathrm{\hspace{0.33em}1}\frac{\lambda }{2}r^2+𝒪(r^3)$$
(2.29)
with a real parameter $`\lambda `$. Using this result, a local Taylor expansion of the Einstein and Dirac equations around $`r=0`$ yields (just as in ) that
$`\alpha (r)`$ $`=`$ $`\alpha _1r+𝒪(r^3),\beta (r)={\displaystyle \frac{1}{2}}(\omega T_0m)\alpha _1r^2+𝒪(r^3)`$ (2.30)
$`A(r)`$ $`=`$ $`1+𝒪(r^2),T(r)=T_0+𝒪(r^2)`$ (2.31)
with two parameters $`\alpha _1`$ and $`T_0>0`$. Using linearity of the Dirac equation, we can always assume that $`\alpha _1>0`$. Furthermore, we demand that our solution has finite ADM mass,
$$\rho :=\underset{r\mathrm{}}{lim}\frac{r}{2\kappa }(1A(r))<\mathrm{},$$
(2.32)
and goes asymptotically to the vacuum solution,
$$\underset{r\mathrm{}}{lim}T(r)=\mathrm{\hspace{0.33em}1},\underset{r\mathrm{}}{lim}(w(r),w^{}(r))=(\pm 1,0),\underset{r\mathrm{}}{lim}(\alpha (r),\beta (r))=(0,0).$$
(2.33)
## 3 The Reciprocal Coupling Limit
Under all realistic conditions, the coupling of Dirac particles to the YM field (describing the weak or strong interactions) is much stronger than the coupling to the gravitational field. Thus we are particularly interested in the case of weak gravitational coupling. In preparation, it is instructive to briefly consider the case without gravitation. In this limit, the Dirac equations read
$`\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{w}{r}}\alpha (m+\omega )\beta `$ (3.1)
$`\beta ^{}`$ $`=`$ $`(m+\omega )\alpha {\displaystyle \frac{w}{r}}\beta .`$ (3.2)
For large $`r`$, these equations go over to a linear system of ODEs with constant coefficients, and the sign of $`m\omega `$ determines whether the solutions of these equations behave oscillatory or exponentially. The normalization condition (2.25) excludes the oscillatory case (as in \[6, Section 5\]) and thus $`m\omega 0`$ . In the case $`m\omega =0`$, the $`\beta `$-equation is independent of $`\alpha `$, and the boundary conditions (2.30) imply that $`\beta 0`$. As a consequence, (2.28) reduces to the homogeneous YM equation
$$r^2w^{\prime \prime }=(1w^2)w.$$
(3.3)
It is well-known that the only solution to this equation satisfying the boundary conditions (2.29),(2.33) is the trivial solution $`w1`$. But then the $`\alpha `$-equation simplifies to
$$\alpha ^{}=\frac{1}{r}\alpha ,$$
whose solution $`\alpha =\alpha _1r`$ violates the normalization condition (2.25). In the case $`m\omega >0`$, on the other hand, the local Taylor expansion (2.30) yields that the $`(\alpha ,\beta )`$-curve lies for small $`r`$ in the fourth quadrant, i.e. $`\beta (r)<0<\alpha (r)`$ for small $`r`$. Using the Dirac equations (3.1),(3.2), one sees that the fourth quadrant is an invariant region, and thus $`\beta (r)<0<\alpha (r)`$ for all $`r`$. But in the fourth quadrant, both $`\alpha (r)`$ and $`\beta (r)`$ are increasing for large $`r`$ (as one sees in (3.1),(3.2) taking into account that $`w/r0`$ for $`r\mathrm{}`$), and thus the normalization condition (2.25) will again be violated.
These considerations show that the gravitational field is essential for the formation of bound states. Nevertheless, for arbitrarily weak gravitational coupling, we can hope to find bound states. It is even conceivable that these bound state solutions might have a well-defined limit when the gravitational coupling tends to zero, if we let the YM coupling go to infinity at the same time. Our idea is that this limiting case might yield a system of equations which is simpler than the full EDYM system, and can thus serve as a physically interesting starting point for the analysis of the coupled interactions described by the EDYM equations. Expressed in dimensionless quantities, we shall thus consider the limits
$$m^2\kappa \mathrm{\hspace{0.33em}0}\text{and}e^2\mathrm{}.$$
(3.4)
Let us determine how the quantities of our EDYM system should behave in this limit. Since we are considering weak gravitational coupling, it is clear that the metric will be close to the Minkowski metric, i.e. $`A1`$ and $`T1`$. Furthermore, the YM potential $`w`$ should have a finite limit. Similar to our flat space consideration at the beginning of this section, one sees that the normalization condition (2.25) can be satisfied only if the function $`m\omega T(r)`$ changes sign, and thus $`\omega m`$ (but both $`m`$ and $`\omega `$ may go to zero or infinity in the limit (3.4)). Putting this information together, we conclude that the Dirac equations (2.23) and (2.24) have a meaningful limit only under the assumptions that $`\alpha `$ converges and that
$$m\beta (r)\widehat{\beta }(r),m^2(T(r)1)\phi ,m(\omega m)E$$
(3.5)
with two real functions $`\widehat{\beta }`$, $`\phi `$ and a real parameter $`E`$. Multiplying (2.24) with $`m`$ and taking the limits (3.5) as well as $`A,T1`$, the Dirac equations become
$`\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{w}{r}}\alpha \mathrm{\hspace{0.25em}2}\widehat{\beta }`$ (3.6)
$`\widehat{\beta }^{}`$ $`=`$ $`(E+\phi )\alpha {\displaystyle \frac{w}{r}}\widehat{\beta }.`$ (3.7)
We next consider the YM equation (2.28). The last term in (2.28) drops out in the limit of weak gravitational coupling (3.4). The second summand converges only under the assumption that
$$\frac{e^2}{m}q$$
(3.8)
with $`q`$ a real parameter, playing the role of an “effective” coupling constant. Together with (3.4), this implies that $`m\mathrm{}`$. The YM equations thus have the limit
$$r^2w^{\prime \prime }=(1w^2)w+qr\alpha \widehat{\beta }.$$
(3.9)
In order to get a well-defined and non-trivial limit of the Einstein equations (2.26),(2.27), we need to assume that the parameter $`m^3\kappa `$ has a finite, non-zero limit. Since this parameter has the dimension of inverse length, we can arrange by a scaling of our coordinates that
$$m^3\kappa \mathrm{\hspace{0.33em}1}.$$
(3.10)
We differentiate the $`T`$-equation (2.27) with respect to $`r`$ and substitute (2.26). Taking the limits (3.5) and (3.10), a straightforward calculation yields the equation
$$r^2\mathrm{\Delta }\phi =\alpha ^2,$$
(3.11)
where $`\mathrm{\Delta }=r^2_r(r^2_r)`$ is the radial Laplacian in Euclidean $`\text{I R}^3`$. Indeed, this equation can be regarded as Newton’s equation with the Newtonian potential $`\phi `$. Thus our limiting case (3.11) for the gravitational field corresponds to taking the Newtonian limit. Finally, the normalization condition (2.25) reduces to
$$_0^{\mathrm{}}\alpha (r)^2𝑑r=\mathrm{\hspace{0.33em}1}.$$
(3.12)
The boundary conditions (2.29)–(2.33) are transformed into
$`w(r)`$ $`=`$ $`1{\displaystyle \frac{\lambda }{2}}r^2+𝒪(r^3),\underset{r\mathrm{}}{lim}w(r)=\pm 1`$ (3.13)
$`\alpha (r)`$ $`=`$ $`\alpha _1r+𝒪(r^3),\widehat{\beta }(r)=𝒪(r^3)`$ (3.14)
$`\phi (r)`$ $`=`$ $`\phi _0+𝒪(r^2),\underset{r\mathrm{}}{lim}\phi (r)<\mathrm{}`$ (3.15)
with the three parameters $`\lambda `$, $`\alpha _1`$, and $`\phi _0`$. We point out that the limiting system contains only one coupling constant $`q`$. According to (3.8) and (3.10), $`q`$ is in dimensionless form given by
$$e^2m^2\kappa q.$$
(3.16)
Hence in dimensionless quantities, our limit (3.4) describes the situation where the gravitational coupling goes to zero, while the YM coupling constant goes to infinity like $`e^2(m^2\kappa )^1`$. Therefore, we call our limiting case the reciprocal coupling limit. The reciprocal coupling system is given by the equations (3.6), (3.7), (3.9), and (3.11) together with the normalization condition (3.12) and the boundary conditions (3.13)–(3.15). According to (3.5), the parameter $`E`$ coincides up to a scaling factor with $`\omega m`$, and thus has the interpretation as the (properly scaled) energy of the Dirac particle. As in Newtonian mechanics, the potential $`\phi `$ is determined only up to a constant $`\mu \text{I R}`$; namely, the reciprocal limit equations are invariant under the transformation
$$\phi \phi +\mu ,EE\mu .$$
(3.17)
Let us consider how the ADM mass behaves in the reciprocal coupling limit. First of all, we can write the quotient $`\rho /m`$ as
$$\frac{\rho }{m}=\underset{r\mathrm{}}{lim}\frac{r}{2\kappa m}(1A(r))=\frac{1}{m}_0^{\mathrm{}}\left(\frac{r}{2\kappa }(1A(r))\right)^{}𝑑r.$$
After substituting the $`A`$-equation (2.26), we can take the limits (3.4) and (3.5) and obtain that
$$\frac{\rho }{m}_0^{\mathrm{}}\alpha (r)^2𝑑r\stackrel{(\text{3.12})}{=}\mathrm{\hspace{0.33em}1}.$$
(3.18)
Thus the ADM mass coincides with the rest mass of the Dirac particle; this shows that the total binding energy $`B:=\rho m`$ goes to zero in our limit. Indeed, the behavior of the total binding energy can be described in more detail as follows. For a solution of the full EDYM system, we can write the binding energy using the normalization condition (2.25) as
$$B=_0^{\mathrm{}}\left(\left(\frac{r}{2\kappa }(1A)\right)^{}m(\alpha ^2+\beta ^2)\frac{T}{\sqrt{A}}\right)𝑑r.$$
We again substitute the $`A`$-equation (2.26) and obtain
$$B=_0^{\mathrm{}}\left(\frac{1}{2e^2}\frac{(1w^2)^2}{r^2}+\frac{1}{e^2}Aw^2+(wT\sqrt{A}m)\frac{T}{\sqrt{A}}(\alpha ^2+\beta ^2)\right)𝑑r.$$
(3.19)
According to (3.16), it is obvious that the first two summands in (3.19) have a finite limit after dividing by $`m^2\kappa `$. In order to treat the last summand, we first multiply the $`T`$-equation (2.27) with $`m^2`$ and take the limits (3.4), (3.5), (3.10),
$$m^2(A1)\mathrm{\hspace{0.33em}2}r\phi ^{}.$$
Using again (3.10), (3.5), and $`\omega m`$, $`T1`$, we obtain that
$`{\displaystyle \frac{1}{m^2\kappa }}(\omega T\sqrt{A}m)`$ $`=`$ $`{\displaystyle \frac{1}{m^3\kappa }}m\omega (\sqrt{A}1)T+{\displaystyle \frac{1}{m^3\kappa }}m(\omega Tm)`$
$``$ $`r\phi ^{}+(\phi +E).`$
From this we conclude that the binding energy (3.19) divided by $`m^2\kappa `$ has a meaningful limit; more precisely
$$\widehat{B}:=\frac{B}{m^2\kappa }_0^{\mathrm{}}\left(\frac{1}{2q}\frac{(1w^2)^2}{r^2}+\frac{1}{q}w^2+\alpha ^2(E+\phi +r\phi ^{})\right)𝑑r.$$
(3.20)
We now describe our method for constructing numerical solutions of our reciprocal limit system. Since it is difficult to take into account the integral condition (3.12) in the numerics, we discard this condition for the construction of the solution; it will be taken care of later via a rescaling technique (see (3.34), (3.35)). This rescaling method requires only that the normalization integral be finite,
$$0<\lambda ^2:=_0^{\mathrm{}}\alpha (r)^2𝑑r<\mathrm{}.$$
(3.21)
According to (3.6), (3.7) and (3.13), (3.15), the behavior of the Dirac spinors at infinity is either oscillatory or exponential. As a consequence, the normalization integral in (3.21) will be finite only if $`\alpha (r)`$ tends to zero for $`r\mathrm{}`$. Furthermore, we can use the transformation (3.17) to set $`\phi (0)=0`$. Hence in the first construction step, we want to find solutions of the modified system
$`\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{w}{r}}\alpha \mathrm{\hspace{0.25em}2}\widehat{\beta }`$ (3.22)
$`\widehat{\beta }^{}`$ $`=`$ $`(E+\phi )\alpha {\displaystyle \frac{w}{r}}\widehat{\beta }`$ (3.23)
$`r^2w^{\prime \prime }`$ $`=`$ $`(1w^2)w+qr\alpha \widehat{\beta }`$ (3.24)
$`r^2\mathrm{\Delta }\phi (r)`$ $`=`$ $`\alpha ^2`$ (3.25)
with the following conditions at the origin,
$`w(r)`$ $`=`$ $`1\lambda r^2+𝒪(r^3),\alpha (r)=\alpha _1r+𝒪(r^2)`$ (3.26)
$`\phi (r)`$ $`=`$ $`𝒪(r^2),\beta (r)=𝒪(r^3),`$ (3.27)
together with the conditions at infinity
$`\underset{r\mathrm{}}{lim}w(r)`$ $`=`$ $`\pm 1,\underset{r\mathrm{}}{lim}\alpha (r)=\mathrm{\hspace{0.33em}0}`$ (3.28)
$`|\phi _{\mathrm{}}|:=|\underset{r\mathrm{}}{lim}\phi (r)|`$ $`<`$ $`\mathrm{}.`$ (3.29)
For any given value of the coupling constant $`q`$, we thus have two free parameters $`\lambda `$ and $`\alpha _1`$ to characterize the solutions near the origin $`r=0`$. Each solution has a unique extension to larger values of $`r`$. Asymptotically for $`r\mathrm{}`$, we must satisfy the two conditions (3.28). Thus we have as many free parameters as asymptotic conditions, and we therefore expect for fixed $`q`$ a discrete set of solutions satisfying (3.26), (3.27), and (3.28). For each solution, we must then verify that the conditions (3.21) and (3.29) are also satisfied.
For the construction of numerical solutions, we enhanced the technique used in to a two-parameter shooting method. Since two-parameter problems are considerably more difficult than one-parameter problems, we describe the method in some detail. For clarity, we first consider the simplified situation where $`\alpha (r)`$ and $`w(r)`$ have prescribed boundary values for a given finite $`r=r_1`$. In this case, one can apply the standard multi-parameter shooting method as e.g. described in . More precisely, one can for given initial data compute $`\alpha (r_1)`$ and $`(w(r_1),w^{}(r_1))`$ numerically, compare with the prescribed boundary conditions, and iteratively adjust the two free parameters $`\lambda `$ and $`\alpha _1`$ until the boundary conditions are satisfied to sufficient accuracy. In our case, the situation is more difficult because we have boundary values not for finite $`r=r_1`$, but for $`r=\mathrm{}`$. In order to deal with this problem, we first choose a finite $`r_1`$. Using an ansatz for the asymptotic form of the solution $`(\alpha ,\widehat{\beta },w,\phi )`$ at infinity, we approximately compute $`\alpha (r_1)`$ and $`(w(r_1),w^{}(r_1))`$ and derive conditions between these functions. Taking these conditions as the boundary conditions at $`r=r_1`$, we can apply the two-parameter shooting method on the finite interval $`(0,r_1]`$ as described above. The so-obtained solution on $`(0,r_1]`$ gives, in combination with the asymptotic formulas on $`(r_1,\mathrm{})`$, an approximate solution for all $`r>0`$. Since our asymptotic description becomes precise only in the limit $`r_1\mathrm{}`$, we must, in order to attain the desired accuracy, choose $`r_1`$ sufficiently large. In order to ensure that $`r_1`$ is appropriately increased during the computation, we modified the two-parameter shooting method in such a way that both the initial data and $`r_1`$ are adjusted in each iteration step. The iteration is stopped when the numerics has stabilized and the accuracy no longer improves. This modified shooting method was implemented in the Mathematica programming language using the standard ODE solver with a working precision of 16 digits. The initial data is adjusted in the iteration with a secant method, and the step size for incrementing $`r_1`$ is determined from the relative error of the numerical solution at the upper boundary $`r_1`$. After the iteration has been stopped and a numerical solution has been found, our program slightly changes the initial data and searches for a nearby solution. In this way, we can automatically trace a one-parameter family of solutions. Finally, we explain our method for describing the asymptotic behavior of the solutions at infinity. According to the asymptotics of the solutions of the ED and EYM equations , we can expect that the spinors $`\alpha `$ and $`\widehat{\beta }`$ will decay exponentially fast at infinity, whereas the potentials $`\phi (r)`$ and $`w(r)`$ for $`r\mathrm{}`$ will behave like rational functions. Therefore it is a reasonable asymptotic approximation to set $`\alpha `$ and $`\widehat{\beta }`$ to zero for $`r>r_1`$. In this approximation, the potential $`\phi `$ is harmonic according to (3.25). The YM equation (3.24), on the other hand, reduces to the vacuum YM equation (3.3). In the new variable $`u=\mathrm{log}r`$, this equation becomes autonomous; namely
$$_u^2w_uw=(1w^2)w.$$
(3.30)
This autonomous equation allows us to derive boundary conditions for $`w`$ as follows. We set
$$x=w\text{and}y=_uw.$$
(3.31)
Then the YM orbits in the $`(x,y)`$ plane are described by the following differential equation,
$$y^{}(x)=\frac{_u^2w}{_uw}\stackrel{(\text{3.30})}{=}\mathrm{\hspace{0.33em}1}\frac{(1x^2)x}{y}.$$
(3.32)
According to the boundary conditions (3.28) and the differential equation (3.30), the variables $`x`$ and $`y`$ must behave in the limit $`r\mathrm{}`$ like either $`x1`$, $`y0`$ or $`x1`$, $`y0`$. In both of these cases, there is a unique YM orbit $`y(x)`$, which can be easily calculated numerically by integrating (3.32). By transforming (3.31) back to the variable $`r`$, we obtain the following mixed boundary conditions for $`w(r)`$ at $`r=r_1`$,
$$w^{}(r_1)=\frac{1}{r_1}y(w(r_1)).$$
(3.33)
We next describe our rescaling method needed to arrange the normalization condition (3.12). Suppose that we have a solution of the modified system (3.22)–(3.29) with finite normalization integral, (3.21). A direct calculation shows that the transformed functions
$`\stackrel{~}{\alpha }(r)`$ $`=`$ $`\lambda ^2\alpha (\lambda ^2r),\stackrel{~}{\widehat{\beta }}(r)=\lambda ^4\beta (\lambda ^2r)`$ (3.34)
$`\stackrel{~}{\phi }(r)`$ $`=`$ $`\lambda ^4\left(\phi (\lambda ^2r)\phi _{\mathrm{}}\right),\stackrel{~}{w}(r)=w(\lambda ^2r)`$ (3.35)
solve our original reciprocal limit system (3.6), (3.7), (3.9), (3.11), and (3.12) with boundary conditions (3.13)–(3.15), if one replaces the energy $`E`$ and coupling constant $`q`$ by
$$\stackrel{~}{E}=\lambda ^4\left(E+\phi _{\mathrm{}}\right)\text{and}\stackrel{~}{q}=\lambda ^4q.$$
(3.36)
We point out that only the rescaled solutions (3.34),(3.35) and rescaled parameters (3.36) have a physical meaning. Therefore, we will in what follows consider only the rescaled tilde solutions; for ease in notation, the tilde will be omitted.
In the remainder of this section, we describe our numerical solutions of the reciprocal limit equations. Just as in the case for the ED and EYM equations , there are solutions for different rotation numbers of the spinors in the $`(\alpha ,\beta )`$-plane and for the YM potential in the $`(w,w^{})`$-plane. For simplicity, we restricted attention to solutions with rotation number zero for the spinors (as for the ground state solutions in ). For the YM potential, we consider only the cases where the $`(w,w^{})`$-curve either makes a half rotation joining the points $`(1,0)`$ and $`(1,0)`$, or makes a full rotation, ending at its starting point $`(1,0)`$. A typical example for a solution of each type is shown in Figures 3 and 3. Because of the similarity of the YM potential to the BM ground state and the BM first excited state, we refer to these two types in what follows as the ground states and the first excited states, respectively. Notice that the curves in the $`(w,w^{})`$-plane are not plotted all the way to their rest points at $`(1,0)`$ or $`(0,1)`$, respectively. The reason is that we plot only the numerical solution on the interval $`(0,r_1]`$. One sees that the spinors have become practically zero for $`r=r_1`$, and it is thus an admissible approximation to smoothly join the $`(w,w^{})`$-curve with a vacuum YM solution by using the boundary conditions (3.33). We first discuss the ground state solutions. In Figure 3, the main characteristics of the solutions are plotted versus the coupling constant $`q`$. As explained above, $`E`$ has the interpretation of the (appropriately scaled) energy of the Dirac particle. Since $`E`$ is negative, the Dirac particle has gained energy by forming the bound state. The parameter $`\widehat{B}`$, (3.20), gives the total binding energy, i.e. the amount of energy which is set free when the binding is broken up. Since $`\widehat{B}`$ is negative, we can expect that solutions of the full EDYM system, which are close to the solutions of the reciprocal coupling equations, should be stable. Finally, $`r_w`$ and $`r_\alpha `$ are the characteristic length scales of the solutions; more precisely, $`r_w`$ is the radius where $`w`$ changes sign, and $`r_\alpha `$ is the radius where $`\alpha `$ has its maximum,
$$w(r_w)=\mathrm{\hspace{0.33em}0}\text{and}\alpha ^{}(r_\alpha )=\mathrm{\hspace{0.33em}0}.$$
(3.37)
The characteristic radii are interesting because they give information about the “size” of the solutions as functions of $`r`$; i.e. they tell whether the fields are spread out in space, or whether they are localized close to the origin. It is also worth considering both radii because $`r_w`$ and $`r_\alpha `$ can behave quite differently (cf. Figure 3).
The plots in Figure 3 have a turning point at $`q8.49`$. Similar to the situation described for the spiral in , this is a bifurcation point which comes about as a consequence of our rescalings. One branch of solutions can be extended up to $`q11.6`$. For solutions close to this end point, the potential $`w`$ leaves the interval $`[1,1]`$ as shown in Figure 6. Since $`r_w`$ and $`r_\alpha `$ both go to zero in this limit, the spinors and YM field are both confined to a smaller and smaller neighborhood of the origin. At the same time, the energy of the Dirac particle and the binding energy become infinite. The other branch of solutions ends near $`q=8.95`$. For solutions near this end point, the $`(w,w^{})`$-curve comes very close to the origin before running into the rest point at $`(1,0)`$, see Figure 6. This makes the numerics rather delicate, and we therefore have not yet analyzed this regime in much detail. It is interesting that $`r_\alpha `$ is bounded near this end point, whereas $`r_w`$ seems to become infinite. This shows that, while the Dirac particle stays in a bounded region of space, the YM field becomes more and more spread out.
For the first excited state, the energy spectrum and characteristic radii are shown in Figure 6. Since in general $`w`$ never equals zero for the first excited state, we define $`r_w`$ via the minimum of $`w`$, i.e.
$$w^{}(r_w)=\mathrm{\hspace{0.33em}0}\text{and}\alpha ^{}(r_\alpha )=\mathrm{\hspace{0.33em}0}.$$
(3.38)
In contrast to the ground state, the solutions can now be extended up to $`q=0`$. In this regime, the YM potential stays close to $`w=1`$; see Figure 9. The solutions have a bifurcation point at $`q9.866`$. The branch coming out at the bifurcation point for larger values of $`E`$ is difficult to study numerically because the $`(w,w^{})`$-curve comes close to the origin, see Figure 9.
It is interesting that for the ground state solutions in Figure 3, the parameter $`q`$ stays bounded away from zero, whereas the plots for the first excited state in Figure 6 could be extended up to $`q=0`$. Let us consider how this can be understood directly from the equations. The parameter $`q`$ enters only the YM equation (3.9). In the limit $`q0`$, this equation goes over to the vacuum YM equation, which has only the trivial solution $`w1`$. Hence if we assume that the spinors have a finite limit for $`q0`$, then $`w(r)`$ must go uniformly in $`r`$ to one. This shows that the solutions can be regular for $`q0`$ only if $`w`$ satisfies the boundary condition $`lim_r\mathrm{}w(r)=+1`$. In particular, our ground state solutions cannot be regular in this limit. We next consider the limit $`q0`$ for the first excited state in more detail. Since $`w`$ converges uniformly in $`r`$ to one, the reciprocal limit equations (3.6), (3.7), and (3.11) go over to the Dirac-Newton equations
$`\alpha ^{}`$ $`=`$ $`{\displaystyle \frac{1}{r}}\alpha \mathrm{\hspace{0.25em}2}\widehat{\beta }`$ (3.39)
$`\widehat{\beta }^{}`$ $`=`$ $`(E+\phi )\alpha {\displaystyle \frac{1}{r}}\widehat{\beta }`$ (3.40)
$`r^2\mathrm{\Delta }\phi (r)`$ $`=`$ $`\alpha ^2.`$ (3.41)
These equations are obtained by taking the nonrelativistic limit of the ED equations , and according to the results obtained in that paper, the equations (3.39)–(3.41) together with the normalization integral (3.12) have a countable number of solutions, characterized by the rotation number of the spinors (called the ground state, the first excited state, etc.). We thus expect that the functions $`(\alpha ,\widehat{\beta },\phi )`$ corresponding to solutions of the reciprocal limit equations should for $`q0`$, go over to a solution of the Dirac-Newton equations. The behavior of the YM potential $`w`$ can now be analyzed in more detail by taking the solution $`(\alpha ,\widehat{\beta },\phi )`$ of the Dirac-Newton equations as a given inhomogeneity in the YM equation (3.9) and performing a perturbation calculation for small $`q`$. More precisely, the ansatz $`w(r)=1+qu(r)`$ to first order in $`q`$, leads to the linear equation
$$r^2u^{\prime \prime }=\mathrm{\hspace{0.33em}2}u+r\alpha \widehat{\beta },$$
which can be solved by integration. Fixing the integration constants with our boundary conditions $`u(0)=u(\mathrm{})=0`$ and $`u^{}(0)=0`$, we obtain the unique solution
$`u(r)=r^2{\displaystyle _0^r}{\displaystyle \frac{ds}{s^4}}{\displaystyle _0^s}t\alpha (t)\widehat{\beta }(t)𝑑t{\displaystyle \frac{r^2}{3}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{1}{t^2}}\alpha (t)\widehat{\beta }(t)𝑑t.`$ (3.42)
This consideration shows that for $`q0`$, the rotation number of $`w`$ is uniquely determined by the rotation number of the spinors. Furthermore, one sees that in the limit $`q0`$, the Dirac wave function is determined by the Dirac-Newton equations (3.39)–(3.41). Thus only the gravitational attraction is responsible for the formation of the bound state, whereas the YM field has no influence on the spinors.
## 4 Solutions of the EDYM Equations
In this section, we shall construct numerical solutions of the full EDYM equations and discuss their properties. Our method is to first find special solutions which are small perturbations of either the BM solutions or solutions to the reciprocal limit equations of the previous section. We then trace these solutions while gradually changing the coupling constants. This yields one-parameter families of solutions which can be extended even to regions in parameter space where the solutions are far from all of the known limiting cases.
In order to simplify the connection between the EDYM equations and the reciprocal limit equations of Section 3, it is useful to introduce a parameter $`\epsilon >0`$ in such a way that the reciprocal limit equations are obtained when $`\epsilon 0`$. To this end, we parametrize the EDYM equations in terms of the new variables $`(\epsilon ,q,E)`$ as follows,
$`\kappa `$ $`=`$ $`(\epsilon q)^{\frac{3}{2}},e^2=\sqrt{{\displaystyle \frac{q}{\epsilon }}}`$
$`m`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\epsilon q}}},\omega ={\displaystyle \frac{1}{\sqrt{\epsilon q}}}+\sqrt{\epsilon q}E.`$
Since the EDYM equations involve three dimensionless parameters (namely $`m^2\kappa `$, $`\omega /m`$, and $`e^2`$), introducing $`(\epsilon ,q,E)`$ is merely a transformation to new independent parameters, prescribing at the same time the gravitational constant (this means that we give up the freedom to rescale $`r`$ by fixing our length scale). In the limit $`\epsilon 0`$, both $`q`$ and $`E`$ go over to the corresponding parameters of the reciprocal limit system (see (3.8) and (3.5)). Also, it is easy to check that the limits (3.4), (3.10), and (3.16) are satisfied if we let $`\epsilon 0`$ and keep $`(q,E)`$ fixed. The parameters $`\epsilon `$ and $`q`$ can be written in dimensionless form as
$$\epsilon =\frac{m^2\kappa }{e^2},q=m^2\kappa e^2.$$
(4.1)
Thus $`\epsilon `$ describes the relative strength of gravity versus the YM interaction, whereas $`q`$ is the product of the gravitational and YM coupling constants. Up to a scale factor, $`E=\omega m`$. Since $`\omega `$ is the relativistic energy and $`m`$ the rest mass of the Dirac particle, $`E`$ can, exactly as in the previous section, be interpreted as the energy of the Dirac particle. Finally, we also describe the binding energy by a parameter which corresponds to our notation for the reciprocal limit system (3.20) and set
$$\widehat{B}=\frac{\rho m}{\sqrt{\epsilon q}}.$$
For the construction of numerical solutions, we use a two-parameter shooting algorithm combined with a rescaling method. Since this technique is quite similar to that described for the reciprocal limit equations in the previous section, we shall merely outline our procedure. In the first step of the construction, we consider the EDYM equations (2.23), (2.24) and (2.26)–(2.28) with the side conditions
$`0`$ $`<`$ $`\lambda ^2:={\displaystyle _0^{\mathrm{}}}(\alpha ^2+\beta ^2){\displaystyle \frac{T}{\sqrt{A}}}𝑑r<\mathrm{}`$ (4.2)
$`0`$ $`<`$ $`\tau :=\underset{r\mathrm{}}{lim}T(r)<\mathrm{}`$ (4.3)
$`\underset{r\mathrm{}}{lim}w(r)`$ $`=`$ $`\pm 1,`$ (4.4)
together with the following expansions near $`r=0`$,
$`\alpha (r)`$ $`=`$ $`\alpha _1r+𝒪(r^3),\beta (r)=𝒪(r^2)`$
$`A(r)`$ $`=`$ $`1+𝒪(r^2),T(r)=\mathrm{\hspace{0.33em}1}+𝒪(r^2).`$
For fixed $`\epsilon `$ and $`q`$, we thus have the two parameters $`\alpha _1`$ and $`E`$ to characterize a solution of this modified EDYM system near the origin $`r=0`$. On the other hand, we must satisfy two conditions at infinity; namely, $`w`$ must converge to $`\pm 1`$, (4.4), and the spinors must go asymptotically to zero in order for the normalization integral to be finite (4.2). Hence, we can apply a two-parameter shooting method as described in Section 3. In order to have optimal boundary conditions at the upper end point $`r=r_1`$, we again match this with the solution of the autonomous vacuum YM equation (see (3.32) and (3.33)). The shooting method was again implemented in Mathematica, using an accuracy of 32 digits. For each solution constructed in this way, we verify that (4.3) is satisfied and that the ADM mass is finite (2.32). Once we have found a solution of the modified equations, we rescale the solution according to
$`\stackrel{~}{\alpha }(r)`$ $`=`$ $`\sqrt{\tau }\lambda ^2\alpha (\lambda ^2r),\stackrel{~}{\beta }(r)=\sqrt{\tau }\lambda ^2\beta (\lambda ^2r)`$ (4.5)
$`\stackrel{~}{A}(r)`$ $`=`$ $`A(\lambda ^2r),\stackrel{~}{T}(r)=\tau ^1T(\lambda ^2r),`$ (4.6)
and transform the parameters $`(m,\omega ,\kappa ,e^2)`$ as follows,
$`\stackrel{~}{m}`$ $`=`$ $`\lambda ^2m,\stackrel{~}{\omega }=\tau \lambda ^2\omega `$ (4.7)
$`\stackrel{~}{\kappa }`$ $`=`$ $`\lambda ^6\kappa ,\stackrel{~}{e}^2=\lambda ^2e^2.`$ (4.8)
A straightforward calculation shows that the so-rescaled solution satisfies the EDYM equations (2.23), (2.24) and (2.26)–(2.28) together with the original side conditions (2.25) and (2.30)–(2.33). The parameters $`(\epsilon ,q,E)`$ transform under the rescalings as
$$\stackrel{~}{\epsilon }=\epsilon ,\stackrel{~}{q}=\lambda ^4q,\stackrel{~}{E}=\lambda ^4(E+(\tau 1)m\omega ).$$
(4.9)
In the limit $`\epsilon 0`$, these transformations coincide with the rescalings of the reciprocal limit equations (3.34)–(3.36). However, we remark that for the ED equations a much different rescaling is used. Namely, in order to get a better correspondence to the reciprocal limit equations, we here scale the gravitational constant $`\kappa `$, whereas in $`\kappa `$ is fixed to be 1 throughout. Clearly, only the rescaled solutions have physical significance. Therefore, in what follows we will consider only the rescaled solutions and again omit the tilde.
In a realistic physical situation, the gravitational coupling is very weak, whereas the YM coupling constant is of order one, i.e. $`m^2\kappa 1`$ and $`e^21`$. Hence, according to (4.1), we will here only investigate the parameter range $`\epsilon 1`$, and we are particularly interested in the situation for small $`q`$. In the limit $`\epsilon 0`$, there are known solutions of our EDYM system, namely the BM solutions (which more precisely are solutions in the limits $`\kappa 0`$ and $`e^2/\kappa 1`$), and the solutions of the reciprocal limit system constructed in Section 3. We take these special solutions as the starting point for the numerics. By varying the parameters $`\epsilon `$ and $`q`$ and tracing the solutions with our shooting and rescaling methods, we obtain a two-parameter family of solutions. In order to reduce the computational workload, we did not step systematically through the two-parameter space, but always kept one parameter fixed while varying the other parameter. Since $`\epsilon `$ remains unchanged under the rescaling (see (4.9)), it is most convenient to construct one-parameter families of solutions for different, fixed values of $`\epsilon `$.
We now describe the solutions we found. Exactly as for the reciprocal limit system in Section 3, we restricted attention to solutions with rotation number zero for the spinors and a rotation angle of $`\pi `$ or $`2\pi `$ for the YM potential. We again refer to these types of solutions as the “ground state” and the “first excited state,” respectively. For the ground state solutions, the energy spectrum and the characteristic radii are in Figures 9 and 10 plotted for different values of the parameter $`\epsilon `$ (the characteristic radii are again defined by (3.37)). The curves A for $`\epsilon =0`$ coincide with the plots for the reciprocal limit system in Figure 3. For small values of the parameter $`\epsilon `$, there are solutions for the EDYM equations which are close to the solutions of the reciprocal limit equations (compare the curves A and B in Figure 9). In this parameter regime, the EDYM solutions look typically as shown in Figure 11. The metric functions $`A`$ and $`T`$ are both close to one; thus the gravitational interaction is weak, in agreement with our considerations after (3.4). The spinors and the YM potential look very similar to the solution of the reciprocal limit equations in Figure 6. We conclude that the reciprocal limit system of Section 3 indeed describes a significant limiting case of the EDYM equations. However, one also sees that even for small $`\epsilon `$, not all the solutions of the EDYM equations are close to the reciprocal limit solutions. More precisely, curve B leaves the vicinity of curve A at $`q10`$ (see Figure 10). If one follows curve B after it branches off from curve A, the parameter $`q`$ first increases up to a turning point, and then decreases to $`q=0`$. If $`\epsilon `$ gets large, the solution curves no longer come so close to the reciprocal limit solutions (see curves C and D). The maximum of $`q`$ decreases (see curve C) and finally disappears (see curve D). Figure 12 shows a typical solution for small $`q`$. We note that in this parameter region, the metric functions $`A`$ and $`T`$ are not near one; this explains why the reciprocal limit equations are no longer a good approximation. Indeed, the potentials $`w`$, $`A`$, and $`T`$ now resemble a BM solution of the EYM equations , and the spinors look like the solution of the Dirac equation in the BM background. Hence $`q0`$ corresponds to the limit of weakly coupled spinors; i.e. spinors in a fixed BM background. Notice that the characteristic radii go to zero and the energies go to infinity in the limit $`q0`$ (see Figure 9). This can be understood from our rescalings. Namely, for the (unscaled) solutions of our modified EDYM system, the BM solutions are easily obtained by taking the limit $`\alpha _10`$ (in which the spinors go uniformly in $`r`$ to zero). In this limit, the normalization integral (4.2) tends to zero, and thus the rescalings (4.5)–(4.9) lead to a singular behavior of the rescaled solutions for $`q0`$. To summarize, there is a one-parameter family of solutions (obtained by continuously changing the coupling constants), connecting the BM solutions to our reciprocal limit solutions
We remark that our plots of the curves B have a small gap at $`q8.7`$. The reason is that in this region the numerics become unstable, and could not be carried out with our methods. But we were able to construct two branches of solutions which approach the problematic region from both sides. We suspect that the instability of the numerics is merely an artifact of our rescaling method, but it might well be an indication for a possible bifurcation point in this region. For the other curves C and D, we analyzed only the branch of solutions which extends towards smaller values of $`q`$.
For the first excited state, the energy spectrum and characteristic radii are plotted in Figures 13 and 14 (the characteristic radii are again defined by (3.38)). The curves A for $`\epsilon =0`$ correspond to the solutions of the reciprocal limit equations in Figure 6. In contrast to the situation for the ground state, the solutions for small $`\epsilon `$ are all close to the reciprocal limit solutions (compare the curves A and B). Figure 15 shows a typical solution for large $`q`$; one sees that the spinors and YM potential look similar to those in Figure 9. The form of the energy spectrum and the characteristic radii gradually change when $`\epsilon `$ is increased; for example, the cusp in the $`(q,r_w)`$-plot becomes smooth (see curve D). It is interesting that for $`q0`$, the curves converge independent of $`\epsilon `$ to a single limit point (see Figure 13). This limit point was already described at the end of Section 3 as the case when the spinors form a bound state due to their gravitational attraction, and the spinors generate a YM field (see (3.39)–(3.42)). This picture is in agreement with our numerics, since the spinors and metric functions, for a solution near this limit point, look similar to the ED solutions in the Newtonian limit, and $`w1`$ (see Figure 16). The fact that this limit point is independent of $`\epsilon `$ follows, because as explained at the end of Section 3, for $`q0`$, the YM equation decouples from the ED equations. For clarity, we point out that it would not be correct to say that the gravitational interaction dominates the YM interaction in the limit $`q0`$. Namely, according to (4.1), the ratio of the gravitational and YM coupling constants is kept fixed, and thus $`q0`$ corresponds to the limit where both coupling constants go to zero at the same rate. Nevertheless, the YM field has for $`q0`$ no influence on the energy spectrum and the characteristic radii.
A main qualitative difference between the ground state and the first excited state is that for the first excited state, we could not continuously join the solutions of the reciprocal limit equations with a BM solution. In order to see how this comes about, we did numerical calculations starting with a Dirac particle in the BM background (similar to that shown in Figure 17) and gradually increased the coupling of the spinors to gravity and to the YM field. For these “deformations of the first excited BM state,” the curves of the energy spectrum and the characteristic radii have spirals, whose size and shape drastically changes when $`\epsilon `$ is increased, see Figures 19 and 19. In the parameter regime where the energy plots spiral around, the spinors have self-intersections similar as observed for the ED solutions , see Figure 20.
We now discuss the stability of our solutions. The relevant parameter for the stability analysis is the total binding energy $`\widehat{B}`$. Namely, if $`\widehat{B}`$ is negative and smaller than the total energies of all other states, then energy is needed to break up the binding or to make a transition to any of the other states, and therefore for physical reasons the solution must be stable. Clearly, this energy argument does not provide a rigorous stability proof, and it also cannot replace the numerical analysis of linear stability (like e.g. in or \[3, Section 8\]), but it gives a strong indication for stability and is therefore commonly used (see e.g. ). Let us first apply this energy argument to the ground state solutions of Figures 9 and 10. One sees that the total energy becomes negative for large $`q`$. For the curves B and C, this region is plotted in more detail in Figure 21. For the solutions on branch b, the total binding energy is minimal, and thus this branch is stable. Applying Conley index methods with $`q`$ as the bifurcation parameter (see ), we obtain, as in , that the two other branches a and c are unstable. Indeed, the instability of branch c follows also from the continuity of the Conley index and the fact that in the limit $`q0`$, this branch goes over to the ground state BM solution which is known to be unstable . When $`\epsilon `$ is increased (see curve D, Figures 9 and 10), only one branch of solutions remains, which comprises the BM solutions as a limiting case and is therefore unstable. More precisely, the one-parameter family has in this case no bifurcation points, and in the limit $`q0`$ the solutions tend to an unstable BM solution. Thus using Conley index techniques, it follows that the entire one-parameter family is unstable. We conclude that for small $`\epsilon `$, there is a stable branch of ground state solutions for which $`q`$ lies in a finite interval away from $`q=0`$; all other ground state solutions are unstable. For the stability of the first excited state, we consider the plots of Figures 13 and 14. Since for $`q0`$, the spinors and metric functions go over to the Newtonian limit of the ED solutions, we conclude from that the branch of solutions starting at $`q=0`$ should be stable. This is in agreement with our above energy argument, because on this branch the total binding energy $`\widehat{B}`$ is negative, and is smaller than the total binding energy of the second branch of solutions, which comes out of the bifurcation point located at the maximum of $`q`$. Again, Conley index theory yields that this second branch is unstable. For the deformations of the first excited BM state, the total binding energy is positive (see Figures 19 and 20), and hence these solutions should be unstable. Indeed, for the branch of solutions which extends up to $`q=0`$ (i.e. before the first bifurcation point), this also follows from the continuity of the Conley index and the instability of the first excited BM solution.
| Felix Finster | | Joel Smoller | | Shing-Tung Yau |
| --- | --- | --- | --- | --- |
| Max Planck Institute MIS | | Mathematics Department | | Mathematics Department |
| Inselstr. 22-26 | | The University of Michigan | | Harvard University |
| 04103 Leipzig, Germany | | Ann Arbor, MI 48109, USA | | Cambridge, MA 02138, USA |
| Felix.Finster@mis.mpg.de | | smoller@umich.edu | | yau@math.harvard.edu | |
warning/0001/cond-mat0001433.html | ar5iv | text | # REFERENCES
## Abstract
The theory of fluctuation conductivity for an arbitrary impurity concentration including ultra-clean limit ($`T\tau \sqrt{\frac{T_c}{TT_c}}`$) is developed. It is demonstrated that the formal divergency of the fluctuation density of states contribution obtained previously for the clean case is removed by the correct treatment of the non-local ballistic electron scattering. We show that in the ultra-clean limit the density-of-states quantum corrections are canceled by the Maki-Thompson term and only classical paraconductivity remains.
Strong compensation of the quantum fluctuation corrections in clean superconductor
D.V.Livanov, G.Savona
Department of Theoretical Physics, Moscow State Institute of Steel and Alloys, Leninski pr. 4, Moscow 117936, Russia
A.A.Varlamov
Istituto Nazionale di Fisica della Materia(INFM)-Unità ”Tor Vergata”,
Dipartimento di Scienze e Tecnologie Fisiche ed Energetiche, Università di Roma ”Tor Vergata”, via di Tor Vergata, 00133 Roma, Italy
PACS: 74.40.+k, 74.50.+r, 74.20.-z
As it is well known, the first order fluctuation corrections to conductivity in the vicinity of superconducting transition are presented by the Aslamazov-Larkin (AL), Maki-Thompson (MT) and density of states (DOS) contributions. First one has the simple physical meaning of the direct charge transfer by the fluctuation pairs themselves and can be easily derived from the phenomenological time dependent Ginzburg-Landau equation . In this sense it is a result characteristic for the classical electron theory, while Maki-Thompson and DOS contributions have the purely quantum origin and can be calculated in the frameworks of the microscopic approach only .
The character of the electron scattering plays a very special role for the manifestation of fluctuation effects. In the BCS theory of superconducting alloys the only criterion of the metal purity exists: it is the ratio between the Cooper pair ”size” (zero temperature coherence length of pure metal $`\xi _0`$) and the electron mean free path $`\mathrm{}.`$ If the alloy is dilute ($`\mathrm{}\xi _0`$) the Cooper pairs motion is ballistic and impurities do not manifest themselves in superconductor properties. In the opposite case $`\mathrm{}\xi _0`$ the Cooper pairs motion has the diffusive character and the role of the effective Cooper pair size plays the renormalized coherence length $`\xi ^{^{}}=\sqrt{\mathrm{}\xi _0}.`$ The relative magnitude of fluctuation effects, which is determined by the Ginzburg-Levanyuk number, is proportional to $`(a/\xi )^n`$ ($`a`$ is an interatomic distance and $`n>0`$ depends on the effective dimensionality of the electron spectrum) and it grows for impure systems.
Dealing with the superconductor electrodynamics in fluctuation regime it is necessary to remember that in the vicinity of the critical temperature the role of fluctuation Cooper pair effective size plays the Ginzburg-Landau coherence length $`\xi _{GL}(T)=\xi _0/\sqrt{\epsilon }`$ (where the reduced temperature $`\epsilon =(TT_c)/T_c`$). So the case of dilute metal ($`\mathrm{}\xi _0`$) in the vicinity of the transition could be formally subdivided on clean, which is still local ($`\xi _0\mathrm{}\xi _{GL}(T)`$) and ultra-clean, non-local ($`\xi _{GL}(T)\mathrm{}`$) limits. In terms of the used in the theory of disordered alloys parameter $`T\tau `$ the same three domains can be written down as $`T\tau 1;`$ $`1T\tau 1/\sqrt{\epsilon }`$ and $`1/\sqrt{\epsilon }T\tau .`$ (We use units $`k_B=\mathrm{}=c=1`$). The latter case was rarely discussed in literature in spite of the fact that it becomes of the first importance already for metals of very modest purity, let us say, $`T\tau 10`$. Really, in this case the condition $`T\tau 1/\sqrt{\epsilon }`$ read for the reduced temperature as $`10^2\epsilon 1`$ practically covers all experimentally accessible range of temperatures for the fluctuation conductivity measurements. What concerns the usually considered local clean case ($`1T\tau 1/\sqrt{\epsilon }`$) for chosen value $`T\tau 10`$ it would not have any range of applicability: indeed, the equivalent condition for the allowed temperature interval $`\epsilon 1/(T\tau )^2`$ almost contradicts to the $`2D`$ thermodynamical Ginzburg-Levanyuk criterion of the mean field approximation applicability ($`Gi\frac{T_c}{E_F}\epsilon `$). Moreover, as it is well known, for transport coefficients the high order corrections become to be comparable with the mean field results much before than for thermodynamical ones, namely at $`\epsilon \sqrt{Gi}`$ . So in practice one can speak about the dirty, intermediate or ultra-clean cases but not about the clean one.
We will restrict our consideration by the most interesting case of 2D electron spectrum, relevant to the high temperature superconductors. As it is known, the classic 2D AL contribution turns out to be independent on the electron mean free path $`\mathrm{}`$ at all :
$`\delta \sigma _{AL}^{(2)}={\displaystyle \frac{e^2}{16}}{\displaystyle \frac{1}{\epsilon }}`$ (1)
Anomalous Maki-Thompson contribution, being induced by the pairing on the Brownian diffusive trajectories , naturally depends on $`T\tau ,`$ but its form changes only when $`\mathrm{}\xi _{GL}(T)`$ ($`T\tau 1/\sqrt{\epsilon }`$). Its calculation, even with the non-local Cooperon vertices but the standard propagator, in the ultra-clean limit leads to the expression less divergent in its temperature dependence but growing as $`T\tau \mathrm{ln}(T\tau )`$ with the increase of $`\mathrm{}`$ :
$$\sigma _{MT(an)}^{(2)}=\frac{e^2}{8}\{\begin{array}{c}\frac{1}{\epsilon \gamma _\phi }\mathrm{ln}(\epsilon /\gamma _\phi ),T_c\tau 1/\sqrt{\epsilon }\\ \frac{8\pi ^2T\tau }{\sqrt{14\zeta (3)}}\frac{1}{\sqrt{\epsilon }}\mathrm{ln}(T_c\tau \sqrt{\epsilon }),1/\sqrt{\epsilon }T_c\tau \end{array},$$
(2)
where $`\gamma _\varphi =\pi /8T\tau _\varphi `$ is the inelastic scattering rate.
The analogous problem takes place in the DOS and regular part of MT contributions: their standard diagrammatic technique calculations lead to the negative correction
$`\sigma _{DOS+MT(reg)}^{(2)}={\displaystyle \frac{e^2}{4s}}\kappa (T\tau )\mathrm{ln}\left({\displaystyle \frac{1}{\epsilon }}\right),`$ (3)
$$\kappa (T\tau )=\frac{\psi ^{^{}}(\frac{1}{2}+\frac{1}{4\pi \tau T})+\frac{1}{2\pi \tau T}\psi ^{^{\prime \prime }}(\frac{1}{2})}{\pi ^2[\psi (\frac{1}{2}+\frac{1}{4\pi \tau T})\psi (\frac{1}{2})\frac{1}{4\pi \tau T}\psi ^{^{}}(\frac{1}{2})]}=\{\begin{array}{c}\frac{56\zeta (3)}{\pi ^4}0.691,T\tau 1\\ \frac{8\pi ^2}{7\zeta (3)}(T\tau )^2,T\tau 1\end{array}$$
(4)
evidently divergent when $`T\tau \mathrm{}.`$
In the derivation of these results the local form of the fluctuation propagator and Cooperons (besides (2)) were used. It is why in view of the mentioned above peculiarity of ultra-clean limit, the extension of their validity for $`T\tau 1/\sqrt{\epsilon }\mathrm{}`$ seems to be doubtful.
One can notice that at the upper limit of the clean case, when $`T\tau 1/\sqrt{\epsilon }`$ both DOS and anomalous MT contributions turn out to be of the same order of value but of the opposite signs. So one can suspect that in the case of correct procedure of the impurities averaging in the ultra-clean case the large negative DOS contribution can be cancelled with the positive anomalous MT one.
The reexamination of all fluctuation corrections of the first order in the case of the arbitrary impurity concentration including non-local electron scattering in the ultra-clean superconductor will be the aim of this communication. The nontrivial cancellation of the contributions, previously divergent in $`T\tau ,`$ will be shown. It results in the reduction of the total fluctuation correction in ultra-clean case to the AL term only.
In purpose to calculate the Cooperon (impurity vertex) $`C(𝐪,ϵ_1,ϵ_2)`$ and fluctuation propagator $`L(𝐪,\omega _\mu )`$ (the two-particle Green function) in the general case of an arbitrary electron mean free path case one needs the explicit expression for the polarization operator $`𝒫(𝐪,ϵ_1,ϵ_2),`$ which due to the elasticity of scattering does not contain the frequency summation and for $`2D`$ spectrum has a form:
$`𝒫(𝐪,ϵ_1,ϵ_2)={\displaystyle \frac{d^2p}{(2\pi )^2}G(𝐩+𝐪,ϵ_1)G(𝐩,ϵ_2)}={\displaystyle \frac{2\pi N(0)\mathrm{\Theta }(ϵ_1ϵ_2)}{\sqrt{v_F^2q^2+(\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2)^2}}},`$ (5)
where $`\mathrm{\Theta }(x)`$ is the Heaviside theta-function, $`\stackrel{~}{ϵ}_n=ϵ_n+\frac{1}{2\tau }sgnϵ_n,`$ $`N(0)`$ and $`v_F`$ are the density of states and the velocity at the Fermi level. Let us stress that this result was carried out without any expansion over the Cooper pair center of mass momentum $`𝐪`$ and is valid for an arbitrary $`\mathrm{}q.`$
The standard ladder consideration results in the following expressions for the Cooperon and fluctuation propagator:
$$C^1(𝐪,ϵ_1,ϵ_2)=1\frac{\mathrm{\Theta }(ϵ_1ϵ_2)}{\tau \sqrt{v_F^2q^2+(\stackrel{~}{ϵ}_1\stackrel{~}{ϵ}_2)^2}}.$$
(6)
and
$`[N(0)L(𝐪,\mathrm{\Omega }_\mu )]^1`$ $`=`$ $`\mathrm{ln}{\displaystyle \frac{T}{T_c}}+{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\{{\displaystyle \frac{1}{n+1/2}}`$ (8)
$`{\displaystyle \frac{1}{\sqrt{\left(n+\frac{1}{2}+\frac{\mathrm{\Omega }_\mu }{4\pi T}+\frac{1}{4\pi T\tau }\right)^2+\frac{v_F^2𝐪^2}{16\pi ^2T^2}}\frac{1}{4\pi T\tau }}}\}`$
Near $`T_c`$ $`\mathrm{ln}\frac{T}{T_c}\epsilon `$ and for the local limit, when just small momenta $`\mathrm{}q1`$ are involved in the final integrations, the Eqs. (6) and (8) can be expanded over $`v_Fq/\mathrm{max}\{T,\tau ^1\}`$ and they are reduced to the well known local expressions.
The Feynman diagrams which contribute to conductivity in the first order of perturbation theory on electron-electron interaction in Cooper channel are presented in Fig.1. Let us start from the discussion of the Maki-Thompson contribution (diagram 6). We restrict our consideration by the vicinity of the critical temperature, where, for the most singular in reduced temperature contribution, static approximation is valid. It means that Cooper pair bosonic frequency can be assumed $`\mathrm{\Omega }_\mu =0`$.
Using the general expressions for Cooperons and propagator (6), (8) after the integration over electronic momentum, one can find:
$`Q^{(6)}\left(\omega _\upsilon \right)`$ $`=`$ $`4\pi N(0)v_F^2e^2T^2{\displaystyle \underset{\epsilon _n}{}}{\displaystyle }{\displaystyle \frac{d^2𝐪}{(2\pi )^2}}L(𝐪,0)\times `$ (10)
$`\left[(\underset{n}{\overset{}{ϵ}},\underset{n+\nu }{\overset{}{ϵ}},𝐪)+(\underset{n+\nu }{\overset{}{ϵ}},\underset{n}{\overset{}{ϵ}},𝐪)\right],`$
where
$$(\alpha ,\beta ,𝐪)=\frac{R_q(2\alpha )R_q(\alpha +\beta )\mathrm{\Theta }(\alpha \beta )R_q(2\alpha )R_q(2\beta )}{(\beta \alpha )^2\left(R_q(2\alpha )\frac{1}{\tau }\right)\left(R_q(2\beta )\frac{1}{\tau }\right)R_q(\alpha +\beta )},$$
(11)
$`R_q(x)=\sqrt{x^2+v_F^2𝐪^2}`$.
The analogous consideration of the main in the clean case DOS diagrams 2 and 4 leads to:
$`Q^{(2+4)}\left(\omega _\upsilon \right)`$ $`=`$ $`4\pi N(0)v_F^2e^2T^2{\displaystyle \underset{\epsilon _n}{}}{\displaystyle }{\displaystyle \frac{d^2𝐪}{(2\pi )^2}}L(𝐪,0)\times `$ (13)
$`\left[𝒟(\underset{n}{\overset{}{ϵ}},\underset{n+\nu }{\overset{}{ϵ}},𝐪)+𝒟(\underset{n+\nu }{\overset{}{ϵ}},\underset{n}{\overset{}{ϵ}},𝐪)\right],`$
with
$`𝒟(\alpha ,\beta ,𝐪)`$ $`=`$ $`(\beta \alpha )^2(R_q(2\alpha ){\displaystyle \frac{1}{\tau }})^2\times `$
$`\left[{\displaystyle \frac{R_q^2(2\alpha )+2\alpha (\alpha \beta )}{R_q(2\alpha )}}{\displaystyle \frac{\mathrm{\Theta }\left(\alpha \beta \right)R_q^2(2\alpha )}{\left(R_q(\alpha +\beta )\frac{1}{\tau }\right)}}\right]`$
One can see that each of expressions for $`Q^{(6)}\left(\omega _\upsilon \right)`$ and $`Q^{(2+4)}\left(\omega _\upsilon \right),`$ in accordance with , in the limit $`T\tau \mathrm{}`$ presents itself the Loran series of the type $`C_2(T\tau )^2+C_1(T\tau )+C_0+C_1(T\tau )^1+\mathrm{}`$ . The careful expansion of the sum of expressions (10) and (13) in the series of such type leads to the exact cancellation of all divergent contributions and even to the coefficient $`C_0=0`$. In result, the leading order of the sum of MT and DOS contributions in the limit of $`T\tau 1`$ turns out to be $`C_1(T\tau )^1`$ and it disappears in the non-local limit. The results of numerical calculation of $`Q^{(6)}\left(\omega _\upsilon \right)+`$ $`Q^{(2+4)}\left(\omega _\upsilon \right)`$ as the function of $`T\tau `$ according to Eqs. (10) and (13) are presented in Fig. 2 for different temperatures. One can convince himself in the rapid decrease of this sum with the increase of $`T\tau `$.
The remained four diagrams among the first order fluctuation corrections to conductivity (see, for example, the Fig.9 in the review article ) are negligible in the vicinity of $`T_c`$. The similar consideration of the remaining two DOS-like diagrams (3) and (5) gives
$`Q^{(3+5)}(\omega _\upsilon )`$ $`=`$ $`4\pi N(0)v_F^2e^2T^2{\displaystyle \underset{\epsilon _n}{}}{\displaystyle }d^2𝐪L(𝐪,0)\times `$ (15)
$`\left[𝒦(\underset{n}{\overset{}{ϵ}},\underset{n+\nu }{\overset{}{ϵ}},𝐪)+𝒦(\underset{n+\nu }{\overset{}{ϵ}},\underset{n}{\overset{}{ϵ}},𝐪)\right],`$
where
$$𝒦(\alpha ,\beta ,𝐪)=\frac{2\beta \mathrm{\Theta }(\alpha \beta )}{(\beta \alpha )\left(R_q(2\alpha )\frac{1}{\tau }\right)^2R_q(2\alpha )}$$
(16)
Evaluation of Eq.(15) demonstrates that for $`T\tau 1`$ the final contribution of diagrams 3 and 5 does not contain $`\tau `$dependence, and is less ($`\mathrm{ln}1/\epsilon `$) singular if compared with paraconductivity, in spite of the fact that each of diagrams 3 and 5 contains the divergent Loran term $`T\tau `$ which cancel each other.
So one can see that the DOS term divergence $`(T\tau )^2,`$ found before for clean case , has a restricted validity and can not be extended to $`T\tau \mathrm{}`$. In the limit of defectless superconductor the total DOS+MT contribution is proportional to $`\mathrm{ln}1/\epsilon `$ and is independent on $`T\tau .`$
Finally let us turn to the discussion of the AL contribution. In this case, as it is well known, even in the vicinity of $`T_c`$ we cannot restrict ourselves by the static approximation and analytical continuation over the external frequency has to be accomplished. The diagram 1 at Fig.1 represents the Aslamazov-Larkin contribution:
$$Q^{AL}(\omega _\upsilon )=\frac{e^2}{2\pi i}\frac{d^2𝐪}{(2\pi )^2}𝑑z\mathrm{coth}\frac{z}{2T}B^2(𝐪,z,\omega _\upsilon )L(𝐪,iz)L(𝐪,iz+\omega _\upsilon ),$$
(17)
where the three Green function blocks have to be calculated with the non-local Cooperons and the expression for the non-local propagators have to be used. This program is hard to be realized in the general form and at present has been tried to be solved with different approximations in the set of papers . We are interested here to study the effect of non-locality on the AL contribution in the first hand sacrificing the $`ac`$ effects ($`\omega 0`$) and being in the vicinity of the transition $`\epsilon 1.`$. So we omit the $`z,\omega _\nu `$ dependence of the Green functions block
$$B(𝐪,z=0,\omega _\upsilon =0)=4\pi N(0)T\upsilon _F^2q\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{\left(\sqrt{4\stackrel{~}{ϵ_n}^2+\upsilon _F^2𝐪^2}\frac{1}{\tau }\right)^2\left(\sqrt{4\stackrel{~}{ϵ_n}^2+\upsilon _F^2𝐪^2}\right)}$$
(18)
and evaluate the AL contribution in this approximation numerically. These calculations show that the temperature dependence of paraconductivity turns out to be close to the classical 2D $`\epsilon ^1`$regime. It is necessary to mention that relatively far from the transition, where our approximation strictly speaking is already not applicable, the calculated curve lies somewhat lower of the AL theory prediction in accordance with the short wave-length fluctuation calculations and experimental findings . We note, that although Eqs. (17), (18) contain dependence on $`\tau `$, the paraconductivity is turned out to be $`\tau `$-independent in the entire range of parameter $`T\tau `$, analogous to the local 2D result (1).
Let us discuss the results obtained. First of all it is necessary to stress the observed strong cancellation of the DOS and MT contributions, which were found previously, within the local fluctuation theory, to be divergent in the limit $`\tau \mathrm{}`$ . As we have demonstrated, the correct account of non-local scattering processes in the ultra-clean limit results in the impurity independent, logarithmic in reduced temperature, contribution negligible in comparison with the more singular AL one.
Let us remind that the fluctuation conductivity in the limiting case $`\tau =\mathrm{}`$ and for the non-zero frequency of the external electromagnetic field was studied in Ref. , where the similar problem of the $`\omega ^2`$-divergence( instead of $`\tau ^2`$ in our case) of the contributions from different DOS and MT-like diagrams aroused. The sum of all relevant diagrams nevertheless was found to be regular and proportional to $`\omega .`$ Moreover, the sum of all DOS and MT diagrams (which correspond to diagrams 6, 2 and 4 in Fig. 1) was shown to be zero in the case $`\tau =\mathrm{}.`$ In the current publication we have confirmed this statement studying the more general case $`\omega 0,\tau \mathrm{}`$ with $`\omega \tau 0`$ and convincing ourselves that for the one-electron type DOS and MT fluctuation processes the final result does not depend on the order of the $`lim_{\omega 0,\tau \mathrm{}}(\sigma ^{DOS}+\sigma ^{MT})`$ calculation. We have evaluated the explicit dependence of the overall fluctuation conductivity on the parameter $`T\tau `$ and have demonstrated its regular character.
What concerns the AL contribution, the careful investigation of its clean limit was done in by means of the analysis of the paraconductivity diagram structure in the coordinate representation. It was shown that the electric field does not interact directly with the fluctuation Cooper pairs, but it produces the effect on the quasiparticles forming these pairs only. The characteristic time of the change of the quasiparticle state is of the order of $`\tau `$ . Consequently the one-particle Drude type conductivity in $`ac`$ field, as it is well known, has the first order pole. What concerns the AL paraconductivity, due to the above mentioned reasons, the pole in it was found to be of the second order :
$$\sigma _{AL}(\omega )=\frac{\sigma _{AL}^{dirty}}{(1i\omega \tau )^2}$$
(19)
In spite of this difference one can see that the AL conductivity, like, the Drude one, vanishes at $`\omega 0,\tau \mathrm{}`$ because in the absence of impurities the interaction of electrons does not produce any effective force acting on the superconducting fluctuations, while the $`dc`$ paraconductivity conserves its usual $`\tau `$independent form.
In the present paper we have approached to the same problem of the investigation of the AL contribution in clean metal studying the general non-local case in q-space and have shown the independence of the $`dc`$ paraconductivity on the material purity.
It is necessary to stress that the non-local forms of the Cooperon and fluctuation propagator have to be accounted not only for the ultra-clean case but in every problem where the relatively large bosonic momenta are involved: account for the dynamical and short wavelength fluctuations beyond the vicinity of critical temperature, the effect of relatively strong magnetic fields on fluctuations and weak localization corrections etc. Recently such approach was developed in the set of studies of the DOS fluctuation effects and the efforts to apply it to the complete microscopic calculation of the magnetoconductivity for arbitrary temperatures and fields is undertaken in .
Authors are grateful to J.Axnas and C.Castellani for valuable discussions. This work was accomplished in the frameworks of the INTAS Grant # 96-0452.
Figure Captions
Fig.1. Feynman diagrams for the leading-order contributions to fluctuation conductivity. Wavy lines are fluctuation propagators, thin solid lines with arrows are impurity-averaged normal-state Green’s functions, shaded semicircles are vertex corrections arising from impurities, dashed lines with central crosses are additional impurity renormalizations and shaded rectangles are impurity ladders. Diagram 1 is the Aslamazov-Larkin term; diagrams 2-5 arise from the corrections to the normal state density of states in the presence of impurity scattering; diagram 6 is the Maki-Thompson term.
Fig. 2. The illustration of the decrease of the sum of DOS and MT contributions with the increase of the mean free path for the different values of reduced temperature: $`\epsilon =0.001;0.01;0.1;0.2.`$ |
warning/0001/gr-qc0001014.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The averaging problem is one of the fundamental problems in cosmology that we have not yet understood sufficiently so far.
It can be summarized as follows: In cosmology we want to understand the whole picture of our universe. However since the structure of the universe is so complicated that we can understand it only with the help of some format of recognition, viz. cosmological models. Thus, cosmology in principle requires a mapping procedure from reality to a model. Alternatively, one may regard this procedure as the procedure of averaging the real, complicated geometry in some manner in order to assign a simplified model-geometry to it. Now the problem is that Einstein equation is highly nonlinear, so that the effective dynamics of the averaged spatial geometry is expected to be highly complicated. Moreover, the effective dynamics of the averaged spatial geometry would not in general match the dynamics of the assigned model. Hence we should first analyze and understand the averaging procedure itself, otherwise we would make serious mistakes in finding out the evolution of the universe and/or matter content of the universe. If we state symbolically, the averaging procedure and the Einstein equation do not commute with each other.
Once we start investigating the averaging procedure itself, we immediately face with the trouble that we do not have a suitable ‘language’ for describing it. In order to formulate the approximation of the real geometry by a certain model, the concepts like ‘closeness’ or ‘distance’ between spaces are indispensable. However there has been no established mathematical theory so far which deals with these concepts and which can be applied to spacetime physics<sup>3</sup><sup>3</sup>3 As a concept which one should recall in this context, there is the Gromov-Hausdorff distance $`d_{GH}(X,Y)`$ between two compact metric spaces . Though it plays a central role in the convergence theory of Riemannian geometry, its abstract nature may be a big obstacle for its effective application to spacetime physics.. Here we will study the averaging problem in cosmology as the problem of defining a distance between spaces.
For this purpose we would like to focus on the spectral representation of spatial geometry as a promising attempt in this direction. The basic idea of the spectral representation is simple: We utilize the ‘sound’ of a space to characterize the geometrical structures of the space. This idea immediately reminds us of a famous problem in mathematics, ‘Can one hear the shape of a drum?’ In imitation of this phrase, we can state that ‘Let us hear the shape of the universe!’
## 2 The spectral distance
We should materialize the idea of ‘hearing the shape of the universe’ in a definite form. For definiteness we confine ourselves to $`(D1)`$-dimensional Riemannian manifolds that are spatial (metric signature $`(+,\mathrm{},+)`$) and compact without boundaries. Let $`Riem`$ denote this class of Riemannian manifolds. Now we set up the eigenvalue problem of the Laplace-Beltrami operator, $`\mathrm{\Delta }f=\lambda f`$. Then we obtain the spectra, viz. the set of eigenvalues (numbered in increasing order), $`\{\lambda _n\}_{n=0}^{\mathrm{}}`$. We note that, on dimensional grounds, the lower (higher) spectrum corresponds to the larger (smaller) scale behavior of geometry.
Suppose we want to compare two geometries $`𝒢`$ and $`𝒢^{}`$. Our strategy is hence to compare the spectra $`\{\lambda _n\}_{n=1}^N`$ for $`𝒢`$ with the spectra $`\{\lambda _n^{}\}_{n=1}^N`$ for $`𝒢^{}`$. However, taking a difference $`\lambda _n^{}\lambda _n`$ simply is not appropriate for our purpose: The simple difference $`\lambda _n^{}\lambda _n`$ would in general count the difference in the higher spectra (corresponding to the smaller scale behavior of geometry) with more weight. In spacetime physics, however, the difference in the larger scale behavior of geometry is of more importance than the one in the smaller scale behavior of geometry. (This is the precise description of the ‘spacetime foam picture’ and the scale-dependent topology .) In addition, the difference $`\lambda _n^{}\lambda _n`$ has a physical dimension $`[\mathrm{Length}^2]`$, which is not very comfortable either. Hence, we should rather take the ratio $`\frac{\lambda _n^{}}{\lambda _n}`$; Then the difference $`\delta \lambda _n:=\lambda _n^{}\lambda _n`$ in the lower spectrum is counted with more weight as $`\frac{\lambda _n^{}}{\lambda _n}=1+\frac{\delta \lambda _n}{\lambda _n}`$.
Now a measure of closeness $`d_N(𝒢,𝒢^{})`$ between two geometries $`𝒢`$ and $`𝒢^{}`$ can be introduced by comparing the spectra $`\{\lambda _n\}_{n=1}^N`$ (for $`𝒢`$) with $`\{\lambda _n^{}\}_{n=1}^N`$ (for $`𝒢^{}`$) as
$$d_N(𝒢,𝒢^{})=\underset{n=1}{\overset{N}{}}\left(\lambda _n^{}/\lambda _n\right).$$
(1)
Here the zero mode $`\lambda _0=\lambda _0^{}=0`$ is not included in the summation, and $`N`$ is the cut-off number which can be treated as a running parameter. The function $`(x)`$ ($`x>0`$) is a suitably chosen function which satisfies $`0`$, $`(1)=0`$, $`(1/x)=(x)`$, and $`(y)>(x)`$ if $`y>x1`$. We also note that the cut-off number $`N`$ characterizes up to which scale two geometries $`𝒢`$ and $`𝒢^{}`$ are compared. In this way, $`d_N(𝒢,𝒢^{})`$ is suitable for the scale-dependent description of the geometry.
At this stage, some comments may be appropriate on the spectral representation in general. It is true that the spectra can be explicitly calculated only for restricted cases. However, still there are several advantages for the spectral representation. First, the concept of the spectra itself is very clear. This is important for practical applications in physics. Second, even when the exact spectra themselves are not known explicitly, the perturbation analysis gives us important information on the spectra. For instance, one can investigate the perturbed spectra around some well-understood spectra, just like one investigates the perturbed metric around the Minkowski metric. Third, in spacetime physics, the lower spectra are more important than the higher spectra, since the former spectra reflect the large scale structure of the universe. Thus, even a few lower-lying spectra, which are in general easier to compute than the higher spectra, carry important information. (For more details, see Ref. .)
We note that the property $`\left(\lambda _n^{}/\lambda _n\right)0`$ as $`n\mathrm{}`$ is required for the convergence of $`d_N`$ as $`N\mathrm{}`$. Thus, it follows that $`\lambda _n^{}/\lambda _n1`$ as $`n\mathrm{}`$ should hold for convergence. Combined with the Weyl’s asymptotic formula , it means that , in the $`N\mathrm{}`$ limit, the dimension and the volume of $`𝒢`$ and $`𝒢^{}`$ should be same in order to give a finite $`d_N(𝒢,𝒢^{})`$ as $`N\mathrm{}`$ . When $`N`$ is kept finite as most of the cases we consider, these conditions need not necessarily to be satisfied for a finite $`d_N`$.
In order to utilize the measure $`d_N(𝒢,𝒢^{})`$ efficiently, it is desirable that $`d_N(𝒢,𝒢^{})`$ satisfies the axioms of distance, or at least some modified version of them:
Positivity: $`d_N(𝒢,𝒢^{})0`$, and $`d_N(𝒢,𝒢^{})=0`$ $``$ $`𝒢𝒢^{}`$, where $``$ means equivalent in the sense of isospectral manifolds ,
Symmetry: $`d_N(𝒢,𝒢^{})=d_N(𝒢^{},𝒢)`$,
Triangle Inequality: $`d_N(𝒢,𝒢^{})+d_N(𝒢^{},𝒢^{\prime \prime })d_N(𝒢,𝒢^{\prime \prime })`$.
Among several possibilities for the choice of $`(x)`$, there is one very important choice:
$`_a(x)=\frac{1}{2}\mathrm{ln}\frac{1}{2}(\sqrt{x}+1/\sqrt{x})`$.
Then Eq.(1) becomes
$$d_N(𝒢,𝒢^{})=\frac{1}{2}\underset{n=1}{\overset{N}{}}\mathrm{ln}\frac{1}{2}\left(\sqrt{\frac{\lambda _n^{}}{\lambda _n}}+\sqrt{\frac{\lambda _n}{\lambda _n^{}}}\right).$$
(2)
It is notable that this form for $`d_N`$ can be related to the reduced density matrix element in quantum cosmology under some circumstances . Namely, a long (short) spectral distance $`d_N(𝒢,𝒢^{})`$ can be interpreted as a strong (weak) quantum decoherence between $`𝒢`$ and $`𝒢^{}`$ for some cases in quantum cosmology. This interpretation of $`d_N`$ gives one motivation for the choice of $`_a(x)`$.
The measure of closeness $`d_N`$ defined in Eq.(2) satisfies (I) and (II) of the distance axioms, but it does not satisfy the triangle inequality (III) . Significantly enough, however, the failure of the triangle inequality turns out to be only a mild one since a universal constant $`c(>0)`$ can be chosen such that $`d_N^{}(𝒢,𝒢^{}):=d_N(𝒢,𝒢^{})+c`$ recovers the triangle inequality . Here $`c`$ is universal in the sense that $`c`$ can be chosen independent of $`𝒢`$, $`𝒢^{}`$ and $`𝒢^{\prime \prime }`$ although it depends on $`N`$. This fact leads to a significant consequence below: $`d_N`$ and its modification $`\overline{d}_N`$ (see below) are closely related to each other, which helps us reveal the nice properties of $`d_N`$.
There is another important choice for $`(x)`$:
$`_b(x):=\frac{1}{2}\mathrm{ln}\mathrm{max}(\sqrt{x},1/\sqrt{x})`$.
This is a slight modification of $`_a`$. In this case, Eq.(1) becomes
$$\overline{d}_N(𝒢,𝒢^{})=\frac{1}{2}\underset{n=1}{\overset{N}{}}\mathrm{ln}\mathrm{max}(\sqrt{\frac{\lambda _n^{}}{\lambda _n}},\sqrt{\frac{\lambda _n}{\lambda _n^{}}}).$$
(3)
Note that $`\overline{d}_N`$ satisfies (I)-(III), so that it is a distance.
Now, each measure of closeness introduced above has its own advantage: $`d_N`$ in (a) has an analytically neat form and it can be related to the quantum decoherence between $`𝒢`$ and $`𝒢^{}`$ in the context of quantum cosmology; However, it does not satisfy the triangle inequality. On the other hand, the measure $`\overline{d}_N`$ in (b) is a distance in a rigorous sense, although its form is not very convenient for practical applications (it contains $`max`$). Quite surprisingly, it turns out that $`d_N`$ and $`\overline{d}_N`$ are deeply related to each other. To discuss this property, we introduce an $`r`$-ball centered at $`𝒢`$ defined by $`d_N`$ in Eq.(2),
$$B(𝒢,r;d_N):=\{𝒢^{}Riem/_{}|d_N(𝒢,𝒢^{})<r\}.$$
Here $``$ indicates the identification of isospectral manifolds. In the same manner, we also introduce an $`r`$-ball centered at $`𝒢`$ defined by $`\overline{d}_N`$, $`B(𝒢,r;\overline{d}_N)`$.
Now we can show that
Theorem 1
The set of balls $`\{B(𝒢,r;d_N)|𝒢Riem/_{},r>0\}`$ and the set of balls $`\{B(𝒢,r;\overline{d}_N)|𝒢Riem/_{},r>0\}`$ generate the same topology on $`Riem/_{}`$.
For the proof of Theorem 1, first we should show that the set of all balls $`\{B(𝒢,r;d_N)|𝒢Riem/_{},r>0\}`$ can actually define a topology (let us call it “$`d_N`$-topology”), viz. the set of all balls can be a basis of open sets. This property is far from trivial, because of the failure of the triangle inequality for $`d_N`$.<sup>4</sup><sup>4</sup>4 On the other hand, a similar set of balls defined by $`\overline{d}_N`$ can define a topology (let us call it $`\overline{d}_N`$-topology), since $`\overline{d}_N`$ is a distance. Next, we need to show that any ball defined by $`d_N`$ (resp. $`\overline{d}_N`$) is an open set in $`\overline{d}_N`$-topology (resp. $`d_N`$-topology) .
From Theorem 1, we immediately obtain
Corollary
The space $`𝒮_N^o:=(Riem,d_N)/_{}`$ is a metrizable space. The distance function for metrization is provided by $`\overline{d}_N`$.
Hence we can extend $`𝒮_N^o`$ to its completion<sup>5</sup><sup>5</sup>5 It is desirable to investigate the structure of $`𝒮_N`$ intensively as a purely mathematical object., $`𝒮_N`$. Due to Theorem 1 and its Corollary, it is justified to treat $`d_N`$ as a distance and to regard $`𝒮_N`$ as a metric space, provided that we resort to the distance function $`\overline{d}_N`$ whenever the triangle inequality is needed in the arguments.
We can also show that $`𝒮_N`$ has several other desirable properties :
Theorem 2
The space $`𝒮_N`$ is paracompact.
Corollary
There exists partition of unity subject to any open covering of $`𝒮_N`$.
Theorem 3
The space $`𝒮_N`$ is locally compact.
Due to this property, we can construct an integral over $`𝒮_N`$ , which is essential to consider, e.g., probability distributions over $`𝒮_N`$.
Corollary
If a sequence of continuous functions on $`𝒮_N`$, $`\{f_n\}_{n=1}^{\mathrm{}}`$, pointwise converges to a function $`f_{\mathrm{}}`$, then $`f_{\mathrm{}}`$ is continuos on a dense subset of $`𝒮_N`$.
Theorem 4
The space $`𝒮_N`$ satisfies the second countability axiom.
These properties of $`𝒮_N`$ suggest that the space $`𝒮_N`$ can serve as a basic arena for spacetime physics. From now on we call $`d_N`$ in Eq.(1) (the form of $`d_N`$ in Eq.(2) in particular) a spectral distance for brevity.
## 3 Model-fitting procedure in cosmology
Now let us come back to the averaging problem in cosmology. Regarding this problem, there are several underling issues as follows:
How to select out a time-slicing for a given spacetime (‘reality’), which in general would possess no symmetry.
Furthermore,
How to incorporate the spatial diffeomorphism invariance,
How to incorporate the scale-dependent aspects of the geometrical structures, and
How to incorporate the apparatus dependence of the observed information
to the averaging procedure of geometry.
Considering its several desirable properties, the spectral representation seems to serve as a suitable ‘language’ for formulating and analyzing these issues. In particular the space $`𝒮_N`$ introduced in the previous section provides an appropriate platform for these discussions.
As a demonstration, let us give a rough sketch of the mapping procedure from reality to a model in terms of the spectral representation.
We here consider how to assign a model spacetime to a given spacetime (‘reality’). We first fix notations.
$`(1^{})`$ We consider a portion of a spacetime $`(,g)`$ bounded by two non-intersecting spatial sections $`\mathrm{\Sigma }_0`$ and $`\mathrm{\Sigma }_1`$ of $`(,g)`$. Let us denote this portion as $`(,g)_{(0,1)}`$.
$`(2^{})`$ Let $`Slice_o(,g)_{(0,1)}`$ be a set of all possible time-slicings of $`(,g)_{(0,1)}`$.
$`(3^{})`$ Hence, a slice $`sSlice_o(,g)_{(0,1)}`$ can be identified with a parameterized set of spatial geometries $`\{(\mathrm{\Sigma },h(\beta ))\}_{0\beta 1}`$. Let $`(,g)_{(0,1),s}`$ denote this set for brevity.
$`(4^{})`$ Let $`\{models\}`$ be a set of model spacetimes, bounded by two non-intersecting spatial sections $`\mathrm{\Sigma }_0^{}`$ and $`\mathrm{\Sigma }_1^{}`$ with a particular time-slicing, $`(^{},g^{})_{(0,1),s^{}}`$. Here, we distinguish between the identical spacetimes $`(^{},g^{})`$ with different choices of two non-intersecting spatial sections ($`\mathrm{\Sigma }_0^{}`$ and $`\mathrm{\Sigma }_1^{}`$) and/or different choices of a time-slicing ($`s^{}`$).
It is notable that the metric-space structure of $`𝒮_N`$ induces the same structure on $`Slice_o(,g)_{(0,1)}`$ also: Let $`s_1:=\{(\mathrm{\Sigma },h_1(\beta ))\}_{0\beta 1}`$ and $`s_2:=\{(\mathrm{\Sigma },h_2(\beta ))\}_{0\beta 1}`$ are any elements in $`Slice_o(,g)_{(0,1)}`$. Then, we can define
$$D_N(s_1,s_2):=_0^1\left[d_N((\mathrm{\Sigma }_1,h_1(\beta )),(\mathrm{\Sigma }_2,h_2(\beta )))\right]𝑑\mu (\beta ),$$
(4)
where $`\mu (\beta )`$ is a positive-definite measure. Clearly $`Slice_o(,g)_{(0,1)}`$ with $`D_N`$ becomes a metrizable space reflecting the same property of $`𝒮_N`$. Thus, we can consider its completion, $`Slice(,g)_{(0,1)}`$.
Now we describe the procedure of assigning a model spacetime to reality.
The choice of time-slicing
Let us choose and fix one model spacetime with a particular time-slicing $`(^{},g^{})_{(0,1),s^{}}`$. We can select the most suitable time-slicing of $`(,g)_{(0,1)}`$ w.r.t. (with respect to) the model $`(^{},g^{})_{(0,1),s^{}}`$ as follows: For each parameter $`\beta `$ ($`0<\beta <1`$), the closeness between the slice $`(\mathrm{\Sigma },h(\beta ))`$ in $`(,g)_{(0,1),s}`$ and the slice $`(\mathrm{\Sigma }^{},h^{}(\beta ))`$ in $`(^{},g^{})_{(0,1),s^{}}`$ can be measured by the spectral distance $`d_{A,\mathrm{\Lambda }}((\mathrm{\Sigma },h(\beta )),(\mathrm{\Sigma }^{},h^{}(\beta )))`$. Here the suffixes $`A`$ and $`\mathrm{\Lambda }`$ indicate, respectively, the elliptic operator (we here consider the Laplacian for simplicity) and the cut-off number (viz. $`N`$ in the previous section) for defining the spectral distance. Physically, $`A`$ and $`\mathrm{\Lambda }`$ symbolize the observational apparatus and the cut-off scale, respectively.
Now we can select out the most suitable time-slicing of $`(,g)_{(0,1)}`$ w.r.t. the model $`(^{},g^{})_{(0,1),s^{}}`$ as
Select $`s_0Slice(,g)_{(0,1)}`$ s.t.
$`D_{A,\mathrm{\Lambda }}((,g)_{(0,1),s},(^{},g^{})_{(0,1),s^{}})`$
$`:={\displaystyle _0^1}\left[d_{A,\mathrm{\Lambda }}((\mathrm{\Sigma },h\left(\beta \right)),(\mathrm{\Sigma }^{},h^{}\left(\beta \right)))\right]𝑑\mu \left(\beta \right)`$
gives the minimum.
On account of the property that $`D_{A,\mathrm{\Lambda }}`$ is bounded from below along with the completeness of $`Slice(,g)_{(0,1)}`$, some time-slicing $`s_0`$ of $`(,g)_{(0,1)}`$ is selected out w.r.t. the model $`(^{},g^{})_{(0,1),s^{}}`$, $`A`$ (apparatus) and $`\mathrm{\Lambda }`$ (scale)<sup>6</sup><sup>6</sup>6 To be more precise, there can be more than one slicings that satisfy the condition. Furthermore, $`s_0`$ can be a limit point of $`Slice_o(,g)_{(0,1)}`$, viz. $`s_0Slice(,g)_{(0,1)}Slice_o(,g)_{(0,1)}`$. In such a case, one would judge that $`(^{},g^{})_{(0,1),s^{}}`$ is not an appropriate model for $`(,g)_{(0,1),s}`$. In any case, it is desirable to investigate the mathematical structure of $`Slice(,g)_{(0,1)}`$ in more detail. .
Assignment of a model to ‘reality’
We can continue the same procedure for every model spacetime $`\{models\}`$ to choose the best-fitted model $`(^{},g^{})_{(0,1),s^{}}`$ and, w.r.t. it, the time-slicing $`s_0^{}`$ of $`(,g)_{(0,1)}`$. Then one can regard $`(_{},g_{})_{(0,1),s^{}}`$ to be the cosmological counterpart of $`(,g)_{(0,1),s_0^{}}`$ w.r.t. $`(A,\mathrm{\Lambda })`$. In this way the spectral representation naturally leads us to the concept of apparatus- and scale-dependent effective evolution of the universe.
## 4 Example:(2+1)-dimensional flat spacetimes
As an illustration for the procedure in the previous section, let us consider a simple example. We choose as ‘reality’ the simplest (2+1)-dimensional flat spacetime with topology $`T^2\times 𝑹`$: We can construct such a spacetime from $`𝑹^3`$ by the identification in space, $`(x+m,y+n)(x,y)`$, where $`m,n𝒁`$. (Here, $`(x,y,t)`$ is the standard coordinates for $`𝑹^3`$.) We can imagine this spacetime as a static spacetime with a spatial section being a regular 2-torus (a torus constructed from a unit square by gluing the edges facing each other), if $`t=const`$ slicing is employed. Now, let $`(,g)_{(0,1)}`$ be a portion of the spacetime defined by $`0t1`$. Then $`Slice_o(,g)_{(0,1)}`$ denotes a set of all slices for the present $`(,g)_{(0,1)}`$, and $`Slice(,g)_{(0,1)}`$ is its completion.
As a set of model spacetimes, $`\{models\}`$, we take a set of all (2+1)-dimensional flat spacetimes of topology $`T^2\times 𝑹`$ with particular slices; For each model, a particular time-slicing is employed by which the line-element is represented as
$$ds^2=dt^2+h_{ab}d\xi ^ad\xi ^b,$$
where
$$h_{ab}=\frac{V}{\tau ^2}\left(\begin{array}{cc}1& \tau ^1\\ \tau ^1& |\tau |^2\end{array}\right).$$
Here $`(\tau ^1,\tau ^2)`$ are the Teichmüller parameters of a 2-torus, and $`\tau :=\tau ^1+i\tau ^2`$, $`\tau ^2>0`$; $`(\tau ^1,\tau ^2)`$ and $`V(>0)`$ are functions of $`t`$ only; The periodicity in the coordinates $`\xi ^1`$ and $`\xi ^2`$ with period 1 are understood. We note that $`(\tau ^1,\tau ^2)`$ represent the shape of a parallelogram<sup>7</sup><sup>7</sup>7 In the present parametrization, the coordinates of four vetices of the parallelogram $`OACB`$ are $`O=(0,0)`$, $`A=(\tau ^1/\sqrt{\tau ^2},\sqrt{\tau ^2})`$, $`B=(1/\sqrt{\tau ^2},0)`$ and $`C=(\frac{1+\tau ^1}{\sqrt{\tau ^2}},\sqrt{\tau ^2})`$ . which forms the 2-torus by the edge-gluing; $`V`$ represents the 2-volume of the 2-torus .
The functional forms for $`\tau ^1`$, $`\tau ^2`$ and $`V`$ are not arbitrary; The evolutions of $`\tau ^1`$, $`\tau ^2`$ and $`V`$ w.r.t. $`t`$ are determined by a simple constrained Hamiltonian system
$`\{(\tau ^1,p_1),(\tau ^2,p_2),(V,\sigma );H0\}`$. Thus, in this example, $`\{models\}`$ is parameterized by distinct initial conditions for the Hamiltonian system. In other words, 4 parameters are required in principle to characterize each model in $`\{models\}`$.
Now, take one model in $`\{models\}`$, and consider its portion characterized by $`0t1`$. This portion of the model corresponds to $`(^{},g^{})_{(0,1),s^{}}`$ in the previous section. We easily get the spectra for each time-slice $`\mathrm{\Sigma }^{}`$ of $`(^{},g^{})_{(0,1),s^{}}`$: The Laplacian in this case becomes $`\mathrm{\Delta }=h^{ab}\frac{}{\xi ^a}`$ $`\frac{}{\xi ^b}`$; The normalized eigenfunctions of the Laplacian are
$$f_{n_1n_2}(\xi ^1,\xi ^2)=\mathrm{exp}(i2\pi n_1\xi ^1)\mathrm{exp}(i2\pi n_2\xi ^2)$$
with the spectra
$`\lambda _{n_1n_2}^{}`$ $`=`$ $`{\displaystyle \frac{4\pi ^2}{V\tau ^2}}|n_2\tau n_1|^2`$ (5)
$`=`$ $`{\displaystyle \frac{4\pi ^2}{V\tau ^2}}(|\tau |^2n_1^22\tau ^1n_1n_2+n_2^2)(n_1,n_2𝐙).`$
On the other hand, the ‘reality’ $`(,g)_{(0,1)}`$ is identical with an element of $`\{models\}`$ when $`t=const`$ slicing $`s_c`$ is employed; viz. $`(,g)_{(0,1),s_c}`$ is identical with the model characterized by $`\tau ^10`$, $`\tau ^21`$ and $`V1`$. Then, for every spatial section $`\mathrm{\Sigma }`$ of $`(,g)_{(0,1),s_c}`$, the spectra become
$$\lambda _{n_1n_2}=4\pi ^2(n_1^2+n_2^2)(n_1,n_2𝐙).$$
(6)
We can measure the spectral distance with the help of Eqs.(5) and (6) $`d_N((\mathrm{\Sigma },h(t)),`$$`(\mathrm{\Sigma }^{},h^{}(t)))`$. It is obvious that $`D_N((,g)_{(0,1),s},`$$`(^{},g^{})_{(0,1),s^{}})`$ gives the absolute minimum, 0, only when the model $`(^{},g^{})_{(0,1),s^{}}`$ is the one characterized by $`\tau ^10`$, $`\tau ^21`$ and $`V1`$, and the slicing of the ‘reality’ is $`s=s_c`$.
## 5 Dynamics of spectra
We have established the spectral distance, which provides a basis for comparing the real spatial geometry with a model spatial geometry. We can now investigate the spectral distance between ‘reality’ and a model as a function of time, which serves as the quantitative analysis of the influence of the averaging procedure on the effective dynamics of the universe. Here we see the usefulness of the spectral distance: On one hand it has a nice mathematical properties, and on the other hand, it can be handled explicitly. Thus, we now need dynamical equations for the spectra.
We first prepare concise notations for specific integrals that appear frequently below. Let $`A()`$ and $`A_{ab}()`$ be any function and any symmetric tensor field, respectively, defined on a spatial section $`\mathrm{\Sigma }`$. Let $`\{f_n\}_{n=0}^{\mathrm{}}`$ be the eigenfunctions of the Laplacian. Then we define
$`A_{mn}:`$ $`=`$ $`{\displaystyle _\mathrm{\Sigma }}f_nA(x)f_m,A_n:=A_{nn},`$
$`A_{ab}_{mn}:`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\lambda _m\lambda _n}}}{\displaystyle _\mathrm{\Sigma }}^af_mA_{ab}(x)^bf_n.`$
In order to derive the spectral evolution equations, we first recall a basic result of the perturbation theory (“Fermi’s golden rule”)
$$\delta \lambda _n=\delta \mathrm{\Delta }_n.$$
(7)
Noting $`\mathrm{\Delta }f=\frac{1}{\sqrt{h}}(\sqrt{h}h^{ab}_bf),_a`$, it is straightforward to get
$$\delta \mathrm{\Delta }_n=\overline{\delta h}_{ab}_n\lambda _n+\frac{1}{2}h\delta h_n\lambda _n,$$
(8)
where $`h\delta h:=h^{ab}\delta h_{ab}`$ and $`\overline{\delta h}_{ab}:=\delta h_{ab}\frac{1}{2}h\delta hh_{ab}`$. Combining Eq.(8) with Eq.(7), simple manipulations lead to a formula
$$\delta \lambda _n=\delta h_{ab}_n\lambda _n+\frac{1}{4}\mathrm{\Delta }(h\delta h)_n.$$
(9)
Now we identify $`\delta h_{ab}`$ in Eq.(9) with the time-derivative of the spatial metric, $`\dot{h}_{ab}`$, w.r.t. the time-slicing: $`\delta h_{ab}`$ should be replaced by $`\dot{h}_{ab}=2NK_{ab}+2D_{(a}N_{b)}`$. Here $`N`$ and $`N_a`$ are the lapse function and the shift vector, respectively; $`K_{ab}`$ is the extrinsic curvature and $`K:=K_a^a`$. After some manipulations , we finally reach the basic formula for the spectral evolution,
$$\dot{\lambda }_n=2NK_{ab}_n\lambda _n+\frac{1}{2}\mathrm{\Delta }(NK)_n.$$
(10)
We note that the shift vector $`N_a`$ does not appear in the final result, Eq.(10). This result comes from the fact that the spectra are spatial diffeomorphism invariant quantities.
For simplicity, let us set $`N1`$. Then we get
$$\dot{\lambda }_n=\frac{2}{D1}K_n\lambda _n+\frac{D3}{2(D1)}\mathrm{\Delta }K_n2ϵ_{ab}_n\lambda _n,$$
(11)
where $`ϵ_{ab}:=K_{ab}\frac{1}{D1}Kh_{ab}`$, and $`D`$ is the spacetime dimension ($`D=4`$ in the ordinary case).
It is not our present aim to go into the detailed analysis on the dynamics of the spectra, which will be done elsewhere . Here we only discuss a simple example as a demonstration of the usefulness of the spectral scheme: Let us investigate the scale-dependence of the effective Hubble parameter $`H_{eff}`$.
Suppose the spacetime is close to the closed Friedman-Robertson-Walker universe. In this case, $`\mathrm{\Delta }K`$ and $`ϵ_{ab}`$ are regarded as small. (More precisely, $`\mathrm{\Delta }K_{mn}`$ and $`ϵ_{ab}_{mn}`$ are small compared to $`K_{mn}\sqrt{\lambda _m\lambda _n}`$ and $`K_{mn}`$, respectively ($`m,n=1,2,\mathrm{}`$).) Then Eq.(11) can be written as (we set $`D=4`$)
$$\dot{\lambda }_n=\frac{2}{3}K_n\left[1\frac{1}{4}\frac{\mathrm{\Delta }K_n}{K_n\lambda _n}+3\frac{ϵ_{ab}_n}{K_n}\right]\lambda _n,$$
(12)
hence
$$(H_{eff})_n=\frac{1}{3}K_n\left[1\frac{1}{4}\frac{\mathrm{\Delta }K_n}{K_n\lambda _n}+3\frac{ϵ_{ab}_n}{K_n}\right].$$
(13)
The last two terms in the bracket describe the influence of inhomogeneity and anisotropy on $`H_{eff}`$ at scale $`\mathrm{\Lambda }`$. Here we note once more that the spectral representation naturally describes the apparatus- and the scale-dependent picture of the universe (in the present example, $`(H_{eff})_n`$).
## 6 Spacetime picture from the viewpoint of the spectral representation
We have discussed the spectral representation of geometrical structures in connection with the averaging problem in cosmology. In particular we have introduced $`𝒮_N`$, the space of all spaces equipped with the spectral distance, and have shown that $`𝒮_N`$ possesses several desirable properties as a basic arena for spacetime physics.
We have sketched the model-fitting procedure in the framework of the spectral representation, and have also suggested how to analyze the dynamical aspects of the averaging procedure within this framework. These arguments imply that the spectral scheme seems to be suitable for the analysis of the averaging problem. In fact, it naturally describes the apparatus- and scale-dependent effective evolution of the universe.
Finally, let us briefly discuss how spaces look like from the viewpoint of the spectral representation.
One may imagine the whole of the geometrical information of a space as a collection of all spectra such as
$$Space=\underset{i}{}(𝒟_i,\left\{\lambda _n^{\left(i\right)}\right\}_{n=0}^{\mathrm{}},\left\{f_n^{\left(i\right)}\right\}_{n=0}^{\mathrm{}}),$$
where $`𝒟_i`$ denotes an elliptic operator and $`\{\lambda _n^{(i)}\}_{n=0}^{\mathrm{}}`$ and $`\{f_n^{(i)}\}_{n=0}^{\mathrm{}}`$ are its spectra and the eigenfunctions, respectively. The index $`i`$ runs over all possible elliptic operators. A single observation is related to a subclass of elliptic operators corresponding to the observational apparatus. Thus we get only a small portion of the whole geometrical information of the space by a single observation. Such incomplete information may not be enough to determine geometry uniquely. Only one has to do then is to conduct other kinds of observation corresponding to other types of elliptic operators in order to get further information on geometry. (This is the physical interpretation of ‘isospectral manifolds’, viz. non-isometric manifolds with the identical spectra of the Laplacian.) It is also tempting to regard the spectral information more fundamental than the concept of Riemannian manifolds. Further investigations are needed to judge to what extent such a viewpoint of spacetime geometry is valid. |
warning/0001/cond-mat0001282.html | ar5iv | text | # Critical dynamics of singlet excitations in a frustrated spin system
\[
## Abstract
We construct and analyze a frustrated quantum spin model with plaquette order, in which the low-energy dynamics is controlled by spin singlets. At a critical value of frustration the singlet spectrum becomes gapless, indicating a quantum transition to a phase with dimer order. The magnetic susceptibility has an activated form throughout the phase diagram, whereas the specific heat exhibits a rich structure and a power law dependence on temperature at the quantum critical point. This generic behavior can be relevant to quantum antiferromagnets on Kagomé and pyrochlore lattices where (almost) all low-energy excitations are spin singlets, as well as to the CaV<sub>4</sub>O<sub>9</sub> lattice and the strongly frustrated antiferromagnets Cu<sub>2</sub>Te<sub>2</sub>O<sub>5</sub>(Cl,Br)<sub>2</sub> and Li<sub>2</sub>VO(Si,Ge)O<sub>4</sub>.
\]
It was argued by Haldane and Chakravarty, Halperin and Nelson that a quantum phase transition between a Néel state and a magnetically disordered phase in a two-dimensional (2D) quantum antiferromagnet is described by the 2+1 dimensional non-linear $`O(3)`$ $`\sigma `$-model. More recently the problem was studied in detail using the $`1/N`$ expansion ($`N=3`$ is the number of components of the order parameter), numerical methods as well as the Brueckner method . In this description the lowest quasiparticle excitation is a spin triplet. There are cases when a low-energy singlet excitation also appears, however it turns out to be irrelevant to the critical dynamics because of its vanishing spectral weight at the quantum critical point.
The above picture based on S=1 excitations describes a wide class of quantum phase transitions . On the other hand in Cr-based S=3/2 Kagomé materials it was found that at low temperatures ($`T35`$ K) the specific heat $`CT^2`$ and is practically independent of magnetic field up to 12 T . If we try to understand these data in terms of quasiparticles, the temperature behavior indicates a linear dependence of the excitation energy on momentum, and the insensitivity to field suggests that the quasiparticle spin is zero. Finite cluster numerical simulations for the S=1/2 Kagomé model indeed show that the lowest excitations are spin singlets with very small or zero gap. Similar behavior with many singlet states inside the triplet gap was also found in the pyrochlore antiferromagnet and the CaV<sub>4</sub>O<sub>9</sub> lattice . In the Kagomé case the singlet quasiparticle picture is far from being certain. It is unclear how to explain in this scenario the rather large zero temperature magnetic susceptibility observed in the Kagomé materials . It is also possible that the material is a spin glass and hence a description in terms of quasiparticles would not be adequate. Nevertheless, it is very important and interesting to analyze the quasiparticle scenario. From this point of view the Kagomé magnets are critical (or close to critical) systems, and they certainly can not be described by the $`O(3)`$ $`\sigma `$-model. In the present work we consider a two-dimensional spin system which exhibits a non-trivial singlet dynamics leading to a quantum critical point controlled by singlet excitations.
Critical singlet dynamics can naturally appear in a 2D quantum spin model with a zero temperature quantum phase transition separating two singlet ground states, both of which are magnetically disordered but have different discrete lattice symmetries. We show that a spin model constructed as an array of weakly coupled frustrated plaquettes has only singlet low-energy degrees of freedom in a certain range of parameters. Many theoretical works have been devoted to the conventional order-disorder quantum phase transition in similar models, mainly in relation to the CaV<sub>4</sub>O<sub>9</sub> compound . Different types of singlet dimer states have also been discussed for the quantum dimer model, the Heisenberg model on a 3D pyrochlore lattice , and the generalized 2D $`J_1J_2`$ model . However, low-energy critical singlet dynamics for these systems has not been studied until now. Very recent numerical work has shown, in the context of the 2D $`J_1J_2`$ model, that a quantum transition governed by singlet dynamics takes place for strong frustration $`J_2J_1/2`$ . It was also recently argued that this regime is experimentally accessible in Cu<sub>2</sub>Te<sub>2</sub>O<sub>5</sub>(Cl,Br)<sub>2</sub> and Li<sub>2</sub>VO(Si,Ge)O<sub>4</sub> , and measurements of the excitation spectrum of these materials, such as Raman spectroscopy, are now actively being pursued . Let us mention that dimer states appear also in various 1D spin models (e.g. in the frustrated $`J_1J_2`$ Heisenberg chain). Their critical behavior is understood in terms of deconfined spin-1/2 excitations or spinons, which do not generally exist in higher dimensions. A 1D system with features similar to the generic critical behavior we analyze in this work is an effective 1D Kagomé lattice model with gapless singlet and gapped triplet excitations. However, unlike the behavior we focus on, this model does not exhibit a quantum phase transition and breaking of lattice symmetry. Thus our future discussion is relevant to models in $`D>1`$, and a 2D case is analyzed as a generic example.
Consider a cluster of four S=1/2 spins (plaquette, see Fig.1(a)):
$$\widehat{}_0=J_1\underset{ij}{}𝐒_i𝐒_j+J_2\underset{(ij)}{}𝐒_i𝐒_j,$$
(1)
where $`ij`$ and $`(ij)`$ denote, respectively, the side and diagonal bonds. We choose the diagonal coupling $`J_2`$ to be close to the side one: $`J_2J_1=J>0`$, and define $`\alpha =J_1J_2J`$. In this regime the spectrum of a plaquette is the following: two close singlet states which we denote by $`|s_A`$ and $`|s_B`$, with energy difference $`ϵ_Bϵ_A=2\alpha `$, three almost degenerate triplet states with excitation energy $`J`$ above the singlets, and one S=2 state with energy $`3J`$. The two singlet wave functions are expressed as: $`|s_A=\frac{1}{\sqrt{3}}\{[1,2][3,4]+[2,3][4,1]\}`$ and $`|s_B=\{[1,2][3,4][2,3][4,1]\}`$, where $`[i,j]`$ denotes a singlet formed by the nearest-neighbor spins $`i`$ and $`j`$ (Fig.1(a)). Both singlets $`|s_A`$ and $`|s_B`$ are invariant under a four-fold rotation of a plaquette, whereas the columnar dimer states, e.g. $`[1,2][3,4]`$, break this symmetry. Consider now two decoupled plaquettes which for $`\alpha =0`$ have a $`2\times 2`$ degenerate ground state. The degeneracy is lifted if a weak interaction between the plaquettes is switched on selecting one of the columnar dimer states. Such a degeneracy lifting is different from the effect of a non-zero $`\alpha `$ and competition between the two leads to a nontrivial singlet dynamics in an array of weakly coupled plaquettes.
To be specific consider a square array with antiferromagnetic couplings $`j_1`$ between nearest-neighbor spins and $`j_2`$ between next-nearest-neighbor (diagonal) spins from different plaquettes, see Fig.1(b). For $`j_1=J_1`$ and $`j_2=J_2`$ it would be equivalent to the translationally-invariant 2D $`J_1J_2`$ model. However now we are interested in the weak-coupling limit: $`j_1,j_2J_2J_1=J`$. The low-energy singlet sector of the Hilbert space has a natural pseudospin representation in terms of the states $`|=|s_A`$ and $`|=|s_B`$. The total Heisenberg spin Hamiltonian of the array of coupled plaquettes is mapped on the following pseudospin Hamiltonian to lowest order in the small parameters $`\alpha `$ and $`j_{1,2}/J`$:
$`\widehat{}`$ $`=`$ $`\mathrm{\Omega }{\displaystyle \underset{i,j}{}}\left[{\displaystyle \frac{1}{6}}S_i^zS_j^z+{\displaystyle \frac{1}{2}}S_i^xS_j^x+{\displaystyle \frac{e^{i𝐐(𝐢𝐣)}}{2\sqrt{3}}}(S_i^zS_j^x+S_i^xS_j^z)\right]`$ (3)
$`\mathrm{\Omega }h{\displaystyle \underset{i}{}}S_i^z.`$
The energy scale $`\mathrm{\Omega }`$ and the effective dimensionless “magnetic field” $`h`$ are obtained in second-order perturbation theory:
$$\mathrm{\Omega }=\frac{1}{2J}(j_1^2+j_2^2),h=\frac{1}{\mathrm{\Omega }}[2\alpha +(j_1^2+j_2^26j_1j_2)/6J].$$
(4)
We have also defined $`𝐐=(0,\pi )`$ and dropped a constant term. The first sum in Eq.(3) is over nearest-neighbor plaquettes since bonds connecting second neighbors contribute to a constant term only and thus do not change the singlet dynamics. We stress that (3) is an exact mapping of the original Heisenberg model in the low-energy sector (excitation energy $`\omega \mathrm{\Omega },\alpha J`$). The Hamiltonian $`\widehat{}`$ describes an anisotropic ferromagnet in an external field. The interaction between pseudospins on adjacent sites is purely Ising, which can be seen by rotating the spin axes. However the Ising axis is staggered in the $`x`$$`z`$ plane deviating by angle $`\pm \pi /6`$ from the $`x`$-direction for horizontal (vertical) pairs.
The form of the Hamiltonian Eq.(2) is unaffected by asymmetry in the coupling constants $`j_1`$ and $`j_2`$, which would only change the position of the reference zero-field level. In fact, the pseudospin Hamiltonian remains the same with parameters given by: $`\mathrm{\Omega }=j_1^2/2J`$ and $`h=[2\alpha +(j_1^23j_1j_2)/6J]/\mathrm{\Omega }`$, even if we switch off one of the diagonal $`j_2`$-bonds between nearest-neighbor plaquettes. In this case our model resembles the frustrated $`\frac{1}{5}`$-depleted square lattice of CaV<sub>4</sub>O<sub>9</sub> , which has next-nearest-neighbor couplings comparable or even exceeding the nearest-neighbor exchange. Thus our results apply to the singlet ground states of this magnet as well.
In zero “magnetic field” $`h=0`$ the ground state of $`\widehat{}`$ has broken Ising symmetry with all spins parallel or antiparallel to the $`\widehat{x}`$ direction. In the language of spin singlets breaking of Z<sub>2</sub> symmetry corresponds to spontaneous dimerization in one of the two columnar patterns. The rotational symmetry of the square lattice is obviously broken. The “magnetic field” tends to orient all spins along the $`\widehat{z}`$ axis, and depending on the sign of $`h`$ favors either $`|s_A`$ or $`|s_B`$. Hence there are two critical fields, a positive and a negative one, for transitions into states with a restored Z<sub>2</sub> symmetry. Since Eq.(2) is invariant under $`hh`$, we consider the case $`h>0`$ only.
To study in more detail the properties of the symmetric phase near the critical point we map the pseudospin Hamiltonian (3) onto a hard-core boson model. Positive and large $`h`$ corresponds to $`J_1>J_2`$. In this case we choose the bare ground state as $`|0=_i|s_A_i`$. A boson creation operator on site $`i`$ is defined by $`|s_B_i=b_i^{}|s_A_i`$. In terms of pseudospins we have $`S_i^z=\frac{1}{2}b_i^{}b_i`$, $`S_i^{}=b_i^{}`$. The resulting boson Hamiltonian is:
$`H_B=ϵ{\displaystyle \underset{i}{}}b_i^{}b_i+t{\displaystyle \underset{ij}{}}(b_i^{}b_j+b_i^{}b_j^{}+h.c.)+`$ (5)
$`g{\displaystyle \underset{ij}{}}\pm (b_i^{}b_j^{}b_i+b_i^{}b_j^{}b_j+h.c.)+V{\displaystyle \underset{ij}{}}b_i^{}b_j^{}b_jb_i.`$ (6)
In the $`g`$-term the sign $`+`$ corresponds to a horizontal link and the sign $``$ to a vertical one. The parameters are:
$$ϵ=\mathrm{\Omega }(\frac{1}{3}+h),t=\frac{\mathrm{\Omega }}{8},g=\frac{\mathrm{\Omega }}{4\sqrt{3}},V=\frac{\mathrm{\Omega }}{6}.$$
(7)
Considered in combination with the hard-core constraint $`(b_i^{})^2=0`$, the Hamiltonian (5) is equivalent to the pseudospin Hamiltonian (3) and hence it is an exact mapping of the original spin problem in its low-energy sector.
The bosonic form of the effective Hamiltonian is convenient for the analysis of the low density (disordered) phase. Diagonalization of Eq.(4) in the quadratic approximation gives the following excitation spectrum: $`\omega _𝐤=\sqrt{(ϵ+4t\gamma _𝐤)^2(4t\gamma _𝐤)^2}`$, where $`\gamma _𝐤=\frac{1}{2}(\mathrm{cos}k_x+\mathrm{cos}k_y)`$. We set the inter-plaquette lattice spacing to one. The excitation gap $`\mathrm{\Delta }=\omega _{𝐤=0}`$ vanishes at the critical point $`h_c=\frac{2}{3}`$. The other critical point is at $`h_c`$. For $`h>h_c`$ the model is in the A-type plaquette phase, for $`h<h_c`$ it is in the B-type plaquette phase, and in between $`h_c<h<h_c`$ there is a phase with a boson condensate at $`𝐤=0`$. The condensate is doubly degenerate because the boson field can have both signs.
The zero-point quantum fluctuations, which exist in the disordered (plaquette) phases, change slightly the critical field $`h_c`$. The simplest way to account for the correlation effects is to use the Brueckner technique developed in Ref. . The cubic and quartic interactions in (5) are treated in the one-loop approximation, and the hard-core constraint is enforced via an infinite repulsion term $`U_ib_i^{}b_i^{}b_ib_i`$, $`U\mathrm{}`$, added to the Hamiltonian (5). This interaction is taken into account in the Brueckner (low-density gas) approximation. The singlet gap obtained from the self-consistent solution of the resulting Dyson equations is plotted in Fig.2. At the critical point the small parameter which justifies the Brueckner approximation is the singlet density $`b_i^{}b_i0.04`$. The renormalized critical value $`h_c0.65`$ is only slightly below the result for non-interacting magnons. The reason is the strong compensation of the hard-core corrections by the one-loop diagrams arising from the cubic terms in Eq.(4). The dependence of the gap on $`h`$ near the transition point agrees with the critical index $`\nu =0.63`$, expected for the $`O(N=1)`$ $`\sigma `$ model and is quite different from the value $`\nu =0.5`$ for non-interacting magnons.
Next we analyze the singlet excitation spectrum in the ordered dimer phase, $`h<h_c`$. We return to the pseudospin representation (3) and use the analogy between the evolution of the dimer phase and the spin reorientation process in an external magnetic field. The Holstein-Primakoff transformation is applied to the pseudospins in the rotating coordinate frame, which to lowest order coincides with the hard-core boson representation. At $`h=0`$ all spins point along the $`\widehat{x}`$-axis. For finite $`h`$ the spins tilt towards the field direction at an angle $`\mathrm{sin}\theta =h/h_c`$. Keeping only quadratic terms we find the following “classical” spectrum of singlet excitations:
$$\omega _𝐤=\mathrm{\Omega }\sqrt{1\gamma _𝐤(1\frac{2}{3}\mathrm{cos}^2\theta )+\frac{1}{\sqrt{3}}\mathrm{sin}2\theta \gamma _{𝐤+𝐐}}.$$
(8)
In accordance with the broken discrete symmetry the singlet excitations have a finite gap $`\mathrm{\Delta }=\omega _{𝐤=0}=\mathrm{\Omega }\sqrt{2(1h^2/h_c^2)/3}`$, plotted in Fig.2. The gap vanishes at the critical point. Proceeding further with the spin analogy we have also calculated the “spin reduction” for the ferromagnet (3). This parameter shows how far is the real ground state wave-function from the approximate mean-field ansatz. At $`h=0`$ we find $`S0.498`$, meaning that the effect of quantum fluctuations is extremely small and the linear spin-wave theory is well justified. At the transition point $`h=h_c`$ the zero-point oscillations are larger $`S0.44`$ but still small enough to justify our approximations leading to Eq.(8).
The broken Z<sub>2</sub> symmetry in the ground state at $`|h|<h_c`$ leads to a finite temperature transition. The whole phase diagram of the system of coupled plaquettes in the $`h`$$`T`$ plane is shown schematically in Fig.3. The critical temperature can be estimated as $`T_c(h=0)\mathrm{\Omega }`$ ($`1.14\mathrm{\Omega }`$ for a pure Ising case). The phase transition between the ordered dimer and the plaquette states belongs to the 2D Ising universality class, and we therefore expect a logarithmic singularity in the specific heat: $`C\mathrm{ln}|TT_c(h)|`$. Below $`T_c`$ the excitations acquire a gap and $`C`$ goes exponentially at low temperatures. The same activated dependence holds also for the symmetric phases at $`|h|>h_c`$. However, when we approach the quantum critical points at $`h=\pm h_c`$ the gap becomes smaller and the low-$`T`$ behavior of the specific heat changes to the quantum critical law $`C_V=\gamma T^2`$. The universal prefactor $`\gamma `$ can be calculated using the Brueckner technique . The result is $`\gamma =\frac{12\zeta (3)}{5\pi c^2}0.92/c^2`$, per plaquette, where $`c=\mathrm{\Omega }/2`$ is the singlet excitation velocity at the critical point (we set $`k_B=\mathrm{}=1`$). This value coincides with the large-N mean-field result . Notice that, unlike $`C_V`$, the magnetic susceptibility is always activated $`\chi \text{exp}(\mathrm{\Delta }_t/T),T0`$, since it is governed by the triplet gap $`\mathrm{\Delta }_tJ`$.
The variation of $`C_V`$ as a function of temperature is schematically presented in Fig.4 for the different parts of the phase diagram. In addition to the singlet contribution, whose form was discussed above, we have also shown the peaks expected to arise from the higher energy triplet and quintiplet states. We remind the reader that $`\mathrm{\Omega }`$ and $`J`$ represent two distinctly different energy scales since according to our weak inter-plaquette coupling assumption $`\mathrm{\Omega }J`$. We have estimated that the singlets contribute about 25% to the entropy $`S=(C_V/T)𝑑T`$ (and 61% and 14% are taken by the S=1 and S=2 states, respectively). The possibility of rich behavior at low T has been debated for some time for the Kagomé antiferromagnet and a sharp structure is believed to exist due to low-energy singlets . Notice that in our model a very peculiar behavior with a logarithmic singularity sets in at $`h<h_c`$ which crosses over to the quantum critical regime at $`h=h_c`$. Such a sharp feature is generic for systems near a transition between two singlet ground states. We also note that near the quantum critical point $`hh_c`$, $`T_c`$ is expected to be small (and ultimately vanish at $`h=h_c`$). Consequently the logarithmic singularity at $`T=T_c\mathrm{\Omega }`$ and the singlet peak at $`T\mathrm{\Omega }`$ should be clearly separated (this has been assumed in Fig.4). Alternatively, for $`h1`$ one has $`T_c\mathrm{\Omega }`$ and therefore the two contributions should merge. A numerically reliable calculation of $`C_V`$ in the temperature region $`T<\mathrm{\Omega }`$ is a separate problem, however for higher temperatures (of order $`J`$) a quantitative description can be easily achieved, and the result is plotted in Fig.5 (for the specific value of $`\mathrm{\Omega }=0.1J`$). The contributions of the triplet and quintiplet states have merged into a broad peak around $`TJ/3`$, while the singlet peak is sharper and centered at $`T\mathrm{\Omega }/2`$ (in this region our calculation is expected to be qualitatively correct).
In conclusion, we have analyzed a quantum spin model with purely singlet low-energy dynamics. There is a quantum phase transition in the model separating a disordered plaquette phase and a columnar dimer phase. Even though the inter-plaquette interaction was assumed to be weak, we expect our results to hold also for stronger interactions as long as there are no other instabilities, and to be applicable, e.g. to CaV<sub>4</sub>O<sub>9</sub>. The broken symmetries in the singlet ground state of CaV<sub>4</sub>O<sub>9</sub> were considered previously in Ref. in the framework of the quantum dimer model, and the possibility of Ising transitions between spin-Peierls and other disordered phases was discussed . Furthermore, recent numerical studies of the square-lattice $`J_1J_2`$ model have indeed found a quantum transition near $`J_2J_1/2`$ of the type discussed in the present work, even though a formally small expansion parameter (justifying the separation of the singlet and triplet dynamics) is not available.
We have found that while the magnetic susceptibility is always activated, the specific heat behavior is very rich and changes substantially in the different regimes. The model was inspired by experimental and numerical data for Kagomé systems and also shares common structure with pyrochlore antiferromagnets, the CaV<sub>4</sub>O<sub>9</sub> lattice, and strongly frustrated square lattice antiferromagnets. The novel quantum critical behavior associated with singlet criticality (or proximity to such a critical point) discussed in this work can be relevant to a wide class of disordered quantum spin systems, and could be detected by measurements of the specific heat as well as the low-T excitation spectrum in the S=0 sector via Raman spectroscopy.
We are grateful to C. Lhuillier, F. Mila, A.P. Ramirez, R.R.P. Singh, S. Sachdev, W.H. Zheng, and P. Lemmens for stimulating discussions. This work was supported by NSF Grant DMR9357474 (V.N.K.), in part by NSF Grant PHY94-07194 (O.P.S.), and by the Swiss National Fund (M.E.Z.). |
warning/0001/astro-ph0001239.html | ar5iv | text | # Modeling the X-ray – UV Correlations in NGC 7469
## 1 Introduction
The study of the physics of Active Galactic Nuclei (AGN) involves length scales much too small to be resolved by current technology or technology of the foreseeable future. As a result, this study is conducted mainly through the theoretical interpretation of the spectral and temporal properties of these systems, much in the way that the study of spectroscopic binaries has been used to deduce the properties of the binary system members and the elements of their orbit. Thus, studies of the spectral energy distribution in AGN have revealed the ubiquitous presence of a broad quasithermal component in the optical – UV part of the spectrum, the so called Big Blue Bump (BBB), as well as soft and hard ($`10`$ keV) X-ray emission, which in cases of sufficiently bright objects, was found to extend to several hundred keV. At lower frequencies, AGN were found to emit roughly the same amount of luminosity in the IR and the far-IR part of the spectrum as in the higher energy bands. This rough equipartition of the AGN luminosity from the far-IR to the X-ray part of the electromagnetic spectrum is an interesting fact for which no apparent, compelling explanation is presently at hand.
It was proposed long ago (Shields 1978) that a feature such as the BBB would signify the presence of a geometrically thin, optically thick accretion disk. It is generally thought that these disks radiate away the locally dissipated accretion energy in black body form at the temperature required to deliver the necessary radiant flux. For a quasar of luminosity $`L10^{46}L_{46}`$ erg s<sup>-1</sup> , associated with a black hole of mass $`M=10^8M_8M_{}`$ the corresponding disk temperature is $`T10^5L_{46}^{1/4}M_8^{1/2}`$ in reasonable agreement with the observed excess flux at the relevant wavelengths. The successful, detailed fits of the BBB feature as described above (Malkan & Sargent 1982; Malkan 1983; Sun & Malkan 1989; Laor & Netzer 1989), has convincingly established the identification of this spectral component with this specific structure of accretion flow onto the black hole. Even though certain of its properties do not completely conform with this notion, (e.g. the absence of prominent Ly $`\alpha `$ edges, and the magnitude of polarization with wavelength (see Koratkar & Blaes 1999), it is generally considered that this spectral component does indicate the presence of a geometrically thin, optically thick disk in the vicinity of the AGN black hole.
The X-ray emission must originate in a tenuous, hot ($`kT10^810^9`$ K) plasma, and its spectrum has been modeled successfully by Comptonization of soft photons by the hot electrons. The source of the soft photons is usually not specified but it is (naturally) considered to be the UV emission of the BBB, while the hot electrons are considered to be located in a corona overlying the BBB thin disk. This corona is thought to be confined and powered either by magnetic loops threading this disk, much the way it is the case with the solar corona (Galeev, Rossner & Vaiana 1979), or to be part of an Advection Dominated Flow (Narayan & Yi 1994) which are hot on their own right. The precise arrangement of these two components is not well defined; however, the form of the AGN spectral luminosity distribution between the optical–UV (BBB) and the X-ray bands suggests that only a small ($`20\%`$) fraction of the soft BBB photons traverse the volume occupied by the overlying hot electrons. This fact then suggests that either: (a) the hot plasma is confined in a small region with extent of only a few Schwarzschild radii around the black hole (while the BBB emission comes from a much larger region) or (b) the X-ray emitting plasma consists of small patches which cover only partially the thin disk (the source of BBB soft photons).
The sites of the IR and far-IR emission are thought to be at much larger distances from the compact object than that of the optical - UV emission, a conclusion reached on the basis of the much longer variability time scales in these components (Edelson & Malkan 1987) compared to those associated with the UV and the X-ray emission in AGN. In the context of unified AGN models, it is considered that the IR and far-IR emission results from reprocessing of the ionizing continuum by a molecular torus, whose inclination to the observer’s line of sight is believed responsible for the apparent dichotomy of Seyfert galaxies in types I and II.
While the above inferences and associations of the various spectral AGN components with specific spacial structures of rather well defined properties and location appear reasonable and make sense in the broader context of AGN physics, one has to bear in mind that they have been deduced mainly on the basis of spectral fits and radiative transfer. In this respect, it is well known (but some times not appreciated) that radiative transfer provides, generally, information only about column densities and optical depths (see e.g. Kazanas, Hua & Titarchuk 1997). However, in order to probe the dynamics and geometry of AGN one needs information about physical lengths and densities; these latter quantities have to be obtained by independent means, usually from time variability. For example, the goal of AGN monitoring (Netzer & Peterson 1987) has been precisely this, namely the determination of the physical size of the AGN Broad Line Region (BLR) using the reverberation mapping technique. The results of this effort have shown that the size of specific AGN components (the BLR region in this case) can in fact be very different from prior estimates based on spectral considerations alone (the size of BLR was found to be off by a factor of $`10`$ and the cloud density by a factor of $`100`$; for a review see Netzer & Peterson 1997).
An additional product of the extensive AGN multiwavelength monitoring efforts - which were aimed primarily in the determination of the BLR size from the continuum- Broad Emission Line correlations - has been the measurement of lags in the cross correlation functions between the optical and the UV continua. These are important because both these bands are part of the BBB and are hence thought to be produced by the putative thin accretion disk responsible for the emission of this spectral component. These studies indicated that the optical and the UV continua vary with much higher synchrony than expected on the basis of simple accretion disk models for the BBB emission. In the simplest models the variations in these two components should be propagating from the low to the high frequencies (the sense of mass inflow) and should be of order of the viscous time scales at the appropriate radii. Much shorter time scales are those associated with the propagation of sound waves traveling on the surface of the disk (in this case the UV variation could preceed that of the optical). Assuming a gas temperature $`T10^410^5`$ K, typical of the values needed to account for the BBB spectral characteristics and a size $`R10^{14}`$ cm, the sound crossing time scales are of order $`\tau R/c_s10^{7.5}10^8`$ s, much longer than the $`1\mathrm{day}10^5`$ s measured (Collier et al. 1998) or upper limits (Krolik et al. 1991) to the lags between the optical and UV wavelengths.
For the above reasons it was conjectured (Krolik et al. 1991) that the correlated optical – UV variability maybe caused by the reprocessing of X-rays, emitted by the hot corona overlying the thin disk, since this process yields signals which propagate much faster (at the speed of light) and might thus account for both the spectral and the temporal properties of these systems. However, the absence (until recently) of well coordinated, simultaneous observations in the X-ray and the UV – optical bands left this conjecture supported only by circumstantial evidence. In fact, timing studies and modeling limited to the correlations between UV and optical bands gave results consistent with such a picture: Rokaki & Magnan (1992) analyzed the results of the NGC 5548 monitoring campaign assuming that the observed variability is due to the reprocessing of an unseen harder spectral component (EUV - X-rays) by an optically thick geometrically thin disk. They found that the correlations between the continuum light curves at $`\lambda `$$`\lambda `$1360, 1840, 2670 and 4870 Å are consistent with their assumption, provided that the X-ray source had a variability they themselves prescribed and it were located 15 $`R_S`$ above the plane of an accretion disk around a black hole of mass $`6\times 10^7`$ M. Nonetheless, the absence of simultaneous X-ray observations, left this entire effort at the level of “reasonable conjecture”.
The launch of RXTE and the simultaneous presence of IUE in orbit made such coordinated observations possible: The active nucleus NGC 7469 was observed simultaneously both in the optical (Collier et al. 1998), in the UV (Wanders et al. 1997), and in the X-ray (Nandra et al. 1998) bands over an interval of roughly thirty days with a sampling rate no smaller than once every other orbit (in the UV and X-rays). The results of this campaign were rather startling and to some extent disappointing: While both the X-ray and the UV band exhibited variability of similar amplitudes ($`50\%`$), there were no apparent, easily understood correlations between the variability of these two bands, at least none that could be attributed reasonably to reprocessing of the X-rays as the cause of the observed optical – UV variability. In the UV, Wanders et al. (1997), measured lags of 0.23, 0.32 and 0.28 days between the variations at $`\lambda `$1315 Å and $`\lambda `$$`\lambda `$1485, 1740 and 1825 Å respectively. With an error of 0.07 days determined through Monte Carlo simulations they were unable to decide conclucively whether these results represented variability in accordance with accretion disk models or variability due to contamination by a very broad delayed emission feature which becomes stronger toward the red part of the spectrum. The optical observations of Collier et al. (1998), in conjuction with those of the UV provided a much larger dynamic range in wavelength which allowed the determination of lags between the UV ($`\lambda `$1315 Å) and the optical ($`\lambda `$$`\lambda `$4815, 6962 Å) with greater confidence. The $`2`$ day relative lags between $`\lambda `$1315 and $`\lambda `$6962 Å measured can in fact be intepreted as due to reprocessing by an accretion disk, as generally considered; however, it is not apparent which part of the spectrum drives the observed variability, while there appears to be, in addition, an inconsistency by more than a factor of 10 between the observed an inferred luminosity of this specific model.
Motivated by these observations we have decided to take a closer look at this particular question through the detailed modeling of the specific situation thought to take place in the innermost regions of AGN. Our approach is straightforward and similar in spirit to the analysis of Rokaki & Magnan (1992) : In §2 we compute the response function of reprocessing X-rays, from a point source above an infinite plane, as a function of the wavelength of the reprocessed radiation $`\lambda `$, the lag time $`\tau `$ and the latitude angle of the observer with respect to the disk plane $`\theta `$. In §3 we fold this response function with the observed X-ray light curve to produce model UV, optical and IR light curves for different values of the black hole mass and inclination angle. We then compute the autocorrelation and cross correlation functions of the model UV, optical and infrared light curves with that of the input X-ray and compare them to those observed. We search the black hole mass $`M`$ \- latitude angle $`\theta `$ parameter space in search of combinations which would result in autocorrelation and cross correlation functions similar to those observed. Finally, in §4 the results are summarized and conclusions are drawn.
## 2 The Response Function
In order to make the problem of X-ray reprocessing from a hot corona into UV-optical radiation by an underlying cool, geometrically thin, optically thick accretion disk as tractable as possible we have made the following idealizations: We have assumed the source of X-rays to be point-like and located at a height $`R_X`$ above an infinite plane representing the accretion disk producing the UV - optical emission associated with the BBB. We believe that the above assumptions approximate adequately the situation under consideration in that the X-ray source does not cover completely the source of UV photons, as required by the spectral fits and discussed in the introduction. This assumption is furthermore confirmed a posteriori by the more rapid variability of the X-ray relative to the UV emission.
The geometric construction associated with the arrangement of the X-ray source and the accretion disk described above is given in Figure 1. We choose Cartesian coordinates $`x^{\prime \prime },y^{\prime \prime },z^{\prime \prime }`$ with the $`x^{\prime \prime },y^{\prime \prime }`$ axes on the plane of the disk, their origin $`O`$ on the compact object (black hole), and the X-ray source $`S`$ at a distance $`R_X`$ above the disk plane in the $`z^{\prime \prime }`$ direction. The observer’s line of sight lies on the $`y^{\prime \prime },z^{\prime \prime }`$ plane, along the $`z`$-direction, as shown in the figure, making an angle $`\theta `$ with the plane of the disk.
The response function $`\mathrm{\Psi }(\lambda ,\tau )`$ of the situation depicted in Figure 1 is obtained from the following considerations: The loci of a given constant temperature on the surface of the disk, due to the reprocessing of X-rays from the source $`S`$, are circles centered around the foot $`O`$ of the vertical from the X-ray source to the disk plane; on the other hand, the loci of constant delay $`\tau `$ between the X-ray source and the observer are paraboloids of revolution about the observer - X-ray source axis (the $`z^{}`$ axis) with the X-ray source as their focal point. The intersection of these paraboloids with the plane of the disk are generally ellipses; the response function $`\mathrm{\Psi }(\lambda ,\tau )`$ is precisely the (thermal) emission by the intersection of the ellipses of constant delay with the circles of constant temperature on the disk. Because it is assumed that the reprocessed radiation is emitted in black body form at a temperature determined by the local X-ray flux, the determination of the response function $`\mathrm{\Psi }(\lambda ,\tau )`$ reduces to computing the area of intersection of the constant delay ellipses with the circles of constant temperature.
To determine this quantity, we choose Cartesian coordinates $`x^{},y^{},z^{}`$ whose origin is located on the extremum of the paraboloids of revolution (which have the source $`S`$ as their focal point) of a given constant delay $`D=c\tau `$, with the $`z^{}`$-axis pointing to the observer, also located on the $`y^{},z^{}`$ plane. In these coordinates one can easily verify that the surfaces of constant delay $`D=c\tau `$ are paraboloids of revolution of the form
$$z^{}=\frac{1}{2D}(x^{\mathrm{\hspace{0.17em}2}}+y^{\mathrm{\hspace{0.17em}2}})$$
(1)
with the distance between the origin of the coordinates (at the extremum of the paraboloids) and the X-ray source, located at their focal point, being $`D/2`$.
We now consider a Cartesian coordinate system $`x,y,z`$ with its origin $`O`$ located on the compact source and with the $`y`$ and $`z`$ axes parallel to $`y^{},z^{}`$ as shown in Figure 1. Let $`\mathrm{\Delta }y`$ and $`\mathrm{\Delta }z`$ be the coordinates of the X-ray source position in this system. Then, the relation between the coordinates $`x^{},y^{},z^{}`$ and $`x,y,z`$ can be easily obatained from the geometric construction of Figure 1. These are
$$z^{}=z(\mathrm{\Delta }zD/2),y^{}=y+\mathrm{\Delta }y,x^{}=x$$
(2)
In terms of these coordinates the equation of paraboloids of constant delay (Eq. 1) reads
$$z^{}=z(\mathrm{\Delta }zD/2)=\frac{1}{2D}[x^2+(y+\mathrm{\Delta }y)^2]$$
(3)
One can now express the coordinates of the source’s position $`\mathrm{\Delta }y`$ and $`\mathrm{\Delta }z`$ in terms of the height of the source above the disk $`R_X`$ and the observer’s latitude $`\theta `$, which read:
$$\mathrm{\Delta }y=R_Xcos\theta ,\mathrm{\Delta }z=R_Xsin\theta $$
(4)
The intersection of the paraboloids of constant delay with the disk plane (i.e. the isodelay curves on the disk plane) are obtained by finding the intersection of the paraboloid (Eq. 3) with the disk plane $`z=ycot\theta `$. This leads to
$$ycot\theta =R_Xsin\theta \frac{D}{2}+\frac{1}{2D}[x^2+(y+R_Xcos\theta )^2],$$
(5)
which after some rearrangement reads
$$\left[y(Dcot\theta R_Xcos\theta )\right]^2+x^2=\frac{D}{sin\theta }\left(\frac{D}{sin\theta }2R_X\right)$$
(6)
This is the equation of a circle on the plane $`x,y`$ (the plane perpendicular to the observer’s line of sight to the X-ray source) of radius square $`D^2/sin^2\theta 2DR_X/sin\theta `$, which, along with the equation for the accretion disk plane $`z=ycot\theta `$, give the parametric equations of the sought curve. Clearly the radius of the circle is non-zero only for sufficiently long lags, i.e. for $`D>2R_Xsin\theta `$, that is for the time it takes the corresponding isodelay surface to “cut” the surface of the accretion disk. For $`D=2R_Xsin\theta `$, the coordinate of the center of this circle is at a distance $`y=R_Xcos\theta `$, i.e. symmetric about the origin of the $`x,y,z`$ system (i.e. the location of the black hole) and the intersection of the observer – X-ray source line on the $`x,y`$ plane. One can easily see that for the above value of the lag $`D`$, the isodelay surfaces are tangent to the plane $`z=ycot\theta `$ at the point $`x=0,y=R_Xcos\theta ,z=R_Xcos^2\theta /sin\theta `$.
In order to compute the equation of the intersection on the plane of the accretion disk, the above equation has to be expressed in terms of the coordinates $`x^{\prime \prime },y^{\prime \prime },z^{\prime \prime }`$ on the plane $`z=ycot\theta `$. These are related to $`x,y,z`$ by the relation $`x^{\prime \prime }=x`$, $`y^{\prime \prime \mathrm{\hspace{0.17em}2}}=y^2+z^2=y^2+y^2cot^2\theta =y^2/sin^2\theta `$, $`z^{\prime \prime }=0`$; i.e. one has to project this circle onto the plane of the accretion disk. This, as expected, yields the equation of an ellipse, namely
$$\left[y^{\prime \prime }sin\theta (Dcot\theta R_Xcos\theta )\right]^2+x^{\prime \prime 2}=\frac{D}{sin\theta }\left(\frac{D}{sin\theta }2R_X\right)$$
(7)
Dropping the primes and defining polar coordinates on this plane, $`x=Rsin\varphi `$, $`y=Rcos\varphi `$, one obtains easily, solving for $`D`$,
$$D=\sqrt{R^2+R_X^2}Rcos\theta cos\varphi +R_Xsin\theta .$$
(8)
This is the equation for the delays used by Rokaki & Magnan (1992), with the difference that their angle $`\theta `$ denotes the colattitude (inclination angle) rather than the lattitude of the observer with respect to the accretion disk, given in our notation.
One can now compute, by setting the value of the angle $`\varphi `$ in Eq. (8) to $`\varphi =0,\pi `$, the range of values of the delay $`D`$ for which a ring of given radius $`R`$ will be illuminated by an instantaneous flash of X-rays emitted by the source $`S`$ at $`t=0`$. These are
$$D_l=\sqrt{R^2+R_X^2}Rcos\theta +R_Xsin\theta ,D_t=\sqrt{R^2+R_X^2}+Rcos\theta +R_Xsin\theta $$
(9)
where the subscripts $`l`$ and $`t`$ refer to the leading and trailing times. Then, the entire interval of the illumination of a ring of radius $`R`$ on the plane of the accretion disk is $`|D_lD_t|=2Rcos\theta `$, while the position at which the isodelay surface of lag $`D`$ touches the plane of the accretion disk is given by the projection of the circle of Eq. (6) onto this plane, i.e. at the point of radius $`R=R_Xcos\theta /sin\theta =R_Xcot\theta `$.
Using Eq. (8) one can now compute the area of overlap between these two families of curves. The simplest way to do this is to form the cross product of the tangent vectors of these two families of curves at the points of their intersection $`(x_i,y_i)`$. These vectors are $`(dx_i/dD,dy_i/dD)`$ and $`(dx_i/dR,dy_i/dR)`$. After a considerable amount of algebra, performed with the use of Mathematica, the expression for the overlap area reduces to the following simple expression
$$A(D,R)=\frac{2R}{\sqrt{(DD_l)(D_tD)}}$$
(10)
with the values of $`D_l,D_t`$ given above.
One can now integrate the above area over all $`D`$ between $`D_l`$ and $`D_t`$. This integration yields $`2\pi R`$ indicating that the Area function as given above is properly normalized. Therefore, the emitting area as a function of the time lag $`\tau `$ is
$$A(\tau )=\{\begin{array}{cc}2R/\sqrt{(\tau \tau _l)(\tau _t\tau )}\hfill & \text{if }\tau _l<\tau <\tau _t\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}.$$
(11)
where $`\tau _l=D_l/c`$ and $`\tau _t=D_t/c`$.
In Figure 2 we present the ratio $`Y(\tau )=A(\tau )/R^3`$ of the overlap function $`A(\tau )`$ divided by the cube of the radius of the disk as a function of the delay $`\tau `$ in units of $`R/c`$, for $`\theta =\pi /3`$ and assuming $`R_X=1`$. This is a quantity of interest as it indicates the reprocessed flux, integrated over all frequencies, contributed at a given lag $`\tau `$ by the appropriate radii, $`R=1,2,3`$ in units of $`R_X`$. The figure makes apparent the rapid decrease of the reprocessed radiation with radius and also the increasing range of delays $`\tau `$ which contribute to the emission at a given radius with the increase of the radius. For the given geometry, most of the reprocessing is effected by the shortest radii $`RR_X`$. Also apparent in the figure are the (integrable) singularities associated with the instants of intersection of the constant temperature circles with the constant delay ellipses.
The response function $`F(\lambda ,\tau )`$ at a wavelength $`\lambda `$ and lag $`\tau `$ will be the flux contributed at a given wavelength $`\lambda `$ by all radii of the disk which emit at the given lag $`\tau `$. Assuming that the reprocessed radiation is emitted locally with a black body spectrum, the response function has the form
$$F(\nu ,\tau )=_{r_{\mathrm{min}}}^{r_{\mathrm{max}}}B_\nu [T(R)]A(\tau )𝑑R$$
(12)
where $`B_\nu [T(R)]`$ is the Planck function of temperature $`T(R)`$. The limits of integration $`r_{\mathrm{min}}`$ and $`r_{\mathrm{max}}`$ are the minimum and maximum radii contributing to a given lag and are computed using Eq. (8) by setting $`c\tau =D_l`$ and $`c\tau =D_t`$.
The temperature $`T(R)`$ is calculated assuming that a fraction $`1𝒜`$ ($`𝒜`$ $`1`$ is the albedo of the disk) of the X-ray flux incident at a given point on the disk is thermalized and re emitted in black body form, i.e.
$$\sigma T(R)^4=\frac{L_X(1𝒜)}{4\pi (R_X^2+R^2)}\frac{R_X}{(R_X^2+R^2)^{1/2}}$$
(13)
or
$$T(R)=\left[\frac{L_X(1𝒜)R_X}{4\pi \sigma }\right]^{1/4}\frac{1}{(R_X^2+R^2)^{3/8}}$$
(14)
If the source of X-rays is not point-like but extended, a different relation between the flux and the distance to the disk would result; Rokaki & Magnan (1992), for example, have used the expression for a spherical source of radius $`R_X`$. This would lead to a slightly different relation between $`T`$ and $`R`$ for $`RR_X`$. However, for large $`R`$ the above relation is essentially correct.
## 3 Reprocessing the X-ray Flux
Having obtained the expression for the response of the disk to X-ray illumination as a function of the frequency of the reprocessed radiation and the lag $`\tau `$, one can now proceed to the computation of the time dependence of the reprocessed radiation at a given wavelength $`\lambda `$. If $`S_X(t)`$ is the light curve of the observed X-ray radiation, then, the reprocessed emission at a wavelength $`\lambda `$ as a function of time, $`f_\lambda (t)`$, will be given by
$`f_\lambda (t)`$ $`=`$ $`{\displaystyle 𝑑\tau F(\lambda ,\tau )S_X(t\tau )}`$ (15)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑RB_\lambda [T(R)]{\displaystyle _{\tau _l}^{\tau _t}}𝑑\tau A(\tau )S_X(t\tau )`$ (16)
It is apparent from the discussion of the previous section that the quantity which sets the scale of the lag is the distance of the X-ray source from the disk, or alternatively, for extended sources the size of the X-ray emitting region, $`R_X`$. We shall assume in the rest of this note, in order to fix our units, that this distance is equal to the Schwarzschild radius of the black hole. Since the scale of the temperature of the reprocessed radiation for the given, observed luminosity of NGC 7469 at a given radius $`R`$ is scaled by the value of $`R_X`$ measured in cm, should $`R_X`$ be larger than the Schwarzschild radius by a given factor, this could simply be reabsorbed in the value of the mass of the black hole, which would have to be decreased by the same factor.
Using as input the observed X-ray flux as a function of time, we employed Eq. (16) to compute the corresponding variations of the reprocessed flux at a number of wavelengths $`\lambda `$, for a set of values for the black hole mass $`M`$ and for two values of the latitude angle of the observer with respect to the plane of the disk $`\theta `$. The (normalized) X-ray light curve in the 2-4 keV range we used is given in Figure 3. This light curve consists of 256 equally spaced points, each representing the average flux in an interval of 10,768 seconds, provided to us by P. Nandra. The flux was obtained by interpolation of the real light curve at the equally spaced intervals given in the figure. The details of these observations are given in Nandra et al. (1998). It is worth noting that in the first 10 days of the observing campaign the sampling interval was roughly half of that given in the figure (once every 90 min., or 5,400 s), while for the remaining of the campaign it was once every 180 minutes or 10,800 seconds. The values of the black hole mass used in our calculations were $`M=10^7,\mathrm{\hspace{0.17em}10}^8,\mathrm{\hspace{0.17em}10}^9\mathrm{M}_{}`$, while for the angle $`\theta `$ we used the values $`\theta =\pi /3,\pi /10`$.
Figure 4 exhibits a representative set of such light curves corresponding to $`M=10^7\mathrm{M}_{}`$ and $`\theta =\pi /3)`$ for four different wavelengths, i.e. $`\lambda `$$`\lambda `$1315Å, 6962Å, 15000Å ($`2\times 10^{14}`$ Hz) and 30000Å ($`10^{14}`$ Hz). The first two wavelengths were chosen to match those at which the UV and optical observations were made, while the last two in order to explore whether model light curves at different wavelengths could reproduce the general characteristics of the observed light curves. The general trend apparent in this figure is the “smoothing” of light curves with increasing wavelength, this “smoothing” becoming apparent already at $`\lambda `$6962 Å. Such a trend is not surprising, as the contribution to the flux at these longer wavelengths comes from increasingly larger radii. In addition to the smoothing of the sharpest features associated with the X-ray light curve, the RMS variability also decreases significantly with increasing wavelength from $`50\%`$ at $`\lambda `$1315 Å, to $`33\%`$ at $`\lambda `$6962 Å, to $`25\%`$ at $`\lambda `$15000 Å ($`2\times 10^{14}`$ Hz), since an increasing range in the lag contributes to the emission at a given wavelength.
Figure 5 exhibits the effect of the black hole mass $`M`$ (or the height of the location of the X-ray source above the disk) on the light curves of the reprocessed radiation, where the light curves at a specific wavelength ($`\lambda `$1315 Å) are shown for three different values of the black hole mass $`(M=10^7,\mathrm{\hspace{0.17em}10}^8,\mathrm{\hspace{0.17em}10}^9\mathrm{M}_{})`$. The most obvious feature in this figure is the increase of the flux with the increase of the black hole mass. This fact can be understood as follows: for a given constant reprocessed luminosity, an increase in the mass corresponds to a decrease in the corresponding disk effective temperature; since the 1315 Å wavelength is in the Rayleigh–Jeans part of the corresponding spectra, a smaller temperature requires, in this regime of the spectrum, a larger flux in order that the same luminosity be radiated away. Increase in the black hole mass leads also to a certain amount of smoothing and shifts in the corresponding light curves, which become most prominent at the largest values of this parameter (see also the figures of the autocorrelation and cross c correlation functions). These trends are expected, as a larger value for the black hole mass corresponds to a larger distance of the X-ray source from the reprocessing disk plane. In this respect, one should note that the increase of the mass $`M`$ from $`10^7`$ to $`10^8`$ M causes very little additional smoothing at $`\lambda `$1315 Å. The flux in this wavelength is emitted by a region of extent $`R_X`$, which, in both these cases, is smaller than the light crossing distance associated with the X-ray light curve sampling interval ($`10^4`$ s), i.e. $`R_X\stackrel{<}{}3\times 10^{14}`$ cm. The smoothing and shifting trends become more pronounced at longer wavelengths as it becomes apparent in Figure 6, where the reprocessed light curves are shown for two values of the wavelength ($`\lambda `$$`\lambda `$= 1315, 6962 Å) and two values of the black hole mass $`(M=10^7,10^8\mathrm{M}_{})`$.
The effects of the inclination angle on the resulting light curves are shown in figure 7 which exhibits these light curves at $`\lambda `$$`\lambda `$6962, 15000 Å for two different values of the observer’s latitude $`(\theta =\pi /3,\pi /10)`$. One can see that the effect of the latitude on the light curve does increase with increasing wavelength: It is barely discernible at $`\lambda `$6962 Å (it is even smaller at $`\lambda `$1315 Å), but it is easily observable (though still small) at $`\lambda `$15000 Å. Qualitatively, the larger inclination angle light curve $`(\theta =\pi /10)`$ preserves slightly more of the high frequency content of the X-ray light curve than does the lower inclination one ($`\theta =\pi /3`$) and it certainly has a faster “turn–on” phase. The reason for that is that in the former case part of the reprocessing takes place almost along the observer’s line of sight and thus does not get “washed-out” by time-of-travel across the source effects.
To cast the arguments given above concerning the “smoothing” and “shifting” of the model light curves with increasing black hole mass and wavelength in more quantitative form, we have computed the autocorrelation (ACF) functions of our model light curves (a sample of which is given in the figures above) as well as their cross correlation (CCF) functions with the input X-rays. We have searched the black hole mass $`M`$ and angle $`\theta `$ parameter space with the intent of comparing these ACFs and CCFs to those of NGC 7469 given in Nandra et al (1998). A slight technical detail should be mentioned here: the model light curves we produce using as input that of the X-rays have a “turn–on” phase; had this section of the light curve been included in the computation of the ACF and CCF, it would lead to unphysical forms for these functions. For instance, this “turn–on” feature would lead to an ACF for the $`\lambda `$1315 Å light curve which is narrower than that of the X-ray light curve, an unphysical result since the latter signal can at best track the former exactly. To avoid such unphysical effects, in computing the ACF and CCF, we have excluded this “turn–on” section from the corresponding light curves. For the computation of the CCFs in particular, we have excluded at the same time an interval of equal length from the X-ray light curve too.
The ACFs are given in figures 8 and 9 for two different values of the black hole mass ($`M=10^7,10^8\mathrm{M}_{})`$, $`\theta =\pi /10`$ and for the four different wavelengths $`\lambda `$$`\lambda `$ 1315Å, 6962Å, 15000Å ($`210^{14}`$ Hz), 30000Å ($`10^{14}`$ Hz) along with the ACF of the X-ray light curve itself depicted by the solid curve. These figures make more quantitative the statements made above concerning the “smoothing” of light curves with wavelength and black hole mass: The ACF of the $`\lambda `$1315 Å light curve is almost identical to that of the X-rays and it changes very little with increasing the black hole mass from $`10^7`$ to $`10^8`$ M. The ACFs become broader with increasing both the black hole mass and the wavelength of the corresponding emission, indicating progressively “smoother” light curves with changes in these parameters, in agreement with the above discussion.
The shift in time of the emission due to X-ray reprosessing relative to the incident X-rays themselves is given by their corresponding cross correlation functions (CCF). These are shown in Figures 10 and 11 respectively along with the observed value of the X-ray–UV CCF (Nandra et al. 1999), for the same values of the parameters used to produce the ACFs of Figures 8 and 9. The trends are similar to those of the ACFs: Increasing black hole masses and wavelengths lead to longer lags between the X-rays and the reprocessed radiation, again in agreement with the qualitative arguments made earlier. An exception to this rule is CCF of the light curve of $`\lambda `$30000 Å for $`M=10^7`$ M and $`\theta =\pi /10`$. The CCF these values of the parameters exhibits a behavior not unlike that corresponding to the CCF between the X-rays and the UV emission at $`\lambda `$1315 Å in the data of Nandra et al. (1998). We believe that this particular feature is related to the specific form of the X-ray light curve. As it is apparent it disappears for a larger value of the black hole mass (see Figure 11).
Finally, in order to exhibit our results in a manner similar to that of Collier et al. (1998), who discussed in detail the UV – optical continuum lags of NGC 7469, and cast them in the parlance of accretion disk physics, we have produced the cross correlation of our model light curves at $`\lambda `$6962 Å with those at $`\lambda `$1315 Å. We have done that in two ways: (a)First we assumed, as done so far, that the variability in both these bands is driven by the X-rays. The resulting cross correlations functions are given in Figure 12 for three values of the black hole mass $`M=10^7,\mathrm{\hspace{0.17em}10}^8,\mathrm{\hspace{0.17em}10}^9\mathrm{M}_{}`$ and $`\theta =\pi /10`$. As expected, since either light curve tracks very closely that of the X-rays, their cross correlation peaks at time scales associated with the X-ray variability. The observed lags are therefore too long to be accounted by the geometric arrangement. (b) In the spirit of the Rokaki & Magnan (1992) treatement (and at the referee’s suggestion) we have assumed that the UV - optical variability is driven by an unseen spectral component other than the X-rays (EUV ?), with a light curve identical to that of the observed UV emission. We then computed the resulting light curves at $`\lambda `$1315 Å and $`\lambda `$6962 Å and their cross correlation functions which are shown in Figure 13, for the same values of the black hole mass as before, namely $`M=10^7,\mathrm{\hspace{0.17em}10}^8,\mathrm{\hspace{0.17em}10}^9\mathrm{M}_{}`$ and for $`\theta =\pi /3`$. The smoother variations of the UV band lead to a broader CCF, which for $`M=10^9\mathrm{M}_{}`$ leads to a lag of about 1 day, consistent with those given in Collier et al. (1998). However, in this specific value of $`M`$, the UV light curve at $`\lambda `$1315 Å is visibly smoother than the observed light curve at the same wavelength. For the smaller values of $`M`$ the incident and reprocessed UV light curves are in fact almost identical but in this case the correponding lags are a lot smaller, as seen in Figure 13. Perhaps, one should have used as an incident light curve one with substantially more power in the high frequencies than the observed UV light curve, as done in Rokaki & Magnan (1992). However, in the absence of such data one can only speculate.
## 4 Discussion, Conclusions
In the sections above we explored the observational effects of X-ray reprocessing by a geometrically thin, optically thick accretion disk in AGN. We assumed in doing this that a point-like X-ray source irradiates the accretion disk (approximated by an infinite plane) from a given distance above it; the disk then re emits the locally incident X-ray flux in black body form, much in the way considered for producing the BBB spectrum. Our goal has been to compare our results to the observations of the recent mulitwavelength campaign of NGC 7469, which sampled its optical, UV and X-ray light curves at a rate sufficiently high to allow meaningful measurements of the lags between these components. We searched the relevant parameter space to examine whether it is possible to account for the observed correlated variability in the optical – UV bands as the result of reprocessing of the observed X-rays by an accretion disk. This is a relevant question, as this process was suggested to be the origin of the optical – UV lags in previous, less well sampled (in time and wavelength) campaigns.
Our “bottom line” results are summarized by Figure 12 which exhibits the relative lags of our model light curves between the UV ($`\lambda `$1315 Å) and the optical ($`\lambda `$6962 Å) bands, in order to compare directly with the same quantity as determined by the NGC 7469 campaign (see figure 5 in Collier et al. 1998): The model lags are much too short to account for the observations under the assumptions of our model (i.e. that they are due to X-ray reprocessing), for any reasonable value of the black hole mass $`M`$ or the angle $`\theta `$. The model lags are shorter than those observed by, at least, a factor five, for any reasonable value of the black hole mass (or, equivalently, height of the X-ray source above the disk). However, the lags reported by Collier et al. (1998) could be accounted within the model, provided that the variability is due to the reprocessing of a continuum component other than the X-rays, with a light curve similar to that observed in the UV.
This outcome, as indicated also by the detailed X-ray reprocessing model light curves, is the result of the short distances associated with black hole - accretion disk systems of these models. For reasonable values of the black hole mass, $`M`$, the $`\lambda `$1315 Å light curves track almost identically that of X-rays. Even for the optical wavelength, $`\lambda `$6962 Å, the reprocessed light curves follow quite closely those of the X-rays as seen in the relevant figures and as shown by the autocorrelation functions. We have also produced model light curves for two additional wavelengths, namely $`\lambda `$15000 Å and $`\lambda `$30000 Å, both well outside the range of observations, in order to explore at what wavelengths our model could yield light curves with autocorrelation functions similar to those observed in the UV. It was only for the longest wavelength ($`\lambda `$30000 Å) and for $`M=10^8\mathrm{M}_{}`$ that we were able to produce autocorrelation functions resembling that associated with the light curve of NGC 7469 at $`\lambda `$1315 Å. However, in this last case the reprocessed emission comes from such large radii that the resulting RMS variability is much smaller than that the UV observations indicate.
In view of our model light curves of Figures 4, 5, 6 and also as discussed in Nandra et al. (1998), the correlated X-ray – UV variability of NGC 7469 appears really puzzling. One might argue that the variability in the UV – optical bands is intrinsic to that of the accretion disk itself; as argued in Collier et al. (1998) the wavelength dependence of the UV – optical lags are consistent with this view. However this would present the following two problems: (a) the observed lag time scales are too short for the standard accretion disk models (as discussed in the introduction); (b) given the roughly equal luminosities in the X-ray and UV bands, there should exist some evidence of a high frequency component in the UV light curve due to the (expected) reprocessing of the observed X-rays, should the geometry assumed in the present note be valid. Perhaps this last constraint could be alleviated if the albedo of the disk were high, i.e. $`𝒜0.9`$, as this could make the amplitude of the reprocessed component very hard to discern. It has been suggested recently that this may very well be the case under certain conditions in X-ray illuminated disks (Nayakshin, Kazanas & Kallman 1999). In this case however, the presence of such a reflected component with amplitude similar to that of the intrinsic X-ray flux should be apparent in the X-ray autocorrelation function at lags $`\tau R/c`$, where $`R`$ is the height of the X-ray source above the disk. Allowing for additional speculation to remove this last conundrum, one could argue that the X-ray emission comes from magnetic loops (Galeev, Rosner, Vaiana 1979) of sizes smaller than the X-ray sampling time multiplied by the speed of light (i.e. $`R10^{14}`$ cm). Such a solution appears contrived and it still does not explain the origin of the UV – optical lags.
In conclusion, the observations of the correlated multiwavelength variability of NGC 7469 seems to be grossly incompatible with the simplified model examined herein, namely that this variability is due to reprocessing of the observed X-rays by an optically thick, geometrically thin disk, with the X-ray source approximated by a point source at a given distance above this disk. It is not apparent to us how simple modifications to this model could lead to results compatible with these observations. The connection between the X-ray and the longer wavelength (BBB) components appears to be much less direct than allowed through this very simplified model, which suggests much more rapid variability, yet these two spectral components should be somehow related, given their similar overall variability amplitudes.
It is not known whether the correlations of the multiwavelength variations observed in NGC 7469 are a general AGN property or particular to this specific object. These intriguing results suggest that additional studies with similar wavelength coverage and sampling rates are badly needed in order to establish whether our present, general notions of AGN structure are indeed sound or in need of a major revision. Our models and their comparison to observations should be viewed only as a case in support of the argument made in the introduction, namely that spectral fits alone are unable to provide unequivocal information concerning the structure of AGN and generally of sources powered by accretion onto compact objects. Time variability observations and successful modeling are an absolutely necessary supplement to the spectral studies.
We would like to thank P. Nandra for providing us with the X-ray light curve of NGC 7469. We would also like to thank Sergei Nayakshin, Hagai Netzer and Brad Peterson for several interesting discussions. |
warning/0001/hep-ph0001055.html | ar5iv | text | # I INTRODUCTION
## I INTRODUCTION
Supersymmetry(SUSY) is one of the most attractive extensions of the Standard Model(SM). It provides an elegant way to stabilize the huge hierarchy between the electroweak and the GUT scales against radiative corrections. Moreover, supersymmetric models offer a natural solution to the Dark Matter problem and allow for a consistent unification of the all known gauge coupling constants in contrast to the SM. Due to the theoretical appealing of SUSY, the search for supersymmetric particles is one of the main issues in the experimental programs at the CERN $`e^+e^{}`$ collider LEP2 and Fermilab Tevatron . It will play an even more important role at the future Large Hadron Collider (LHC) and the Next $`e^+e^{}`$ Linear Collider.
Although the colored supersymmetric particles, squarks and gluinos, can be searched for most efficiently at hadron colliders, for a precise determination of the underlying SUSY parameters lepton colliders will be necessary. For the experimental search it is useful to predict the production rates of these particles precisely incorporating radiative corrections. Up to now, many works have been devoted to the QCD corrections to various sparticle production rates. QCD corrections to colored sparticle (except stop) production at hadron colliders were discussed in detail by W. Beenakker et al. . The corresponding corrections to the top squark production were given in another paper. The QCD and SUSY-QCD corrections to non-colored sparticle production at hadron colliders were given in and those to colored sparticle production at $`e^+e^{}`$ colliders were given in .
In this paper, we consider the electroweak corrections to the third generation diagonal squark pair production in $`e^+e^{}`$ annihilation, $`e^+e^{}\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i},\stackrel{~}{b}_i\overline{\stackrel{~}{b}}_i`$, due to large Yukawa couplings. Our framework is the Minimal Supersymmetric Standard Model(MSSM). As is well known that there are five physical Higgs bosons in the MSSM, two CP-even neutral Higgs bosons, one CP-odd neutral Higgs boson and a pair of charged Higgs bosons. Their supersymmetric partners, higgsinos, are components of two charginos and four neutralinos in the MSSM. The top and bottom squarks have Yukawa couplings with these Higgs bosons and higgsinos, which are proportional to $`m_t\mathrm{cot}\beta `$ or $`m_b\mathrm{tan}\beta `$, where $`\mathrm{tan}\beta =v_2/v_1`$ and $`v_1`$, $`v_2`$ are the vacuum expectation values of the two Higgs doublets. These interaction terms are large in the region of small or large $`\mathrm{tan}\beta `$ and they can even be leading electroweak corrections for $`\mathrm{tan}\beta 1`$ or $`\mathrm{tan}\beta m_t/m_b`$. On the other hand, the internal gauge bosons may also have large corrections enhanced by large masses due to virtual heavy particle loops such as the top or stop loops. For consistency, we also include such corrections in our calculations. We calculate in the ’t Hooft-Feynman gauge. We find the total corrections are quite large in some regions of the MSSM parameter space allowed by present experiments, which can be larger than the SUSY-QCD corrections to the same process due to gluino exchanges.
This paper is organized as follows. In Sec. II we present the renormalization scheme adopted in our calculation. Some analytic results are given in Sec. III and the numerical results are discussed in Sec IV. Finally we summarize the conclusion in Sec V. The relevant pieces of the Lagrangian are presented in Appendix A and some analytic expressions are collected in Appendix B.
## II RENORMALIZATION SCHEME
In this section we briefly discuss the renormalization scheme adopted in our calculations. To calculate the electroweak corrections to the process $`e^+e^{}\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}(\stackrel{~}{b}_i\overline{\stackrel{~}{b}}_i)`$ at one loop level, we do not need to consider the renormalization of the Higgs sector after imposing the vanishing of the tadpole terms. (However, we adopt an approximate Higgs mass formula including radiative corrections.) Thus, the renormalization scheme is focused on the gauge sector. It differs only slightly from that given by M. Böhm. Another complexity arises from the renormalization of the squark mixing angle.
### A Gauge boson renormalization
The diagonal production of squark pairs proceeds through S-channal photon and $`Z`$ boson exchange at tree level(see Fig. 1). The longitudinal part of $`Z`$ boson does not give any contribution to the process. In the MSSM the photon and $`Z`$ boson may mix with the CP-odd neutral Higgs boson $`A^0`$ and the neutral Goldstone boson $`G^0`$ at one loop level. However, for diagonal production of squark pairs, such mixing does not give any contribution either. So only the renormalization of the transverse part of the gauge bosons is needed.
To respect gauge symmetry explicitly, each gauge multiplet is associated with one renormalization constant:
$`W_\mu ^a(Z_2^W)^{1/2}W_\mu ^a,B_\mu (Z_2^B)^{1/2}B_\mu ,`$ (1)
$`g_2Z_1^W(Z_2^W)^{3/2}g_2,g_1Z_1^B(Z_2^B)^{3/2}g_1.`$ (2)
The Weinberg angle $`\theta _W`$ is defined by the on-shell condition $`\mathrm{cos}\theta _W=\frac{M_W}{M_Z}`$, where $`M_W`$ and $`M_Z`$ are the masses of $`W`$ and $`Z`$ bosons. Now we denote
$$c_W=\mathrm{cos}\theta _W,s_W=\mathrm{sin}\theta _W$$
(3)
as abbreviations throughout the paper and
$`\delta Z_i^\gamma `$ $`=`$ $`s_W^2\delta Z_i^W+c_W^2\delta Z_i^B,\delta Z_i^Z=s_W^2\delta Z_i^B+c_W^2\delta Z_i^W,`$ (4)
$`\delta Z_i^{\gamma Z}`$ $`=`$ $`c_Ws_W(\delta Z_i^W\delta Z_i^B)`$ (5)
as the renormalization constants for the photon, $`Z`$ boson and $`\gamma Z`$ mixing terms respectively. Then we get
$$\left(\genfrac{}{}{0pt}{}{Z}{A}\right)\left(\begin{array}{cc}1+\frac{1}{2}\delta Z_2^Z& \delta Z_1^{\gamma Z}+\delta Z_2^{\gamma Z}\\ \delta Z_1^{\gamma Z}2\delta Z_2^{\gamma Z}& 1+\frac{1}{2}\delta Z_2^\gamma \end{array}\right)\left(\genfrac{}{}{0pt}{}{Z}{A}\right)$$
(6)
from which we can see the $`\gamma Z`$ mixing term.
The renormalization constants $`Z_{1,2}^W,Z_{1,2}^B`$ are fixed by the following on-shell conditions
$`\widehat{\mathrm{\Sigma }}_T^W(M_W^2)=\widehat{\mathrm{\Sigma }}_T^Z(M_Z^2)=\widehat{\mathrm{\Sigma }}_T^{\gamma Z}(0)=0,`$ (7)
$`\widehat{\mathrm{\Gamma }}_\mu ^{\gamma ee}(k^2=0,p/=q/=0)=ie\gamma _\mu ,`$ (8)
$`{\displaystyle \frac{1}{k^2}}\widehat{\mathrm{\Sigma }}^\gamma (k^2)|_{k^2=0}=0`$ (9)
where the $`\widehat{\mathrm{\Sigma }}_T`$s represent the renormalized self-energies and the $`\widehat{\mathrm{\Gamma }}^{\gamma ee}`$ represents the renormalized photon-electron vertex. From $`M_W=g_2/2\sqrt{v_1^2+v_2^2}`$ and $`M_Z=\frac{1}{2}\sqrt{g_1^2+g_2^2}\sqrt{v_1^2+v_2^2}`$ we get
$`{\displaystyle \frac{\delta M_W^2}{M_W^2}}{\displaystyle \frac{\delta M_Z^2}{M_Z^2}}=s_W^2\left[(2\delta Z_1^W3\delta Z_2^W)(2\delta Z_1^B3\delta Z_2^B)\right].`$ (10)
Throughout this paper we shall keep only corrections proportional to a large mass $`M>M_Z`$. All terms independent of the large masses or depending on them only logarithmically will be ommitted. It is found that no terms proportional to large masses enter the expressions $`\frac{\mathrm{\Sigma }^\gamma (k^2)}{k^2}|_{k^2=0}`$ and $`\mathrm{\Sigma }^{\gamma Z}(0)`$. The same is true for $`\delta Z_1^\gamma `$ determined from (2.6). Taking into account these facts we obtain from the renormalization conditions (2.5)-(2.7)
$`s_W^2\delta Z_2^W+c_W^2\delta Z_2^B`$ $`=`$ $`\delta Z_2^\gamma ={\displaystyle \frac{\mathrm{\Sigma }^\gamma (k^2)}{k^2}}|_{k^2=0}=0,`$ (11)
$`\delta Z_1^{\gamma Z}+\delta Z_2^{\gamma Z}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Sigma }^{\gamma Z}(0)}{M_Z^2}}=0,`$ (12)
$`\delta Z_1^\gamma `$ $`=`$ $`0.`$ (13)
The calculations here are similar to those in . The four equations in (10) and (11) completely determine all wave function renormalization constants as only four of them are independent. We get from these equations
$`\delta Z_2^Z`$ $`=`$ $`{\displaystyle \frac{c_W^2s_W^2}{s_W^2}}\left({\displaystyle \frac{\delta M_Z^2}{M_Z^2}}{\displaystyle \frac{\delta M_W^2}{M_W^2}}\right),`$ (14)
$`\delta Z_2^{\gamma Z}`$ $`=`$ $`{\displaystyle \frac{c_W}{s_W}}\left({\displaystyle \frac{\delta M_Z^2}{M_Z^2}}{\displaystyle \frac{\delta M_W^2}{M_W^2}}\right).`$ (15)
The self-energies of gauge bosons $`\mathrm{\Sigma }^\gamma `$, $`\mathrm{\Sigma }^{\gamma Z}`$ and $`\mathrm{\Sigma }_T^Z`$(Fig. 2a) for $`k^20`$ may contain terms proportional to large masses. However, it turns out that their contribution to renormalized gauge boson propagators can be neglected. As an example, let us look into the top-quark loop correction to the renormalized $`Z`$ boson propagator. This propagator can be written as
$`ig_{\mu \nu }\left({\displaystyle \frac{1}{k^2M_Z^2}}{\displaystyle \frac{1}{k^2M_Z^2}}(\mathrm{\Sigma }_T^Z(k^2)\mathrm{\Sigma }_T^Z(M_Z^2)){\displaystyle \frac{1}{k^2M_Z^2}}{\displaystyle \frac{\delta Z_2^Z}{k^2M_Z^2}}\right).`$ (16)
The contribution of the top-quark loop to $`\mathrm{\Sigma }_T^Z(k^2)`$ is written down in (B.20). Although $`\mathrm{\Sigma }_T^Z`$ contains terms proportional to $`m_t^2`$, it can be checked that the combination $`\frac{(\mathrm{\Sigma }_T^Z(k^2)\mathrm{\Sigma }_T^Z(M_Z^2))}{k^2M_Z^2}`$ is not enhanced by $`m_t^2`$ and can be negelected compared to $`\delta Z_2^Z`$ in (16) for all values of $`k^2`$ (See discussion following (B.21)). Thus, $`Z`$ and $`\gamma Z`$ boson propagator can be written as
$$\frac{ig^{\mu \nu }}{k^2M_Z^2}(1\delta Z_2^Z)$$
(17)
and
$$\frac{ig^{\mu \nu }}{k^2M_Z^2}\delta Z_2^{\gamma Z}$$
(18)
respectively.
The analytic expressions for the gauge boson mass corrections calculated from Fig. 3 are given in Appendix B. By using (B.15)-(B.19) we have checked that the divergences in individual terms of the expression (B.4) for $`\frac{\delta M_W^2}{M_W^2}\frac{\delta M_W^2}{M_Z^2}`$ are cancelled out after omitting terms not enhanced by a large mass $`M>M_Z`$. Hence $`\delta Z_2^Z`$ and $`\delta Z_2^{\gamma Z}`$ obtained from (14) are finite. This is also confirmed numerically. It should be noted that the full expression for $`\delta Z_2^\gamma `$ and $`\delta Z_2^{\gamma Z}`$ are of course divergent. The finite results obtained here are consequences of omitting divergent terms independent of large masses $`M>M_Z`$. After cancellation of divergences in the full expression such terms can not induce finite corrections proportional to a large mass.
### B Renormalization of squark wave function
There are two scalar partners $`\stackrel{~}{q}_L`$ and $`\stackrel{~}{q}_R`$ for every quark $`q`$ in SUSY theories. They mix and form two mass eigenstates $`\stackrel{~}{q}_1`$, $`\stackrel{~}{q}_2`$ which are related to the original fields by
$$\left(\genfrac{}{}{0pt}{}{\stackrel{~}{q}_L}{\stackrel{~}{q}_R}\right)=Z_{\stackrel{~}{q}}\left(\genfrac{}{}{0pt}{}{\stackrel{~}{q}_1}{\stackrel{~}{q}_2}\right)$$
(19)
where
$$Z_{\stackrel{~}{q}}=\left(\begin{array}{cc}\mathrm{cos}\theta _{\stackrel{~}{q}}& \mathrm{sin}\theta _{\stackrel{~}{q}}\\ \mathrm{sin}\theta _{\stackrel{~}{q}}& \mathrm{cos}\theta _{\stackrel{~}{q}}\end{array}\right).$$
(20)
We will adopt a scheme in which both stops and sbottoms are defined on shell. We give the formulas for stops here while those for sbottoms are similar.
The complexity of the squark wave function renormalization is due to the fact that the two diagonalized states $`\stackrel{~}{q}_1`$ and $`\stackrel{~}{q}_2`$ mix again at one loop level(See Fig. 4). Write the bare stop fields as
$$\stackrel{~}{t}_i^0=\left(1+\frac{1}{2}\delta Z_i^{\stackrel{~}{t}}\right)\stackrel{~}{t}_i+\delta Z_{ij}^{\stackrel{~}{t}}\stackrel{~}{t}_j,ji.$$
(21)
(We use $`\delta Z_i^{\stackrel{~}{q}}`$ and $`\delta Z_{ij}^{\stackrel{~}{q}}`$ to represent the wave function renormalization constants and $`Z_{\stackrel{~}{q}}`$ the mixing matrix.) The on-shell renormalization conditions require that the mass parameters are the physical masses, the residues of the squark propagators on shell are one and the mixing between on-shell squarks should be absent, i.e.
$`\widehat{\mathrm{\Sigma }}^{1i}(m_{\stackrel{~}{t}_1}^2)=0,\widehat{\mathrm{\Sigma }}^{2i}(m_{\stackrel{~}{t}_2}^2)=0,i=1,2,`$ (22)
$`{\displaystyle \frac{d}{dp^2}}\widehat{\mathrm{\Sigma }}^{11}(p^2)|_{p^2=m_{\stackrel{~}{t}_1}^2}=0,{\displaystyle \frac{d}{dp^2}}\widehat{\mathrm{\Sigma }}^{22}(p^2)|_{p^2=m_{\stackrel{~}{t}_2}^2}=0.`$ (23)
From the above equations, we get
$`\delta m_{\stackrel{~}{t}_i}^2`$ $`=`$ $`\mathrm{\Sigma }(m_{\stackrel{~}{t}_i}^2),`$ (24)
$`\delta Z_i^{\stackrel{~}{t}}`$ $`=`$ $`\mathrm{\Sigma }^{}(m_{\stackrel{~}{t}_i}^2),\delta Z_{ij}^{\stackrel{~}{t}}={\displaystyle \frac{\mathrm{\Sigma }^{ji}(m_{\stackrel{~}{t}_j}^2)}{m_{\stackrel{~}{t}_i}^2m_{\stackrel{~}{t}_j}^2}},`$ (25)
where $`\mathrm{\Sigma }^{}(p^2)`$ is the derivative of $`\mathrm{\Sigma }(p^2)`$ with respect to $`p^2`$.
The wave function renormalization constant matrix can be decomposed into a symmetric and an antisymmetric part
$$\sqrt{Z}=\left(\begin{array}{cc}1+\frac{1}{2}\delta Z_1^{\stackrel{~}{t}}& \frac{1}{2}(\delta Z_{12}^{\stackrel{~}{t}}+\delta Z_{21}^{\stackrel{~}{t}})\\ \frac{1}{2}(\delta Z_{12}^{\stackrel{~}{t}}+\delta Z_{21}^{\stackrel{~}{t}})& 1+\frac{1}{2}\delta Z_2^{\stackrel{~}{t}}\end{array}\right)\left(\begin{array}{cc}1& \frac{1}{2}(\delta Z_{12}^{\stackrel{~}{t}}\delta Z_{21}^{\stackrel{~}{t}})\\ \frac{1}{2}(\delta Z_{12}^{\stackrel{~}{t}}\delta Z_{21}^{\stackrel{~}{t}})& 1\end{array}\right),$$
(26)
where the off-diagonal elements of the symmetric part are ultraviolet finite and the antisymmetric part can be interpreted as a rotation matrix in the first order. Besides the wave function and gauge coupling constant renormalization defined above, an additional renormalization of the stop mixing angle $`\theta _{\stackrel{~}{t}}\theta _{\stackrel{~}{t}}+\delta \theta _{\stackrel{~}{t}}`$ must be introduced to make the $`Z\stackrel{~}{t}_i\overline{\stackrel{~}{t}_j}`$ vertex part finite beyond the tree level. We choose $`\delta \theta _{\stackrel{~}{t}}`$ such that this additional rotation just cancels the last factor in (26), that is,
$$\delta \theta _{\stackrel{~}{t}}=\frac{\delta Z_{12}^{\stackrel{~}{t}}\delta Z_{21}^{\stackrel{~}{t}}}{2}.$$
(27)
This is the same scheme as used in . It is found that with this choice of mixing angle renormalization the ultraviolet divergence in the vertex graph is exactly cancelled.
The analytic expressions for the self energies $`\mathrm{\Sigma }^{ij}`$s calculated from Fig. 4 are given in Appendix B.
### C Renormalization of the gauge boson and squark vertex
With the choice of (27) the complete one-loop electroweak corrected Lagrangian for the gauge boson and stop interaction vertex is given by
$`_{\gamma \stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}}`$ $`=`$ $`\left\{{\displaystyle \frac{2}{3}}ie\left(1+{\displaystyle \frac{\delta e}{e}}+\delta Z_{\stackrel{~}{t}_i}+{\displaystyle \frac{1}{2}}\delta Z_2^\gamma \right){\displaystyle \frac{ie}{s_Wc_W}}\left({\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{t}}^{1i})^2{\displaystyle \frac{2}{3}}s_W^2\right)(\delta Z_1^{\gamma Z}+\delta Z_2^{\gamma Z})\right\}\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}A_\mu ,`$ (28)
$`_{Z\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}}`$ $`=`$ $`{\displaystyle \frac{ie}{s_Wc_W}}\left[{\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{t}}^{1i})^2{\displaystyle \frac{2}{3}}s_W^2\right]\left(1+{\displaystyle \frac{\delta e}{e}}{\displaystyle \frac{\delta \mathrm{cos}\theta _W}{\mathrm{cos}\theta _W}}{\displaystyle \frac{\delta \mathrm{sin}\theta _W}{\mathrm{sin}\theta _W}}+\delta Z_{\stackrel{~}{t}_i}+{\displaystyle \frac{1}{2}}\delta Z_2^Z\right)\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}Z_\mu `$ (32)
$`{\displaystyle \frac{ie}{s_Wc_W}}\left(\mathrm{sin}\theta _{\stackrel{~}{t}}\mathrm{cos}\theta _{\stackrel{~}{t}}{\displaystyle \frac{\delta Z_{12}^{\stackrel{~}{t}}+\delta Z_{21}^{\stackrel{~}{t}}}{2}}\right)\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}Z_\mu `$
$`{\displaystyle \frac{2}{3}}ie(\delta Z_1^{\gamma Z}+\delta Z_2^{\gamma Z})\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}Z_\mu .`$
By using (11), the above two equations are reduced to
$`_{\gamma \stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}}`$ $`=`$ $`{\displaystyle \frac{2}{3}}ie(1+\delta Z_{\stackrel{~}{t}_i})\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}A_\mu ,`$ (33)
$`_{Z\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}}`$ $`=`$ $`{\displaystyle \frac{ie}{s_Wc_W}}\left[{\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{t}}^{1i})^2{\displaystyle \frac{2}{3}}s_W^2\right](1+\delta Z_{\stackrel{~}{t}_i})\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}Z_\mu `$ (35)
$`{\displaystyle \frac{ie}{s_Wc_W}}\left(\mathrm{sin}\theta _{\stackrel{~}{t}}\mathrm{cos}\theta _{\stackrel{~}{t}}{\displaystyle \frac{\delta Z_{12}^{\stackrel{~}{t}}+\delta Z_{21}^{\stackrel{~}{t}}}{2}}\right)\stackrel{~}{t_i}^{}\stackrel{}{^\mu }\stackrel{~}{t_i}Z_\mu .`$
The corresponding Lagrangian for the sbottoms is similar.
## III ANALYTIC RESULTS
### A Vertex corrections
All the quantum effects, including the vertex corrections and the corrections to the $`Z`$ gauge boson propagator but not the $`\gamma Z`$ mixing contributions, can be written in concise forms by defining two effective coupling constants $`D^\gamma `$ and $`D^Z`$. Let
$`\mathrm{\Gamma }_{eff}^{\gamma \stackrel{~}{t}\stackrel{~}{t}}`$ $`=`$ $`ie\left[{\displaystyle \frac{2}{3}}(1+\delta Z_{\stackrel{~}{t}_i})+\mathrm{\Lambda }^\gamma \right]=ieD^\gamma ,`$ (36)
$`\mathrm{\Gamma }_{eff}^{Z\stackrel{~}{t}\stackrel{~}{t}}`$ $`=`$ $`{\displaystyle \frac{ie}{s_Wc_W}}\left[\left({\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{t}}^{1i})^2{\displaystyle \frac{2}{3}}s_W^2\right)\left(1+\delta Z_{\stackrel{~}{t}_i}\delta Z_2^Z\right)\mathrm{sin}\theta _{\stackrel{~}{t}}\mathrm{cos}\theta _{\stackrel{~}{t}}{\displaystyle \frac{(\delta Z_{12}^{\stackrel{~}{t}}+\delta Z_{21}^{\stackrel{~}{t}})}{2}}+\mathrm{\Lambda }^Z\right]`$ (37)
$`=`$ $`{\displaystyle \frac{ie}{s_Wc_W}}D^Z`$ (38)
where $`\mathrm{\Lambda }^\gamma `$ and $`\mathrm{\Lambda }^Z`$ are the vertex corrections by exchanging virtual Higgs bosons, charginos and neutralinos to the $`\gamma `$ and $`Z`$ vertices respectively, as depicted in Fig. 5. It should be noted that the contribution from Fig. 5(f) is zero. The contributions coming from Fig. 5(d) and 5(e) cancel each other because the $`\stackrel{~}{q}\stackrel{~}{q}A^0`$ coupling changes signs when the momentum of the squark changes signs. Our results are analytically and numerically confirmed by that the UV divergence are cancelled precisely as it should be. The results have been verified by testing the Ward identity. The analytic expressions for $`\mathrm{\Lambda }^\gamma `$ and $`\mathrm{\Lambda }^Z`$ are listed in Appendix B.
### B cross section
We need to consider two types of contributions, one is the vertex corrections and the other is the corrections to the internal propagators shown in (17) and (18). We denote the amplitude due to the $`\gamma Z`$ mixing as $`T^{\gamma Z}`$ and $`T^{Z\gamma }`$ where $`T^{\gamma Z}`$ corresponds to the photon propagator on the left and the $`Z`$ boson propagator on the right in Fig. 2 and $`T^{\gamma Z}`$ the vice versa. We use $`T^\gamma `$ and $`T^Z`$ to represent the amplitudes calculated from (36) and (37). Then the cross section can be expressed as
$`\sigma `$ $`=`$ $`{\displaystyle \frac{\pi \alpha ^2}{s}}\left(1{\displaystyle \frac{4m_{\stackrel{~}{t}_i}^2}{s}}\right)^{\frac{3}{2}}\left[L+M{\displaystyle \frac{s^2}{(sM_Z^2)^2}}+N{\displaystyle \frac{s}{(sM_Z^2)}}\right]`$ (39)
where $`s=(p_1+p_2)^2`$ is the s-channal Mandelstam variable(Fig. 1) and $`L`$, $`M`$ and $`N`$ are
$`L`$ $`=`$ $`8\left(D^\gamma \right)^2,`$ (40)
$`M`$ $`=`$ $`(14s_W^2+8s_W^2)\left[\left({\displaystyle \frac{D^Z}{s_W^2c_W^2}}\right)^2+{\displaystyle \frac{\delta Z^{\gamma Z}D^\gamma D^Z}{4s_W^3c_W^3}}\right]+{\displaystyle \frac{\delta Z^{\gamma Z}(D^Z)^2}{2s_W^3c_W^3}}(14s_W^2),`$ (41)
$`N`$ $`=`$ $`{\displaystyle \frac{4D^ZD^\gamma }{s_W^2c_W^2}}(14s_W^2)+{\displaystyle \frac{\delta Z^{\gamma Z}\left(D^\gamma \right)^2}{2c_Ws_W}}(14s_W^2)+{\displaystyle \frac{2\delta Z^{\gamma Z}D^ZD^\gamma }{s_Wc_W}}.`$ (42)
$`L`$ comes from the square of $`T^\gamma `$. The first term in the square bracket of $`M`$ comes from the square of $`T^Z`$ while the second term in it comes from the interference between $`T^{Z\gamma }`$ and $`T^Z`$. The last term in $`M`$ comes from the interference between $`T^{\gamma Z}`$ and $`T^Z`$. The first term in $`N`$ is due to the interference between $`T^\gamma `$ and $`T^Z`$, the second term is due to the interfernce between $`T^{Z\gamma }`$ and $`T^\gamma `$ and the last term is due to the interference between $`T^{\gamma Z}`$ and $`T^\gamma `$.
Throwing away the one-loop corrections, we regain the tree-level formula.
### C Higgs boson mass formula
The Higgs sector is strongly constrained by supersymmetry. In the tree level a light Higgs boson exists with an upper mass bound $`M_Z`$. Radiative corrections can considerably change the Higgs mass spectrum. In our calculations we adopt an approximate Higgs mass formula which incorporates the one loop radiative corrections. It is given by
$`M_{H^0,h^0,eff}^2`$ $`=`$ $`{\displaystyle \frac{M_{A^0}^2+M_Z^2+\omega _t}{2}}`$ (44)
$`\pm \sqrt{{\displaystyle \frac{\left(M_{A^0}^2+M_Z^2\right)^2+\omega _{t}^{}{}_{}{}^{2}}{4}}M_{A^0}^2M_Z^2\mathrm{cos}^22\beta +{\displaystyle \frac{\omega _t\mathrm{cos}2\beta }{2}}(M_{A^0}^2M_Z^2)}`$
where
$`\omega _t`$ $`=`$ $`{\displaystyle \frac{N_cG_Fm_t^4}{\sqrt{2}\pi ^2\mathrm{sin}^2\beta }}(\mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{t}_1}m_{\stackrel{~}{t}_2}}{m_t^2}}+{\displaystyle \frac{A_t(A_t+\mu \mathrm{cot}\beta )}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}}\mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{t}_1}^2}{m_{\stackrel{~}{t}_2}}}`$ (46)
$`+{\displaystyle \frac{A_t^2(A_t+\mu \mathrm{cot}\beta )^2}{(m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2)^2}}(1{\displaystyle \frac{m_{\stackrel{~}{t}_1}^2+m_{\stackrel{~}{t}_2}^2}{m_{\stackrel{~}{t}_1}^2m_{\stackrel{~}{t}_2}^2}}\mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{t}_1}}{m_{\stackrel{~}{t}_2}}})).`$
Let $`\omega _t=0`$ we return to the tree level formula of neutral Higgs boson masses. Although the correction about Higgs masses is a two-loop effect to the squark pair production, we find it can greatly affect the numerical results.
The charged Higgs boson mass is given by
$$M_{H^\pm }^2=M_{A^0}^2+M_W^2.$$
(47)
## IV NUMERICAL RESULTS
Now we turn to discuss the numerical results. Since the cross section of squark pair production is sensitive to the squark masses, we use the two stop masses, $`m_{\stackrel{~}{t}_1}`$ and $`m_{\stackrel{~}{t}_2}`$, as input parameters. Making the following assumptions for simplicity
$`m_{\stackrel{~}{Q}_L}`$ $`=`$ $`m_{\stackrel{~}{t}_R}=m_{\stackrel{~}{b}_R},`$ (48)
$`A_t`$ $`=`$ $`A_b,`$ (49)
where the $`m`$s and $`A`$s are the scalar masses and trilinear soft breaking parameters, we are left with only two free parameters in the squark sector. The sbottom masses are then determined by $`m_{\stackrel{~}{t}_1}`$, $`m_{\stackrel{~}{t}_2}`$, $`\mu `$ and $`\mathrm{tan}\beta `$. For simplicity we also assume the GUT relation $`m_1=(5/3)\mathrm{tan}^2\theta _Wm_2`$ where $`m_1`$ and $`m_2`$ are U(1) and SU(2) gaugino masses respectively. The chargino and neutralino sectors are determined by taking $`m_2`$ as another free parameter. $`m_{A^0}`$ and $`\mathrm{tan}\beta `$ determine the MSSM Higgs sector. These free parameters are constrained by the experimental mass bounds. We impose $`m_{h^0}>90GeV`$, $`m_{\chi _1^0}>35GeV`$, $`m_{\chi _1^+}>95GeV`$ and $`m_{\stackrel{~}{b}_1}>150GeV`$. To discuss the large Yukawa couplings we focus our attention on the regions of small and large $`\mathrm{tan}\beta `$. The MSSM may seem unnatural for these values of $`\mathrm{tan}\beta `$. However, they are actually not excluded by present experiments even for $`m_{h^0}90GeV`$(See e.g. Ref. ). Other parameters are taken as $`\alpha =1/128`$, $`M_W=80.4GeV`$, $`M_Z=91.2GeV`$, $`m_t=174GeV`$, $`m_b=4.7GeV`$ and $`\mathrm{sin}^2\theta _W=0.223`$.
In Fig. 6, we show the cross section $`\sigma (e^+e^{}\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i},\stackrel{~}{b}_i\overline{\stackrel{~}{b}}_i)`$ as a function of the collision energy $`\sqrt{s}`$ for $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$ and $`\mu =m_{A^0}=m_2=400GeV`$ for small and large $`\mathrm{tan}\beta `$ scenarios. For $`\mathrm{tan}\beta =1.5`$ the two sbottom quarks are almost degenerate. However, for $`\mathrm{tan}\beta =30`$ the lighter sbottom can be as light as $`\stackrel{~}{t}_1`$ and its production rates are much larger than those of the heavier one.
We then calculate the corrections to the cross section $`\sigma (e^+e^{}\stackrel{~}{t}_1\stackrel{~}{t}_1)`$ at $`\sqrt{s}=206GeV`$ at which LEP2 can run in 2000 . In Fig. 7, we show $`\delta \sigma /\sigma `$ as a function of the parameter $`\mu `$ by taking $`m_{\stackrel{~}{t}_1}=92GeV`$, which is slightly heavier than the present lower limit, and $`m_{\stackrel{~}{t}_2}=350GeV`$ for $`\mathrm{tan}\beta =1.5`$, $`m_{A^0}=400,800GeV`$ and $`m_2=200,800GeV`$. We can see that the corrections are not sensitive to $`m_2`$. For large $`\mu `$ the corrections can be quite large, which is reasonable since $`\mu `$ directly enters the Higgs boson and squark coupling vertices. They are generally larger than the SUSY-QCD corrections due to gluino exchanges. Fig. 8 shows that $`\delta \sigma /\sigma `$ is also sensitive to $`m_{A^0}`$ and is more sensitive for smaller $`\mathrm{tan}\beta `$.
In Figs. 9–11, we present $`\delta \sigma /\sigma `$ for $`\sqrt{s}=500GeV`$, $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$. For these mass values of the stops $`\sqrt{s}=500GeV`$ is close to the peak for $`\stackrel{~}{t}_1`$ pair production and also to that for $`\stackrel{~}{b}_1`$ pair production at large $`\mathrm{tan}\beta `$ as shown in Fig. 6. Fig. 9 shows $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`\mathrm{tan}\beta =1.5`$ and Fig. 10 shows $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`\mathrm{tan}\beta =0.6`$. We can see that for $`\mathrm{tan}\beta <1`$ the cross section for stop production is greatly suppressed. We find the cusps in the two figures are a threshold effect mainly coming from Fig. 5(i) when $`m_\chi ^{}250GeV`$. Fig. 11 shows the correction as a function of $`m_{A^0}`$ for $`\mu =400GeV`$ and several values of $`\mathrm{tan}\beta `$. From this figure we also see that the corrections are negative and the cross section is suppressed for $`\mathrm{tan}\beta <1`$.
We then discuss a scenario with large SUSY parameters at $`\sqrt{s}=2000GeV`$. We take $`m_{\stackrel{~}{t}_1}=400GeV`$ and $`m_{\stackrel{~}{t}_2}=800GeV`$ in the following discussions. Figs. 12 and 13 show the ratio of the corrections to the tree level result for $`\sigma (e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1},\stackrel{~}{t}_2\overline{\stackrel{~}{t}_2})`$ as a function of $`\mu `$ and $`m_{A^0}`$ respectively when $`m_2=1000GeV`$. The cusp in Fig. 13 mainly comes from Fig. 5(c) when $`m_{H^+}1000GeV`$. The corrections are large for small $`\mathrm{tan}\beta `$. The corrections for $`\sigma (e^+e^{}\stackrel{~}{t}_2\overline{\stackrel{~}{t}_2})`$ show a singularity which stems from the wave function renormalization for $`\stackrel{~}{t}_2`$ at $`m_{A^0}=400GeV`$, where $`m_{\stackrel{~}{t}_2}=m_{\stackrel{~}{t}_1}+m_{A^0}`$. Such singularity was also mentioned in . The corrections can even reach up to 20% in this case.
In Fig. 14 we give the corrections to the sbottom production rates for small and large $`\mathrm{tan}\beta `$ scenarios. When $`\mathrm{tan}\beta =30`$ the corrections for $`\stackrel{~}{b}_1`$ and $`\stackrel{~}{b}_2`$ are both positive and those for $`\stackrel{~}{b}_2`$ can be larger than 20%. For $`\mathrm{tan}\beta =2`$ the corrections tend to increase the $`\stackrel{~}{b}_2`$ production and decrease the $`\stackrel{~}{b}_1`$ production.
Finally, we will compare the contributions from three classes of diagrams, i.e., (1) the vertex and squark wave function corrections from Higgs bosons, (2) those from charginos and neutralinos, and (3) corrections to the gauge boson propagator. Especially we will show the importance of the third class. We find this contribution is sensitive to the mass difference between the two scalar top quark mass eigenstates. In Fig. 15 we give three classes of contributions as functions of parameter $`\mu `$ for $`\sqrt{s}=1500GeV`$, $`\mathrm{tan}\beta =3`$, $`m_{A^0}=400GeV`$, $`m_2=800GeV`$, $`m_{\stackrel{~}{t_1}}=300GeV`$ and $`m_{\stackrel{~}{t_2}}=500GeV,800GeV`$. In fig. 16 we give the same quantities as functions of parameter $`m_{\stackrel{~}{t_1}}`$ by fixing $`m_{\stackrel{~}{t_2}}=800GeV`$, $`\mu =300GeV`$ and take other parameters the same as those in Fig. 15. We find that the contribution from corrections to gauge boson propagator can be as large as about 7% of the total cross section. It can be seen from these figures that the contribution from corrections to the gauge boson propagator is larger than those from charginos and neutralinos in a large region of the parameter space and it is opposite in sign and not much smaller than the contributions from Higgs bosons.
## V SUMMARY AND CONCLUSION
In summary, we have calculated the large Yukawa coupling corrections to the diagonal stop and sbottom pair production in $`e^+e^{}`$ annihilation. We include also terms of the self-energy corrections to gauge bosons enhanced by large masses. We discuss the corrections as functions of different SUSY parameters. They are found to be quite significant and are larger than the SUSY-QCD corrections by gluino exchanges in a large region of the MSSM parameter space. They can even be comparable to the conventional QCD corrections which is about 20%. The corrections can be both positive and negative. We find the corrections are quite sensitive to the parameters $`\mu `$, $`m_{A^0}`$ and $`\mathrm{tan}\beta `$. They are not sensitive to the gaugino mass $`m_2`$. In conclusion, when one consider the third generation squark production in the MSSM such corrections should not be ignored if precision prediction is needed.
After we finished the work we became aware of the work of H. Eberl et al in which a large part of this work had been done. However, they did not include the contributions coming from corrections to the gauge boson propagator. As discussed at the end of the last section this contribution is sizable in a large region of parameter space and should be taken into account for consistency. The corrections to the gauge boson propagator has been calculated in a different renormalization scheme in connection to the process of chargino pair production. The effects of charginos, neutralinos and Higgs bosons in the loop are not considered in that paper. Apart from this difference we agree with the analytical formulas given in their paper. Our numerical results contain studies of effects of varying different supersymmetric parameters which are not given in ref. We have checked that by taking the same parameter values as in and neglecting the corrections to the gauge boson propagator we obtain numerical results close to theirs.
## ACKNOWLEDGMENTS
Y-B. Dai’s work is supported by the National Science Foundation of China.
Appendix A
In this appendix we list the relevant pieces of the SUSY Lagrangian in terms of the mass eigenstates. We follow the conventions of ref, where the full Lagrangian and the complete set of Feynman rules for the MSSM are given. Some abbreviations of the vertex couplings are defined here, which will appear in the analytic expressions in next appendix.
$`_{A^0\stackrel{~}{t}_i^{}\stackrel{~}{t}_j}`$ $`=`$ $`(ij){\displaystyle \frac{g_2m_t}{2M_W}}(\mu A_t\mathrm{cot}\beta ),`$ (A.1)
$`_{G^0\stackrel{~}{t}_i^{}\stackrel{~}{t}_j}`$ $`=`$ $`(ji){\displaystyle \frac{g_2m_t}{2M_W}}(\mu \mathrm{cot}\beta +A_t),`$ (A.2)
$`_{H_k^0\stackrel{~}{t}_{i}^{}{}_{}{}^{}\stackrel{~}{t}_j}`$ $`=`$ $`ig_2[{\displaystyle \frac{2}{3}}M_W\mathrm{tan}^2\theta _WB_R^k(\delta ^{ij}+{\displaystyle \frac{38s_W^2}{4s_W^2}}Z_{\stackrel{~}{t}}^{1i}Z_{\stackrel{~}{t}}^{1j})+{\displaystyle \frac{m_t^2}{M_W\mathrm{sin}\beta }}Z_R^{2k}\delta ^{ij}`$ (A.4)
$`{\displaystyle \frac{m_t}{2M_Ws_W}}(Z_{\stackrel{~}{t}}^{1i}Z_{\stackrel{~}{t}}^{2j}+Z_{\stackrel{~}{t}}^{1j}Z_{\stackrel{~}{t}}^{2i})(A_tZ_R^{2k}+\mu Z_R^{1k})]`$
$`=`$ $`ig_2\mathrm{\Gamma }_{ijk},`$ (A.5)
$`_{H^+\stackrel{~}{t}_{i}^{}{}_{}{}^{}\stackrel{~}{b}_j}`$ $`=`$ $`ig_2[{\displaystyle \frac{1}{\sqrt{2}}}(M_W\mathrm{sin}2\beta +{\displaystyle \frac{m_b^2}{M_W}}\mathrm{tan}\beta +{\displaystyle \frac{m_t^2}{M_W}}\mathrm{cot}\beta )Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{t}}^{1i}+{\displaystyle \frac{m_tm_b}{\sqrt{2}M_Ws_Wc_W}}Z_{\stackrel{~}{b}}^{2j}Z_{\stackrel{~}{t}}^{2i}`$ (A.7)
$`+(\mu A_t\mathrm{cot}\beta ){\displaystyle \frac{m_t}{\sqrt{2}M_W}}Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{t}}^{2i}+(\mu A_b\mathrm{tan}\beta ){\displaystyle \frac{m_b}{\sqrt{2}M_W}}Z_{\stackrel{~}{b}}^{2j}Z_{\stackrel{~}{t}}^{1i}]`$
$`=`$ $`ig_2D_{ij},`$ (A.8)
$`_{G^+\stackrel{~}{t}_{i}^{}{}_{}{}^{}\stackrel{~}{b}_j}`$ $`=`$ $`ig_2[{\displaystyle \frac{1}{\sqrt{2}}}(M_W\mathrm{cos}2\beta +{\displaystyle \frac{m_t^2}{M_W}}{\displaystyle \frac{m_b^2}{M_W}})Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{t}}^{1i}`$ (A.10)
$`(\mu \mathrm{cot}\beta +A_t){\displaystyle \frac{m_t}{\sqrt{2}M_W}}Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{t}}^{2i}+(\mu \mathrm{tan}\beta +A_b){\displaystyle \frac{m_b}{\sqrt{2}M_W}}Z_{\stackrel{~}{b}}^{2j}Z_{\stackrel{~}{t}}^{1i}]`$
$`=`$ $`ig_2D_{ij}^{},`$ (A.11)
$`_{A^0\stackrel{~}{b}_i^{}\stackrel{~}{b}_j}`$ $`=`$ $`(ij){\displaystyle \frac{g_2m_b}{2M_W}}(\mu A_b\mathrm{tan}\beta ),`$ (A.12)
$`_{G^0\stackrel{~}{b}_i^{}\stackrel{~}{b}_j}`$ $`=`$ $`(ji){\displaystyle \frac{g_2m_b}{2M_W}}(\mu \mathrm{tan}\beta +A_b),`$ (A.13)
$`_{H_k^0\stackrel{~}{b}_{i}^{}{}_{}{}^{}\stackrel{~}{b}_j}`$ $`=`$ $`ig_2[{\displaystyle \frac{M_W}{3}}\mathrm{tan}^2\theta _WB_R^k(\delta ^{ij}+{\displaystyle \frac{34s_W^2}{2s_W^2}}Z_{\stackrel{~}{b}}^{1i}Z_{\stackrel{~}{b}}^{1j}){\displaystyle \frac{m_b^2}{M_W\mathrm{cos}\beta }}Z_R^{1k}\delta ^{ij}`$ (A.15)
$`+{\displaystyle \frac{m_b}{2M_W\mathrm{cos}\beta }}(Z_{\stackrel{~}{b}}^{1i}Z_{\stackrel{~}{b}}^{2j}+Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{b}}^{2i})(A_bZ_R^{1k}+\mu Z_R^{2k})],`$
$`_{\overline{\chi }_j^0\stackrel{~}{t}_it}`$ $`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}[({\displaystyle \frac{Z_{\stackrel{~}{t}}^{1i}}{c_W}}({\displaystyle \frac{Z_N^{1j}s_W}{3}}+Z_N^{2j}c_W)+{\displaystyle \frac{m_t}{M_W\mathrm{sin}\beta }}Z_{\stackrel{~}{t}}^{2i}Z_N^{4j})P_L`$ (A.17)
$`+({\displaystyle \frac{4\mathrm{tan}\theta _W}{3}}Z_{\stackrel{~}{t}}^{2i}Z_N^{1j}+{\displaystyle \frac{m_t}{M_W\mathrm{sin}\beta }}Z_{\stackrel{~}{t}}^{1i}Z_N^{4j})P_R]`$
$`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}\left[R_{ij}P_L+S_{ij}P_R\right],`$ (A.18)
$`_{\overline{\chi }_j^{}\stackrel{~}{t}_i^{}b}`$ $`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}[(\sqrt{2}Z_{\stackrel{~}{t}}^{1i}Z_{}^{+}{}_{}{}^{1j}+{\displaystyle \frac{m_t}{M_W\mathrm{sin}\beta }}Z_{\stackrel{~}{t}}^{2i}Z_{}^{+}{}_{}{}^{2j})P_L`$ (A.20)
$`+\left({\displaystyle \frac{m_b}{M_W\mathrm{cos}\beta }}Z_{\stackrel{~}{t}}^{2i}Z_{}^{}{}_{}{}^{2j}\right)P_R]`$
$`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}\left[U_{ij}P_L+V_{ij}P_R\right],`$ (A.21)
$`_{\overline{\chi }_j^0\stackrel{~}{b}_ib}`$ $`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}[({\displaystyle \frac{Z_{\stackrel{~}{b}}^{1i}}{c_W}}({\displaystyle \frac{Z_N^{1j}s_W}{3}}Z_N^{2j}c_W)+{\displaystyle \frac{m_b}{M_W\mathrm{cos}\beta }}Z_{\stackrel{~}{b}}^{2i}Z_N^{3j})P_L`$ (A.23)
$`+({\displaystyle \frac{2\mathrm{tan}\theta _W}{3}}(Z_{\stackrel{~}{b}}^{2i}Z_N^{1j})+{\displaystyle \frac{m_b}{M_W\mathrm{sin}\beta }}Z_{\stackrel{~}{b}}^{1i}Z_N^{3j})P_R],`$
$`_{\overline{\chi }_j^+\stackrel{~}{b}_it}`$ $`=`$ $`{\displaystyle \frac{ig_2}{\sqrt{2}}}[(\sqrt{2}Z_{\stackrel{~}{b}}^{1i}Z_{}^{}{}_{}{}^{1j}+{\displaystyle \frac{m_b}{M_W\mathrm{cos}\beta }}Z_{\stackrel{~}{b}}^{2i}Z_{}^{}{}_{}{}^{2j})P_L`$ (A.25)
$`+\left({\displaystyle \frac{m_t}{M_W\mathrm{sin}\beta }}Z_{\stackrel{~}{b}}^{1i}Z_{}^{+}{}_{}{}^{2j}\right)P_R].`$
In the above expressions, $`Z_R`$, $`Z_N`$ and $`Z^+`$ and $`Z^{}`$ are the mixing matrices for the two neutral CP-even Higgs bosons, neutralinos and charginos respectively.
$`Z_R`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{cos}\alpha \hfill & \mathrm{sin}\alpha \hfill \\ \mathrm{sin}\alpha \hfill & \mathrm{cos}\alpha \hfill \end{array}\right),`$ (A.28)
$`B_R^k`$ $`=`$ $`\{\begin{array}{cc}\mathrm{cos}(\alpha +\beta ),\hfill & k=1,\hfill \\ \mathrm{sin}(\alpha +\beta ),\hfill & k=2,\hfill \end{array}`$ (A.31)
and
$$\mathrm{tan}2\alpha =\mathrm{tan}2\beta \frac{m_{A^0}^2+M_Z^2}{m_{A^0}^2M_Z^2}.$$
(A.33)
Appendix B
In this appendix we give some analytic results in our calculations. The vertex corrections in Eq. (36) and (37) are given by
$`\mathrm{\Lambda }^\gamma `$ $`=`$ $`{\displaystyle \frac{g_2^2}{(4\pi )^2}}\{{\displaystyle \frac{2}{3}}[\left({\displaystyle \frac{m_t}{2M_W}}\right)^2(\mu A_t\mathrm{cot}\beta )^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{A^0}^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\alpha }^2]`$ (B.16)
$`+\left({\displaystyle \frac{m_t}{2M_W}}\right)^2(\mu \mathrm{cot}\beta +A_t)^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,M_Z^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\alpha }^2]`$
$`+(\mathrm{\Gamma }_{i\alpha k}\mathrm{\Gamma }_{i\alpha k})(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{H_k^0}^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\alpha }^2]]`$
$`+(D_{ij})^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{b}_j}^2,m_{H^+}^2,m_{H^+}^2]`$
$`+(D_{ij}^{})^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{b}_j}^2,M_W^2,M_W^2]`$
$`{\displaystyle \frac{1}{3}}(D_{ij})^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{H^+}^2,m_{\stackrel{~}{b}_j}^2,m_{\stackrel{~}{b}_j}^2]`$
$`{\displaystyle \frac{1}{3}}(D_{ij}^{})^2(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,M_W^2,m_{\stackrel{~}{b}_j}^2,m_{\stackrel{~}{b}_j}^2]`$
$`{\displaystyle \frac{2}{3}}[2(R_{ij}S_{ij})m_{\chi _j^0}m_t(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^0}^2,m_t^2,m_t^2]`$
$`+(R_{ij}^2+S_{ij}^2)`$
$`((m_t^2+m_{\stackrel{~}{t}_i}^2+m_{\chi _j^0}^2)C_1+m_{\chi _j^0}^2C_0)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^0}^2,m_t^2,m_t^2]+{\displaystyle \frac{1}{2}}B_0[s,m_t^2,m_t^2]]`$
$`[2(U_{ij}V_{ij})m_{\chi _j^+}m_b(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_b^2,m_{\chi _j^+}^2,m_{\chi _j^+}^2]`$
$`+(U_{ij}^2+V_{ij}^2)`$
$`((m_b^2+m_{\stackrel{~}{t}_i}^2+m_{\chi _j^+}^2)C_1+m_b^2C_0)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_b^2,m_{\chi _j^+}^2,m_{\chi _j^+}^2]+{\displaystyle \frac{1}{2}}B_0[s,m_{\chi _j^+}^2,m_{\chi _j^+}^2]]`$
$`+{\displaystyle \frac{1}{3}}[2(U_{ij}V_{ij})m_{\chi _j^+}m_b(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^+}^2,m_b^2,m_b^2]`$
$`+(U_{ij}^2+V_{ij}^2)`$
$`((m_b^2+m_{\stackrel{~}{t}_i}^2+m_{\chi _j^+}^2)C_1+m_{\chi _j^+}^2C_0)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^+}^2,m_b^2,m_b^2]+{\displaystyle \frac{1}{2}}B_0[s,m_b^2,m_b^2]]\}`$
$`\mathrm{\Lambda }^Z`$ $`=`$ $`{\displaystyle \frac{g_2^2}{(4\pi )^2}}\{\left({\displaystyle \frac{m_t}{2M_W}}\right)^2(\mu A_t\mathrm{cot}\beta )^2F_{\alpha \alpha }(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{A^0}^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\alpha }^2]`$ (B.36)
$`+\left({\displaystyle \frac{m_t}{2M_W}}\right)^2(\mu \mathrm{cot}\beta +A_t)^2F_{\alpha \alpha }(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,M_Z^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\alpha }^2]`$
$`+(\mathrm{\Gamma }_{i\alpha k}\mathrm{\Gamma }_{i\beta k})F_{\alpha \beta }\left(C_0+2C_1\right)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{H_k^0}^2,m_{\stackrel{~}{t}_\alpha }^2,m_{\stackrel{~}{t}_\beta }^2]`$
$`+(D_{ij})^2(0.5s_W^2)(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{b}_j}^2,m_{H^+}^2,m_{H^+}^2]`$
$`+(D_{ij}^{})^2(0.5s_W^2)(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{b}_j}^2,M_W^2,M_W^2]`$
$`(D_{ij})^2G_{ij}(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{H^+}^2,m_{\stackrel{~}{b}_j}^2,m_{\stackrel{~}{b}_j}^2]`$
$`(D_{ij}^{})^2G_{ij}(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,M_W^2,m_{\stackrel{~}{b}_j}^2,m_{\stackrel{~}{b}_j}^2]`$
$`{\displaystyle \frac{1}{2}}[((Z_{\stackrel{~}{t}}^{2i})^2(Z_{\stackrel{~}{t}}^{1i})^2)Z_N^{4j}Z_N^{4k}(Z_N^{4j}Z_N^{4k}Z_N^{3j}Z_N^{3k})\left({\displaystyle \frac{m_t}{m_W\mathrm{sin}\beta }}\right)^2((m_{\chi _j^0}m_{\chi _k^0}m_t^2m_{\stackrel{~}{t}_i}^2)`$
$`C_1m_t^2C_0)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_t^2,m_{\chi _k^0}^2,m_{\chi _j^0}^2]{\displaystyle \frac{1}{2}}B_0[s,m_{\chi _k^0}^2,m_{\chi _j^0}]]`$
$`{\displaystyle \frac{1}{2}}[(R_{ij}S_{ij})m_{\chi _j^0}m_t(1{\displaystyle \frac{8}{3}}s_W^2)(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^0}^2,m_t^2,m_t^2]`$
$`+(R_{ij}^2{\displaystyle \frac{4}{3}}s_W^2(R_{ij}^2+S_{ij}^2))`$
$`\left(\left((m_{\stackrel{~}{t}_i}^2+m_{\chi _j^0}^2)C_1+m_{\chi _j^0}^2C_0\right)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^0}^2,m_t^2,m_t^2]+{\displaystyle \frac{1}{2}}B_0[s,m_t^2,m_t^2]\right)`$
$`+(S_{ij}^{}{}_{}{}^{2}{\displaystyle \frac{4}{3}}s_W^2(R_{ij}^{}{}_{}{}^{2}+S_{ij}^{}{}_{}{}^{2}))m_t^2C_1[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^0}^2,m_t^2,m_t^2]]`$
$`{\displaystyle \frac{1}{2}}\mathrm{cos}(2\theta _W)[2(U_{ij}V_{ij})m_{\chi _j^+}m_b(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_b^2,m_{\chi _j^+}^2,m_{\chi _j^+}^2]`$
$`+(U_{ij}^2+V_{ij}^2)`$
$`((m_b^2+m_{\stackrel{~}{t}_i}^2+m_{\chi _j^+}^2)C_1+m_b^2C_0)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_b^2,m_{\chi _j^+}^2,m_{\chi _j^+}^2]+{\displaystyle \frac{1}{2}}B_0[s,m_{\chi _j^+}^2,m_{\chi _j^+}^2]]`$
$`+{\displaystyle \frac{1}{2}}[(U_{ij}V_{ij})m_{\chi _j^+}m_b(1{\displaystyle \frac{4}{3}}s_W^2)(C_0+2C_1)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^+}^2,m_b^2,m_b^2]`$
$`+(U_{ij}^2{\displaystyle \frac{2}{3}}s_W^2(U_{ij}^2+V_{ij}^2))`$
$`\left(\left((m_{\stackrel{~}{t}_i}^2+m_{\chi _j^+}^2)C_1+m_{\chi _j^+}^2C_0\right)[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^+}^2,m_b^2,m_b^2]+{\displaystyle \frac{1}{2}}B_0[s,m_b^2,m_b^2]\right)`$
$`+(V_{ij}^{}{}_{}{}^{2}{\displaystyle \frac{2}{3}}s_W^2(U_{ij}^{}{}_{}{}^{2}+V_{ij}^{}{}_{}{}^{2}))m_b^2C_1[m_{\stackrel{~}{t}_i}^2,s,m_{\stackrel{~}{t}_i}^2,m_{\chi _j^+}^2,m_b^2,m_b^2]]\}.`$
The analytic expressions for stop self energies are
$`\mathrm{\Sigma }^{ij}(p^2)`$ $`=`$ $`{\displaystyle \frac{g_2^2}{(4\pi )^2}}\{(1\delta ^{ij})\left({\displaystyle \frac{m_t}{2M_W}}\right)^2(\mu A_t\mathrm{cot}\beta )^2B_0[p^2,m_{\stackrel{~}{t}_j}^2,m_{A^0}^2]`$ (B.47)
$`+(1\delta ^{ij})\left({\displaystyle \frac{m_t}{2M_W}}\right)^2\left(\mu \mathrm{cot}\beta +A_t\right)^2B_0[p^2,m_{\stackrel{~}{t}_j}^2,m_Z^2]`$
$`+(\mathrm{\Gamma }_{i\alpha k}\mathrm{\Gamma }_{j\alpha k})B_0[p^2,m_{\stackrel{~}{t}_\alpha }^2,m_{H_k^0}^2]+(D_{j\alpha }D_{\alpha i})B_0[p^2,m_{\stackrel{~}{b}_\alpha }^2,m_{H^+}^2]+(D_{j\alpha }^{}D_{\alpha i}^{})B_0[p^2,m_{\stackrel{~}{b}_\alpha }^2,m_W^2]`$
$`2\mathrm{sin}\theta _{\stackrel{~}{t}}\mathrm{cos}\theta _{\stackrel{~}{t}}\delta ^{ij}\left({\displaystyle \frac{3+2s_W^2}{12c_W^2}}\mathrm{cos}(2\beta ){\displaystyle \frac{m_t^2\mathrm{cot}\beta }{2m_W^2}}+{\displaystyle \frac{m_b^2\mathrm{tan}^2\beta }{2m_W^2}}\right)A_0[m_{H^\pm }^2]`$
$`+\mathrm{sin}\theta _{\stackrel{~}{t}}\mathrm{cos}\theta _{\stackrel{~}{t}}{\displaystyle \frac{38s_W^2}{12s_W^2}}\delta ^{ij}\left(\mathrm{cos}(2\beta )A_0[m_{A^0}^2]+\mathrm{cos}(2\alpha )(A_0[m_{H^0}^2]A_0[m_{h^0}^2])\right)`$
$`[m_{\chi _k^0}m_t(R_{jk}S_{ik}+S_{jk}R_{ik})B_0[p^2,m_t^2,m_{\chi _k^0}^2]`$
$`+\left(R_{jk}R_{ik}+S_{jk}S_{ik}\right)\left(A_0[m_{\chi _k^0}^2]+m_t^2B_0[p^2,m_t^2,m_{\chi _k^0}^2]\right)`$
$`+(R_{jk}R_{ik}+S_{jk}S_{ik})p^2B_1[p^2,m_t^2,m_{\chi _k^0}^2]]`$
$`[m_{\chi _k^0}m_b(U_{jk}V_{ik}+V_{jk}U_{ik})B_0[p^2,m_b^2,m_{\chi _k^{}}^2]`$
$`+\left(U_{jk}U_{ik}+V_{jk}V_{ik}\right)\left(A_0[m_{\chi _k^{}}^2]+m_b^2B_0[p^2,m_b^2,m_{\chi _k^{}}^2]\right)`$
$`+(U_{jk}U_{ik}+V_{jk}V_{ik})p^2B_1[p^2,m_b^2,m_{\chi _k^{}}^2]]\}`$
The analytic expressions for the gauge boson mass corrections are
$`{\displaystyle \frac{\delta M_Z^2}{M_Z^2}}{\displaystyle \frac{\delta M_W^2}{M_W^2}}`$ $`=`$ $`{\displaystyle \frac{N_cg_2^2m_t^2}{64\pi ^2M_W^2}}{\displaystyle \frac{g_2^2}{16\pi ^2M_W^2}}\{\mathrm{sin}^2(\alpha \beta )B_{00}[M_Z^2,m_{A^0}^2,m_{H^0}^2]`$ (B.62)
$`+\mathrm{cos}^2(\alpha \beta )B_{00}[M_Z^2,m_{A^0}^2,m_{h^0}^2]`$
$`\left[\mathrm{sin}^2(\alpha \beta )B_{00}[M_W^2,m_{H^+}^2,m_{H^0}^2]+\mathrm{cos}^2(\alpha \beta )B_{00}[M_W^2,m_{H^+}^2,m_{h^0}^2]\right]`$
$`B_{00}[M_W^2,m_{H^+}^2,m_{A^0}^2]+{\displaystyle \frac{1}{2}}A_0[m_{H^+}^2]`$
$`+4\left[{\displaystyle \frac{Z_{\stackrel{~}{t}}^{1i}Z_{\stackrel{~}{t}}^{1j}}{2}}{\displaystyle \frac{2s_W^2\delta ^{ij}}{3}}\right]^2B_{00}[M_Z^2,m_{\stackrel{~}{t}_i}^2,m_{\stackrel{~}{t}_j}^2]`$
$`2\left[{\displaystyle \frac{4s_W^4}{9}}+{\displaystyle \frac{38s_W^2}{12}}(Z_{\stackrel{~}{t}}^{1i})^2\right]A_0[m_{\stackrel{~}{t}_i}^2]`$
$`+4\left[{\displaystyle \frac{Z_{\stackrel{~}{b}}^{1i}Z_{\stackrel{~}{b}}^{1j}}{2}}{\displaystyle \frac{s_W^2\delta ^{ij}}{3}}\right]^2B_{00}[M_Z^2,m_{\stackrel{~}{b}_i}^2,m_{\stackrel{~}{b}_j}^2]`$
$`2\left[{\displaystyle \frac{s_{W}^{}{}_{}{}^{4}}{9}}+{\displaystyle \frac{34s_W^2}{12}}(Z_{\stackrel{~}{b}}^{1i})^2\right]A_0[m_{\stackrel{~}{b_i}}^2]`$
$`2(Z_{\stackrel{~}{b}}^{1i}Z_{\stackrel{~}{t}}^{1j})^2B_{00}[M_Z^2,m_{\stackrel{~}{t}_j}^2,m_{\stackrel{~}{t}_i}^2]`$
$`+{\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{t}}^{1i})^2A_0[m_{\stackrel{~}{t}_i}^2]+{\displaystyle \frac{1}{2}}(Z_{\stackrel{~}{b}}^{1i})^2A_0[m_{\stackrel{~}{b}_i}^2]`$
$`(Z_N^{4i}Z_N^{4j}Z_N^{3i}Z_N^{3j})^2(A0[m_{\chi _i^0}^2](m_{\chi _i^0}m_{\chi _j^0}+m_{\chi _j^0}^2)B_0[M_Z^2,m_{\chi _i^0}^2,m_{\chi _j^0}^2]`$
$`+2B_{00}[M_Z^2,m_{\chi _i^0}^2,m_{\chi _j^0}^2])`$
$`+((Z_N^{4i}Z_{}^{+}{}_{}{}^{2j})^2+(Z_N^{3i}Z_{}^{}{}_{}{}^{2j})^2)(2B_{00}[M_W^2,m_{\chi _j^{}}^2,m_{\chi _i^0}^2]A0[m_{\chi _i^0}^2]`$
$`m_{\chi _j^{}}^2B_0[M_W^2,m_{\chi _j^{}}^2,m_{\chi _i^0}^2])`$
$`2m_{\chi _i^0}m_{\chi _j^{}}(Z_N^{4i}Z_{}^{+}{}_{}{}^{2j}Z_N^{3i}Z_{}^{}{}_{}{}^{2j})B_0[M_W^2,m_{\chi _j^{}}^2,m_{\chi _i^0}^2]\}`$
$`\mathrm{\Gamma }_{ijk}`$, $`D_{ij}`$, $`D_{ij}^{}`$, $`R_{ij}`$, $`S_{ij}`$, $`U_{ij}`$ and $`V_{ij}`$ in the above expressions are the vertex couplings defined in Appendix A. In the concrete calculations we only keep the higgsino sector in the $`R_{ij}`$, $`S_{ij}`$, $`U_{ij}`$ and $`V_{ij}`$. The other two coupling constants
$`F_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Z_{\stackrel{~}{t}}^{1j}Z_{\stackrel{~}{t}}^{1j}{\displaystyle \frac{2}{3}}s_W^2\delta ^{ij},`$ (B.63)
$`G_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}Z_{\stackrel{~}{b}}^{1j}Z_{\stackrel{~}{b}}^{1j}{\displaystyle \frac{1}{3}}s_W^2\delta ^{ij}`$ (B.64)
are the couplings of $`Z`$ boson to top squark and bottom squark respectively. The relevant scalar functions are defined as follows
$`A_0(m^2)=(i\pi ^2)^1(2\pi \mu )^{4D}{\displaystyle d^Dq(q^2m^2)^1},`$ (B.65)
$`B_0(p_1^2,m_0^2,m_1^2)=(i\pi ^2)^1(2\pi \mu )^{4D}{\displaystyle d^Dq[(q^2m_0^2)((q+p_1)^2m_1^2)]^1},`$ (B.66)
$`C_0(p_1^2,p_{12},p_2^2,m_0^2,m_1^2,m_2^2)`$ (B.67)
$`\text{123}=(i\pi ^2)^1(2\pi \mu )^{4D}{\displaystyle d^Dq[(q^2m_0^2)((q+p_1)^2m_1^2)((q+p_2)^2m_2^2)]^1},`$ (B.68)
in which $`p_{ij}=(p_ip_j)^2`$.
The definitions of the tensor integrals and the relevant decompositions are given below
$`T_{\mu _1\mathrm{}\mu _p}(p_1,\mathrm{},p_{N1},m_0,\mathrm{},m_{N1})={\displaystyle \frac{(2\pi \mu )^{4D}}{i\pi ^2}}{\displaystyle d^Dq\frac{q_{\mu _1}\mathrm{}q_{\mu _n}}{D_0D_1\mathrm{}D_{N1}}},`$ (B.69)
with the denominator factors $`D_0=q^2m_0^2,D_i=(q+p_i)^2m_i^2`$ (i=1,$`\mathrm{}`$,N-1) and $`T=B,C,D\mathrm{}`$ corresponding to $`N=2,3,4\mathrm{}`$.
$`B_\mu =p_{1}^{}{}_{\mu }{}^{}B_1,`$ (B.70)
$`B_{\mu \nu }=g_{\mu \nu }B_{00}+p_{1}^{}{}_{\mu }{}^{}p_{1}^{}{}_{\nu }{}^{}B_{11},`$ (B.71)
$`C_\mu =p_{1}^{}{}_{\mu }{}^{}C_1+p_{2}^{}{}_{\mu }{}^{}C_2={\displaystyle \underset{i=1}{\overset{2}{}}}p_{i\mu }C_i,`$ (B.72)
$`C_{\mu \nu }=g_{\mu \nu }C_{00}+p_{1}^{}{}_{\mu }{}^{}p_{1}^{}{}_{\nu }{}^{}C_{11}+p_{2}^{}{}_{\mu }{}^{}p_{2}^{}{}_{\nu }{}^{}C_{22}+(p_{1}^{}{}_{\mu }{}^{}p_{2}^{}{}_{\nu }{}^{}+p_{2}^{}{}_{\mu }{}^{}p_{1}^{}{}_{\nu }{}^{})C_{12}`$ (B.73)
$`\text{123}=g_{\mu \nu }C_{00}+{\displaystyle \underset{i,j=1}{\overset{2}{}}}p_{i}^{}{}_{\mu }{}^{}p_{j}^{}{}_{\nu }{}^{}C_{ij}.`$ (B.74)
The analytic expressions of $`A_0(m^2)`$, $`B_0(p^2,m_0^2,m_1^2)`$ and $`B_1(p^2,m_0^2,m_1^2)`$ can be easily obtained and the corresponding divergences are:
$`A_0(m^2)`$ $`=`$ $`m^2\left({\displaystyle \frac{2}{4D}}\gamma _E+\mathrm{ln}4\pi \right)+\mathrm{},`$ (B.75)
$`B_0(p^2,m_0^2,m_1^2)`$ $`=`$ $`{\displaystyle \frac{2}{4D}}\gamma _E+\mathrm{ln}4\pi +\mathrm{},`$ (B.76)
$`B_1(p^2,m_0^2,m_1^2)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{2}{4D}}\gamma _E+\mathrm{ln}4\pi \right)+\mathrm{}.`$ (B.77)
$`B_{00}`$ can be expressed by $`A_0`$, $`B_1`$ and $`B_0`$ as follows
$`B_{00}(p^2,m_0^2,m_1^2)`$ $`=`$ $`{\displaystyle \frac{1}{6}}\{A_0(m_1^2)+2m_0^2B_0(p^2,m_0^2,m_1^2)+(p^2+m_0^2m_1^2)B_1(p^2,m_0^2,m_1^2)`$ (B.79)
$`+m_0^2+m_1^2{\displaystyle \frac{p^2}{3}}\}`$
and we can extract the divergent part
$`B_{00}(p^2,m_0^2,m_1^2)=\left[{\displaystyle \frac{1}{4}}(m_1^2+m_0^2){\displaystyle \frac{1}{12}}p^2\right]\left({\displaystyle \frac{2}{4D}}\gamma _E+\mathrm{ln}4\pi \right)+\mathrm{}.`$ (B.80)
The contribution of the top-quark loop to the self-energy of the $`Z`$ gauge boson is
$`\mathrm{\Sigma }_T^Z(k^2)`$ $`=`$ $`{\displaystyle \frac{\alpha }{4\pi }}[{\displaystyle \frac{4}{3}}({\displaystyle \frac{1}{8}}+{\displaystyle \frac{4}{9}}s_W^4{\displaystyle \frac{2}{3}}s_W^2)(k^2\mathrm{\Delta }_t+(k^2+2m_t^2)F(k^2,m_t,m_t){\displaystyle \frac{k^2}{3}})`$ (B.82)
$`{\displaystyle \frac{3}{8s_W^2c_W^2}}m_t^2(\mathrm{\Delta }_t+F(k^2,m_t,m_t))]`$
where
$`F(k^2,m_t^2,m_t^2)`$ $`=`$ $`{\displaystyle _0^1}𝑑x\mathrm{ln}{\displaystyle \frac{x^2k^2xk^2+m_t^2iϵ}{m_t^2}}.`$ (B.83)
For $`k^2<m_t^2`$ $`F(k^2,m_t^2,m_t^2)`$ can be expanded as a convergent power series in $`\frac{k^2}{m_t^2}`$ with $`\frac{k^2}{6m_t^2}`$ as the first term. Therefore, $`m_t^2F(k^2,m_t^2,m_t^2)`$ is not large in this region. For $`k^2>m_t^2`$ the term $`\frac{\mathrm{\Sigma }_t^Z(k^2)\mathrm{\Sigma }_t^Z(M_Z^2)}{k^2M_Z^2}`$ is not enhanced either. Therefore $`\frac{\mathrm{\Sigma }_t^Z(k^2)\mathrm{\Sigma }_t^Z(M_Z^2)}{k^2M_Z^2}`$ is not large for all values of $`k^2`$.
## Figure Captions
* The tree-level Feynman diagram for the process $`e^+e^{}\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}`$.
* Self-energies and counter terms for internal gauge bosons. Note that each graph represents four different combinations.
* Feynman diagrams for gauge boson mass corrections.
* Feynman diagrams for self energies of squarks and their mixing at one loop order.
* Corrections to vertex $`\gamma (Z)\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}`$ due to Higgs boson, neutralino and chargino exchanges.
* The cross section $`\sigma (e^+e^{}\stackrel{~}{t}_i\overline{\stackrel{~}{t}_i}`$,$`\stackrel{~}{b}_i\overline{\stackrel{~}{b}}_i)`$ as a function of the collision energy $`\sqrt{s}`$ for $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$, $`\mu =m_{A^0}=m_2=400GeV`$ and (a) $`\mathrm{tan}\beta =1.5`$, (b) $`\mathrm{tan}\beta =30`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ at $`\sqrt{s}=206GeV`$ for $`m_{\stackrel{~}{t}_1}=92GeV`$, $`m_{\stackrel{~}{t}_2}=350GeV`$, $`\mathrm{tan}\beta =1.5`$ and several values of $`m_{A^0}`$ and $`m_2`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`m_{A^0}`$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ at $`\sqrt{s}=206GeV`$ for $`m_{\stackrel{~}{t}_1}=92GeV`$, $`m_{\stackrel{~}{t}_2}=350GeV`$, $`\mu =600GeV`$, $`m_2=600GeV`$ and different $`\mathrm{tan}\beta `$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ at $`\sqrt{s}=500GeV`$ for $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$, $`\mathrm{tan}\beta =1.5`$, $`m_2=600GeV`$ and several values of $`m_{A^0}`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ at $`\sqrt{s}=500GeV`$ for $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$, $`\mathrm{tan}\beta =0.6`$, $`m_2=600GeV`$ and several values of $`m_{A^0}`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`m_{A^0}`$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1}`$ at $`\sqrt{s}=500GeV`$ for $`m_{\stackrel{~}{t}_1}=150GeV`$, $`m_{\stackrel{~}{t}_2}=450GeV`$, $`\mu =400GeV`$, $`m_2=600GeV`$ and several values of $`\mathrm{tan}\beta `$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1},\stackrel{~}{t}_2\overline{\stackrel{~}{t}_2}`$ at $`\sqrt{s}=2000GeV`$ for $`m_{\stackrel{~}{t}_1}=400GeV`$, $`m_{\stackrel{~}{t}_2}=800GeV`$, $`\mathrm{tan}\beta =2`$, $`m_2=1000GeV`$ and $`m_{A^0}=500GeV,900GeV,1300GeV`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`m_{A^0}`$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}_1},\stackrel{~}{t}_2\overline{\stackrel{~}{t}_2}`$ at $`\sqrt{s}=2000GeV`$ for $`m_{\stackrel{~}{t}_1}=400GeV`$, $`m_{\stackrel{~}{t}_2}=800GeV`$, $`\mu =1200GeV`$, $`m_2=1000GeV`$ and $`\mathrm{tan}\beta =1.5,2,30`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{b}_1\overline{\stackrel{~}{b}}_1,\stackrel{~}{b}_2\overline{\stackrel{~}{b}}_2`$ at $`\sqrt{s}=2000GeV`$ for $`m_{\stackrel{~}{t}_1}=400GeV`$, $`m_{\stackrel{~}{t}_2}=800GeV`$, $`m_2=1000GeV`$, $`m_{A^0}=500GeV,900GeV,1300GeV`$ and (a) $`\mathrm{tan}\beta =2`$, (b) $`\mathrm{tan}\beta =30`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`\mu `$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}}_1`$ due to the three classes of contributions, that is coming from, (1) corrections to the vertex and squark lines by Higgs bosons (2) by charginos and neutralinos and (3) corrections to the gauge boson propagators respectively. The parameters are taken to be $`\sqrt{s}=1500GeV`$, $`\mathrm{tan}\beta =3`$, $`m_{A^0}=400GeV`$, $`m_2=800GeV`$, $`m_{\stackrel{~}{t_1}}=300GeV`$ and $`m_{\stackrel{~}{t_2}}=500GeV,800GeV`$.
* Corrections $`\delta \sigma /\sigma `$ as a function of $`m_{\stackrel{~}{t_1}}`$ for $`e^+e^{}\stackrel{~}{t}_1\overline{\stackrel{~}{t}}_1`$ due to the three classes of contributions, that is coming from, (1) corrections to the vertex and squark lines by Higgs bosons (2) by charginos and neutralinos and (3) corrections to the gauge boson propagators respectively. The parameters are taken to be $`\sqrt{s}=1500GeV`$, $`\mathrm{tan}\beta =3`$, $`m_{A^0}=400GeV`$, $`m_2=800GeV`$, $`\mu =300GeV`$ and $`m_{\stackrel{~}{t_2}}=800GeV`$. |
warning/0001/gr-qc0001092.html | ar5iv | text | # 1 Introduction
## 1 Introduction
### 1.1 Conventions
We will use spacetime definitions and conventions from . In particular this means that the metric $`g_{ab}`$ is assumed to have signature $`(+)`$. We will use spinors for our calculations, but as all results are local in nature there is no need to postulate the existence of a global spinor structure on spacetime. All spinor dyads $`(o^A,\iota ^A)`$ will be assumed to be normalized i.e., $`o_A\iota ^A=1`$. $`_{AA^{}}`$ denotes the Levi-Civita connection i.e., the uniquely defined metric and torsion-free (symmetric) connection on spacetime.
### 1.2 Outline
When looking for exact solutions of Einstein’s field equations in NP- or GHP-formalism, a common approach is to assume the existence of a geodesic shear-free expanding null congruence. In vacuum spacetimes, the Goldberg-Sachs theorem tells us that this condition is necessary and sufficient for the Weyl curvature spinor to be algebraically special. In this paper we will relax the vacuum condition and instead assume that the GHP-components of the Ricci spinor satisfy
$$\mathrm{\Phi }_{00}=\mathrm{\Phi }_{01}=\mathrm{\Phi }_{02}=0$$
or equivalently
$$\mathrm{\Phi }_{ABA^{}B^{}}o^A^{}o^B^{}=0$$
where $`l^a=o^Ao^A^{}`$ denotes the geodesic shear-free expanding congruence. This assumption is still sufficient to ensure algebraic speciality of the Weyl spinor. For technical reasons we will also assume
$$\mathrm{\Lambda }=\mathrm{constant}.$$
We then show that all GHP equations containing $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ can be integrated explicitly, something which is done in for the vacuum case.<sup>1</sup><sup>1</sup>1Kerr has also integrated the GHP equations using this method for the case $`\mathrm{\Phi }_{ABA^{}B^{}}=0`$ i.e., Einstein spacetimes. We also show that the remaining GHP equations take on a very simple appearance after the integration procedure is completed and discuss redundancy in the final system of equations.
It is worth noting that a similar integration in NP-formalism for a non-vacuum subclass of these spacetimes is performed in and and, under even more restrictive assumptions, this NP integration have been performed in e.g., , , and
## 2 Spacetimes admitting a geodesic shear-free expanding null congruence
### 2.1 The general case
We will consider spacetimes admitting a geodesic, shear-free null congruence $`l^a=o^Ao^A^{}`$ and whose Ricci spinor satisfies the condition
$$\mathrm{\Phi }_{ABA^{}B^{}}o^Ao^B=0.$$
(1)
Take $`o^A`$ as the first spinor of a spinor dyad. In GHP-formalism the above conditions are equivalent to
$$\mathrm{\Phi }_{00}=\mathrm{\Phi }_{01}=\mathrm{\Phi }_{02}=0,\kappa =\sigma =0$$
(2)
By the Goldberg-Sachs theorem we obtain
$$\mathrm{\Psi }_0=\mathrm{\Psi }_1=0$$
(3)
so the spacetime is algebraically special.
In addition we will assume that $`\rho 0`$ so that $`l^a`$ is expanding. Then we can use a null rotation about $`o^A`$ to achieve $`\tau =0`$, and the Ricci equations then imply that also
$$\tau ^{}=\sigma ^{}=0.$$
(4)
We introduce Held’s modified operators which can be written
$$\stackrel{~}{\stackrel{ ̵}{}}=\frac{1}{\overline{\rho }}\stackrel{ ̵}{},\stackrel{~}{\stackrel{ ̵}{}}^{}=\frac{1}{\rho }\stackrel{ ̵}{}^{},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}+\frac{p}{2\rho }(\mathrm{\Psi }_2+2\mathrm{\Lambda })+\frac{q}{2\overline{\rho }}(\overline{\mathrm{\Psi }}_2+2\mathrm{\Lambda })$$
(5)
in this dyad. Note that the definition of $`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}`$ is slightly modified from Held’s in our non-vacuum case. The purpose of using Held’s modified operators is simply to reduce the length of calculations; in particular the new operators have the nice properties
$$\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}\right]=\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]=0$$
(6)
and
$$[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}]\eta =\left[\frac{1}{2\rho }(\mathrm{\Psi }_2+2\mathrm{\Lambda })\frac{1}{2\overline{\rho }}(\overline{\mathrm{\Psi }}_2+2\mathrm{\Lambda })\right]\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\eta $$
(7)
so if $`\eta ^{}`$ satisfies $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\eta ^{}=0`$ (a degree sign will, throughout the paper, be used to denote a quantity that is killed by $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$) then
$$\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\stackrel{~}{\stackrel{ ̵}{}}^{}\eta ^{}=\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]\eta ^{}=0$$
and the same result is true if $`\stackrel{~}{\stackrel{ ̵}{}}^{}`$ is replaced with $`\stackrel{~}{\stackrel{ ̵}{}}`$ or $`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}`$.
From the Ricci equations we obtain the equations
$$\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\rho =\rho ^2,\stackrel{~}{\stackrel{ ̵}{}}\rho =0$$
(8)
The first of these equations can be used to ‘$`\rho `$-integrate’ some of the Ricci- and Bianchi equations. However, we first need to calculate $`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho `$. Therefore, define the twist of the congruence $`l^a`$ as
$$\mathrm{\Omega }^{}=\frac{1}{\overline{\rho }}\frac{1}{\rho }$$
(9)
We then obtain
$$\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Omega }^{}=\frac{\overline{\rho }^2}{\overline{\rho }^2}+\frac{\rho ^2}{\rho ^2}=0$$
so the notation is consistent. We also obtain
$$\stackrel{~}{\stackrel{ ̵}{}}^{}\rho =\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}.$$
(10)
The Ricci- and Bianchi equations involving $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$, but not $`\mathrm{\Phi }_{22}`$ are
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\rho ^{}`$ $`=`$ $`\overline{\rho }\rho ^{}\mathrm{\Psi }_22\mathrm{\Lambda }`$ (19)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_3\mathrm{\Phi }_{21}`$ (28)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}(\mathrm{\Psi }_2+2\mathrm{\Lambda })`$ $`=`$ $`3\rho \mathrm{\Psi }_2+2\rho \mathrm{\Phi }_{11}`$ (37)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_3\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{21}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}(\mathrm{\Psi }_2+2\mathrm{\Lambda })`$ $`=`$ $`2\rho \mathrm{\Psi }_32\overline{\rho }\mathrm{\Phi }_{21}`$ (54)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_4\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{21}`$ $`=`$ $`\rho \mathrm{\Psi }_4`$ (63)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}(\mathrm{\Phi }_{11}+3\mathrm{\Lambda })`$ $`=`$ $`2(\rho +\overline{\rho })\mathrm{\Phi }_{11}`$ (72)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{21}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}(\mathrm{\Phi }_{11}3\mathrm{\Lambda })`$ $`=`$ $`(\rho +2\overline{\rho })\mathrm{\Phi }_{21}`$ (81)
In order to integrate these with respect to $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ we need to know the dependence of the Ricci scalar $`\mathrm{\Lambda }`$ on $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$. Therefore we will make the additional assumption that $`\mathrm{\Lambda }`$ is constant<sup>2</sup><sup>2</sup>2It would in fact be sufficient to assume that $`\mathrm{\Lambda }`$ is killed by $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$, $`\stackrel{ ̵}{}`$ and $`\stackrel{ ̵}{}^{}`$, but unless $`\rho =\overline{\rho }`$ the commutators would imply that also $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Lambda }=0`$, so we choose to ignore this possibility. It may also be sufficient to assume some other explicit $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$, $`\stackrel{ ̵}{}`$ and $`\stackrel{ ̵}{}^{}`$dependence of $`\mathrm{\Lambda }`$, $`\mathrm{\Psi }_2`$, or $`\mathrm{\Phi }_{11}`$.. Then these equations reduce to
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\rho ^{}`$ $`=`$ $`\overline{\rho }\rho ^{}\mathrm{\Psi }_22\mathrm{\Lambda }`$ (90)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_3\mathrm{\Phi }_{21}`$ (99)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_2`$ $`=`$ $`3\rho \mathrm{\Psi }_2+2\rho \mathrm{\Phi }_{11}`$ (108)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_3\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{21}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2`$ $`=`$ $`2\rho \mathrm{\Psi }_32\overline{\rho }\mathrm{\Phi }_{21}`$ (125)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_4\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{21}`$ $`=`$ $`\rho \mathrm{\Psi }_4`$ (134)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{11}`$ $`=`$ $`2(\rho +\overline{\rho })\mathrm{\Phi }_{11}`$ (143)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{21}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}`$ $`=`$ $`(\rho +2\overline{\rho })\mathrm{\Phi }_{21}`$ (152)
The sixth of these equations is equivalent to
$$0=\frac{1}{\rho ^2\overline{\rho }^2}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{11}\left(\frac{2}{\rho \overline{\rho }^2}+\frac{2}{\rho ^2\overline{\rho }}\right)\mathrm{\Phi }_{11}=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left(\frac{\mathrm{\Phi }_{11}}{\rho ^2\overline{\rho }^2}\right)$$
Thus, we obtain
$$\mathrm{\Phi }_{11}=\rho ^2\overline{\rho }^2\mathrm{\Phi }_{11}^{}$$
(153)
where, as usual $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{11}^{}=0`$. Then the third equation can be written
$$0=\frac{1}{\rho ^3}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Psi }_2\frac{3}{\rho ^2}\mathrm{\Psi }_22\overline{\rho }^2\mathrm{\Phi }_{11}^{}=\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left(\frac{\mathrm{\Psi }_2}{\rho ^3}2\overline{\rho }\mathrm{\Phi }_{11}^{}\right).$$
This gives us
$$\mathrm{\Psi }_2=\rho ^3\mathrm{\Psi }_2^{}+2\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}.$$
(154)
In a similar way, the remaining of the above equations can be integrated, to give
$`\rho ^{}`$ $`=`$ $`\overline{\rho }\rho ^{}{\displaystyle \frac{1}{2}}(\rho ^2+\rho \overline{\rho })\mathrm{\Psi }_2^{}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}+{\displaystyle \frac{1}{\overline{\rho }}}\mathrm{\Lambda }`$
$`\kappa ^{}`$ $`=`$ $`\kappa ^{}\rho \mathrm{\Psi }_3^{}{\displaystyle \frac{1}{2}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Phi }_{21}`$ $`=`$ $`\rho \overline{\rho }^2\mathrm{\Phi }_{21}^{}+\rho ^2\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\rho ^3\overline{\rho }^2\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_3`$ $`=`$ $`\rho ^2\mathrm{\Psi }_3^{}+\rho ^3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{3}{2}}\rho ^4\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\rho ^2\overline{\rho }\mathrm{\Phi }_{21}^{}+2\rho ^3\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+3\rho ^4\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_4`$ $`=`$ $`\rho \mathrm{\Psi }_4^{}+\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Psi }_2^{}+2\mathrm{\Psi }_3^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{1}{2}}\rho ^4\left(\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\right)`$ (155)
$`+{\displaystyle \frac{3}{2}}\rho ^5\mathrm{\Psi }_2^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{21}^{}+\rho ^3\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Phi }_{11}^{}+\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`+\rho ^4\overline{\rho }\left(\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\right)+3\rho ^5\overline{\rho }\mathrm{\Phi }_{11}^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2.`$
The remaining Ricci- and Bianchi equations are
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\rho `$ $`=`$ $`\rho \overline{\rho }^{}\mathrm{\Psi }_22\mathrm{\Lambda }`$ (164)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\rho ^{}\stackrel{ ̵}{}\kappa ^{}`$ $`=`$ $`\rho ^2+\mathrm{\Phi }_{22}`$ (173)
$`\stackrel{ ̵}{}^{}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_4`$
$`\stackrel{ ̵}{}^{}\rho ^{}`$ $`=`$ $`(\overline{\rho }\rho )\kappa ^{}\mathrm{\Psi }_3+\mathrm{\Phi }_{21}`$
$`\stackrel{ ̵}{}\mathrm{\Psi }_2`$ $`=`$ $`2\rho \mathrm{\Phi }_{12}`$
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Psi }_2\stackrel{ ̵}{}\mathrm{\Psi }_3\stackrel{ ̵}{}\mathrm{\Phi }_{21}+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{22}`$ $`=`$ $`3\rho ^{}\mathrm{\Psi }_2+\overline{\rho }\mathrm{\Phi }_{22}+2\rho ^{}\mathrm{\Phi }_{11}`$ (190)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Psi }_3\stackrel{ ̵}{}\mathrm{\Psi }_4\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Phi }_{21}+\stackrel{ ̵}{}^{}\mathrm{\Phi }_{22}`$ $`=`$ $`4\rho ^{}\mathrm{\Psi }_33\kappa ^{}\mathrm{\Psi }_22\overline{\rho }^{}\mathrm{\Phi }_{21}+2\kappa ^{}\mathrm{\Phi }_{11}`$ (207)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{22}+\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Phi }_{11}\stackrel{ ̵}{}^{}\mathrm{\Phi }_{12}\stackrel{ ̵}{}\mathrm{\Phi }_{21}`$ $`=`$ $`(\rho +\overline{\rho })\mathrm{\Phi }_{22}+2(\rho ^{}+\overline{\rho }^{})\mathrm{\Phi }_{11}`$ (224)
We will now use the expressions (155), to rewrite these equations. Starting with the first one, we obtain
$$\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho =\rho ^2\overline{\rho }^{}\frac{1}{2}\rho ^2\overline{\rho }\overline{\mathrm{\Psi }}_2^{}\frac{1}{2}\rho ^3\mathrm{\Psi }_2^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}+\frac{\rho }{\overline{\rho }}\mathrm{\Lambda }.$$
(225)
From this equation we obtain the relation
$$\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Omega }^{}=\overline{\rho }^{}\rho ^{},$$
and by applying the $`[\stackrel{ ̵}{},\stackrel{ ̵}{}^{}]`$-commutator to $`\rho `$ we obtain the useful relation
$$\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=2\mathrm{\Omega }^{}\overline{\rho }^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{}.$$
Similarly, the substitution of known expressions into the third, fourth and fifth equation of (224) easily give us
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_4^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho ^{}`$ $`=`$ $`\mathrm{\Omega }^{}\kappa ^{}\mathrm{\Psi }_3^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_2^{}`$ $`=`$ $`2\overline{\mathrm{\Phi }}_{21}^{}`$ (226)
As we have already calculated the derivative operators action on $`\mathrm{\Omega }^{}`$ we can replace their action on $`\rho `$ with their action on the real $`(1,1)`$-quantity $`\frac{1}{\rho }+\frac{1}{\overline{\rho }}`$
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`2`$ (235)
$`\stackrel{~}{\stackrel{ ̵}{}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`(\rho ^{}+\overline{\rho }^{})+\rho \mathrm{\Psi }_2^{}+\overline{\rho }\overline{\mathrm{\Psi }}_2^{}+2\rho \overline{\rho }\mathrm{\Phi }_{11}^{}{\displaystyle \frac{2\mathrm{\Lambda }}{\rho \overline{\rho }}}.`$ (244)
We will next use the last equation of (224) to $`\rho `$-integrate for $`\mathrm{\Phi }_{22}`$, but before that we give two of the commutators that will be used in the integration
$`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$ $`=`$ $`({\displaystyle \frac{\kappa ^{}}{\rho }}+\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\rho ^2\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\overline{\rho }\mathrm{\Phi }_{21}^{}+\rho \overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}`$ (262)
$`+\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}+p(\kappa ^{}+\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`[\stackrel{~}{\stackrel{ ̵}{}},\stackrel{~}{\stackrel{ ̵}{}}^{}]`$ $`=`$ $`\left({\displaystyle \frac{\overline{\rho }^{}}{\overline{\rho }}}{\displaystyle \frac{\rho ^{}}{\rho }}+{\displaystyle \frac{\rho }{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\mathrm{\Psi }_2^{}{\displaystyle \frac{\overline{\rho }}{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\overline{\mathrm{\Psi }}_2^{}+\mathrm{\Omega }^{}(\rho \overline{\rho }\mathrm{\Phi }_{11}^{}\mathrm{\Lambda })\right)\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ (280)
$`+\mathrm{\Omega }^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}+p\left(\rho ^{}+\mathrm{\Omega }^2\mathrm{\Lambda }\right)q(\overline{\rho }^{}+\mathrm{\Omega }^2\mathrm{\Lambda }).`$
The last equation of (224) requires considerably more work to $`\rho `$-integrate than the ones integrated previously. However, after a very long calculation involving these commutators, we obtain
$`\mathrm{\Phi }_{22}`$ $`=`$ $`\rho \overline{\rho }\mathrm{\Phi }_{22}^{}+\rho ^2\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}\overline{\mathrm{\Phi }}_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})+\rho \overline{\rho }^2(\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})`$ (298)
$`+\rho ^3\overline{\rho }\overline{\mathrm{\Phi }}_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }^2(\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{})\rho \overline{\rho }^3\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}`$
$`+\rho ^3\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{}\rho ^2\overline{\rho }^3\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }^3\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
Next we will look at the sixth equation of (224). We start by subtracting it from the last equation of (224) to get rid of the $`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Phi }_{22}`$-term
$$\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Phi }_{11}\stackrel{ ̵}{}^{}\mathrm{\Phi }_{12}\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}^{}\mathrm{\Psi }_2+\stackrel{ ̵}{}\mathrm{\Psi }_3=3\rho ^{}\mathrm{\Psi }_2+2\overline{\rho }^{}\mathrm{\Phi }_{11}+\rho \mathrm{\Phi }_{22}.$$
(299)
By using the previous equations along with the $`[\stackrel{~}{\stackrel{ ̵}{}},\stackrel{~}{\stackrel{ ̵}{}}^{}]`$-commutator it becomes
$$\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_3^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Psi }_2^{}=\mathrm{\Phi }_{22}^{}$$
(300)
It remains to check the second and seventh equation of (224). Using the equations obtained so far, the second gives us
$$\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho ^{}\stackrel{~}{\stackrel{ ̵}{}}\kappa ^{}=\mathrm{\Lambda }(2\mathrm{\Omega }^{}\rho ^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{}).$$
(301)
Checking the seventh equation of (224) is a very long and tedious calculation where both of the commutators (280) are used. The end result is that
$$\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_4^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Psi }_3^{}=\mathrm{\Lambda }(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}2\mathrm{\Phi }_{21}^{}2\mathrm{\Omega }^{}\mathrm{\Psi }_3^{})$$
(302)
As could be expected, the system of equations that we have obtained contains considerable redundancy. As an example of this we note that the imaginary part of equation (300) and the whole of equation (302) are actually consequences of the other equations.
## 3 Summary and conclusions
### 3.1 Summary
In this section we will collect the resulting equations in one place.
First of all, the equations for the GHP-operators acting on $`\rho `$
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\rho `$ $`=`$ $`\rho ^2`$ (311)
$`\stackrel{~}{\stackrel{ ̵}{}}\rho `$ $`=`$ $`0`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho `$ $`=`$ $`\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho `$ $`=`$ $`\rho ^2\overline{\rho }^{}{\displaystyle \frac{1}{2}}\rho ^2\overline{\rho }\overline{\mathrm{\Psi }}_2^{}{\displaystyle \frac{1}{2}}\rho ^3\mathrm{\Psi }_2^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}+{\displaystyle \frac{\rho }{\overline{\rho }}}\mathrm{\Lambda }`$ (320)
can be split into the relations
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`2`$ (329)
$`\stackrel{~}{\stackrel{ ̵}{}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)`$ $`=`$ $`(\rho ^{}+\overline{\rho }^{})+\rho \mathrm{\Psi }_2^{}+\overline{\rho }\overline{\mathrm{\Psi }}_2^{}+2\rho \overline{\rho }\mathrm{\Phi }_{11}^{}{\displaystyle \frac{2\mathrm{\Lambda }}{\rho \overline{\rho }}}`$ (338)
$`\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}\mathrm{\Omega }^{}`$ $`=`$ $`0`$ (347)
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Omega }^{}`$ $`=`$ $`\overline{\rho }^{}\rho ^{}`$ (356)
From the commutators on $`\rho `$ we have the useful relation
$$\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}=2\mathrm{\Omega }^{}\overline{\rho }^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{}.$$
(357)
The curvature scalars and the spin coefficients are
$`\rho ^{}`$ $`=`$ $`\overline{\rho }\rho ^{}{\displaystyle \frac{1}{2}}(\rho ^2+\rho \overline{\rho })\mathrm{\Psi }_2^{}\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}+{\displaystyle \frac{1}{\overline{\rho }}}\mathrm{\Lambda }`$
$`\kappa ^{}`$ $`=`$ $`\kappa ^{}\rho \mathrm{\Psi }_3^{}{\displaystyle \frac{1}{2}}\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}{\displaystyle \frac{1}{2}}\rho ^3\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\rho \overline{\rho }\mathrm{\Phi }_{21}^{}\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_2`$ $`=`$ $`\rho ^3\mathrm{\Psi }_2^{}+2\rho ^3\overline{\rho }\mathrm{\Phi }_{11}^{}`$
$`\mathrm{\Psi }_3`$ $`=`$ $`\rho ^2\mathrm{\Psi }_3^{}+\rho ^3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{3}{2}}\rho ^4\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\rho ^2\overline{\rho }\mathrm{\Phi }_{21}^{}+2\rho ^3\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+3\rho ^4\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Psi }_4`$ $`=`$ $`\rho \mathrm{\Psi }_4^{}+\rho ^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho ^3(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Psi }_2^{}+2\mathrm{\Psi }_3^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})+{\displaystyle \frac{1}{2}}\rho ^4\left(\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}\right)`$
$`+{\displaystyle \frac{3}{2}}\rho ^5\mathrm{\Psi }_2^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2+\rho ^2\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{21}^{}+\rho ^3\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Phi }_{11}^{}+\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`+\rho ^4\overline{\rho }\left(\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^2\mathrm{\Omega }^{}+3\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\right)+3\rho ^5\overline{\rho }\mathrm{\Phi }_{11}^{}(\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})^2`$
$`\mathrm{\Phi }_{11}`$ $`=`$ $`\rho ^2\overline{\rho }^2\mathrm{\Phi }_{11}^{}`$
$`\mathrm{\Phi }_{21}`$ $`=`$ $`\rho \overline{\rho }^2\mathrm{\Phi }_{21}^{}+\rho ^2\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\rho ^3\overline{\rho }^2\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}`$
$`\mathrm{\Phi }_{22}`$ $`=`$ $`\rho \overline{\rho }\mathrm{\Phi }_{22}^{}+\rho ^2\overline{\rho }(\stackrel{~}{\stackrel{ ̵}{}}^{}\overline{\mathrm{\Phi }}_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})+\rho \overline{\rho }^2(\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{21}^{}{\displaystyle \frac{1}{2}}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Phi }_{11}^{})`$ (375)
$`+\rho ^3\overline{\rho }\overline{\mathrm{\Phi }}_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+{\displaystyle \frac{1}{2}}(\stackrel{~}{\stackrel{ ̵}{}}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{})\rho \overline{\rho }^3\mathrm{\Phi }_{21}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}`$
$`+\rho ^3\overline{\rho }^2\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{}\rho ^2\overline{\rho }^3\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}\rho ^3\overline{\rho }^3\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}.`$
The remaining Ricci and Bianchi equations are
$`\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\rho ^{}\stackrel{~}{\stackrel{ ̵}{}}\kappa ^{}`$ $`=`$ $`(2\mathrm{\Omega }^{}\rho ^{}+\mathrm{\Psi }_2^{}\overline{\mathrm{\Psi }}_2^{})\mathrm{\Lambda }`$ (384)
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\kappa ^{}`$ $`=`$ $`\mathrm{\Psi }_4^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}^{}\rho ^{}`$ $`=`$ $`\mathrm{\Omega }^{}\kappa ^{}\mathrm{\Psi }_3^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_2^{}`$ $`=`$ $`2\overline{\mathrm{\Phi }}_{21}^{}`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_3^{}+\stackrel{~}{\stackrel{ ̵}{}}^{}\overline{\mathrm{\Psi }}_3^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}(\mathrm{\Psi }_2^{}+\overline{\mathrm{\Psi }}_2^{})`$ $`=`$ $`2\mathrm{\Phi }_{22}^{}.`$ (393)
The commutators become
$`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}\right]`$ $`=`$ $`0`$ (402)
$`\left[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$ $`=`$ $`0`$ (411)
$`[\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array},\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}]`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}\rho ^2\mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\overline{\rho }^2\overline{\mathrm{\Psi }}_2^{}+\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}+\rho \overline{\rho }^2\mathrm{\Phi }_{11}^{}+\mathrm{\Lambda }\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\right)\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ (436)
$`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}\right]`$ $`=`$ $`({\displaystyle \frac{\overline{\kappa }^{}}{\overline{\rho }}}+\overline{\mathrm{\Psi }}_3^{}+{\displaystyle \frac{1}{2}}\overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\overline{\mathrm{\Psi }}_2^{}{\displaystyle \frac{1}{2}}\overline{\rho }^2\overline{\mathrm{\Psi }}_2^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{}+\rho \overline{\mathrm{\Phi }}_{21}^{}+\rho \overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Phi }_{11}^{}`$ (454)
$`\rho \overline{\rho }^2\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{})\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}+q(\overline{\kappa }^{}\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Omega }^{})`$
$`\left[\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{},\stackrel{~}{\stackrel{ ̵}{}}^{}\right]`$ $`=`$ $`({\displaystyle \frac{\kappa ^{}}{\rho }}+\mathrm{\Psi }_3^{}+{\displaystyle \frac{1}{2}}\rho \stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Psi }_2^{}+{\displaystyle \frac{1}{2}}\rho ^2\mathrm{\Psi }_2^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{}+\overline{\rho }\mathrm{\Phi }_{21}^{}+\rho \overline{\rho }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Phi }_{11}^{}`$ (472)
$`+\rho ^2\overline{\rho }\mathrm{\Phi }_{11}^{}\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}+p(\kappa ^{}+\mathrm{\Lambda }\stackrel{~}{\stackrel{ ̵}{}}^{}\mathrm{\Omega }^{})`$
$`[\stackrel{~}{\stackrel{ ̵}{}},\stackrel{~}{\stackrel{ ̵}{}}^{}]`$ $`=`$ $`\left({\displaystyle \frac{\overline{\rho }^{}}{\overline{\rho }}}{\displaystyle \frac{\rho ^{}}{\rho }}+{\displaystyle \frac{\rho }{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\mathrm{\Psi }_2^{}{\displaystyle \frac{\overline{\rho }}{2}}\left({\displaystyle \frac{1}{\rho }}+{\displaystyle \frac{1}{\overline{\rho }}}\right)\overline{\mathrm{\Psi }}_2^{}+\mathrm{\Omega }^{}(\rho \overline{\rho }\mathrm{\Phi }_{11}^{}\mathrm{\Lambda })\right)\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}`$ (490)
$`+\mathrm{\Omega }^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}+p\left(\rho ^{}+\mathrm{\Omega }^2\mathrm{\Lambda }\right)q(\overline{\rho }^{}+\mathrm{\Omega }^2\mathrm{\Lambda })`$
### 3.2 Conclusions
The most obvious application of these results is when searching for exact non-vacuum solutions of Einstein’s field equations using the GHP integration procedure suggested by Held , and developed by Edgar and Ludwig , , . It is worth noting that among the remaining Ricci- and Bianchi equations (393) we can view the second one as the definition of $`\mathrm{\Psi }_4^{}`$, the third as the definition of $`\mathrm{\Psi }_3^{}`$, the fourth as the definition of $`\mathrm{\Phi }_{21}^{}`$ and the fifth as the definition of $`\mathrm{\Phi }_{22}^{}`$. This leaves only a system of equations consisting of the equations (356), (490) and the first equation of (393) for the unknown functions $`\frac{1}{\rho }+\frac{1}{\overline{\rho }}`$, $`\mathrm{\Omega }^{}`$, $`\kappa ^{}`$, $`\rho ^{}`$, $`\mathrm{\Psi }_2^{}`$, $`\mathrm{\Phi }_{11}^{}`$ and the constant $`\mathrm{\Lambda }`$. Of these, $`\frac{1}{\rho }+\frac{1}{\overline{\rho }}`$, $`\mathrm{\Phi }_{11}^{}`$ and $`\mathrm{\Lambda }`$ are real while $`\mathrm{\Omega }^{}`$ is purely imaginary and the others are complex.
There are however two important points to keep in mind.
* The GHP integration program requires that the commutators are applied to four real $`(0,0)`$-weighted quantities and one non-trivially weighted complex quantity. Only then is all the information extracted from the commutators.
* We have not yet imposed Einstein’s field equations. When we do this some of the equations taken as definitions of the Ricci components will actually become constraints e.g., if we impose the condition that the spacetime be vacuum we obtain the extra conditions
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_2^{}`$ $`=`$ $`0`$
$`\stackrel{~}{\stackrel{ ̵}{}}\mathrm{\Psi }_3^{}\stackrel{~}{\begin{array}{c}\mathrm{I}\hfill \\ \mathrm{P}\hfill \end{array}}^{}\mathrm{\Psi }_2^{}`$ $`=`$ $`0.`$
Also, impositions of other features such as symmetries, special Petrov types etc., may result in extra constraint equations.
Another important application is to quasi-local momentum in these spacetimes. Recall that a symmetric spinor $`L_{ABCA^{}}`$ is said to be a Lanczos potential of $`\mathrm{\Psi }_{ABCD}`$ if it satisfies the Weyl-Lanczos equation i.e.,
$$\mathrm{\Psi }_{ABCD}=2_{(A}{}_{}{}^{A^{}}L_{BCD)A^{}}^{}$$
This equation can be translated into GHP-formalism (see e.g., ) and if we look for solutions satisfying
$$L_{ABCA^{}}o^A^{}=0$$
the GHP Weyl-Lanczos equations are easy to $`\rho `$-integrate. Such Lanczos potentials turn out to be very useful when looking for metric asymmetric curvature-free connections in these spacetimes. Such a connection have been used in the Kerr spacetime to construct quasi-local momentum. For details of this construction and some generalizations, see , , , and . This application is further developed in .
## Acknowledgements
Special thanks are due to docent S. Brian Edgar for helpful suggestions and discussions. |
warning/0001/hep-ph0001088.html | ar5iv | text | # DO–TH 00/01LNF-00/001 (P) Analyzing 𝜀'/𝜀 in the 1/𝑁_𝑐 Expansion
## 1 Introduction
Direct CP violation in $`K\pi \pi `$ decays was recently observed by the KTeV and NA48 collaborations. The present world average for the parameter $`\epsilon ^{}/\epsilon `$ is $`\text{Re}\epsilon ^{}/\epsilon =(21.2\pm 4.6)10^4`$. In the standard model CP violation originates in the CKM phase, and direct CP violation is governed by loop diagrams of the penguin type. The main source of uncertainty in the calculation of $`\epsilon ^{}/\epsilon `$ is the QCD non-perturbative contribution related to the hadronic nature of the $`K\pi \pi `$ decay. Using the $`\mathrm{\Delta }S=1`$ effective hamiltonian,
$$_{eff}^{\mathrm{\Delta }S=1}=\frac{G_F}{\sqrt{2}}\lambda _u\underset{i=1}{\overset{8}{}}c_i(\mu )Q_i(\mu )(\mu <m_c),$$
(1)
the non-perturbative contribution, contained in the hadronic matrix elements of the four-quark operators $`Q_i`$, can be separated from the perturbative Wilson coefficients $`c_i(\mu )=z_i(\mu )+\tau y_i(\mu )`$ (with $`\tau =\lambda _t/\lambda _u`$ and $`\lambda _q=V_{qs}^{}V_{qd}`$). Introducing $`Q_i_I(\pi \pi )_I|Q_i|K`$, the CP ratio can be written as
$$\frac{\epsilon ^{}}{\epsilon }=\frac{G_F}{2}\frac{\omega \text{ Im}\lambda _t}{|\epsilon |\text{Re}A_0}\left[\left|\underset{i}{}y_iQ_i_0\right|\left(1\mathrm{\Omega }_{\eta +\eta ^{}}\right)\frac{1}{\omega }\left|\underset{i}{}y_iQ_i_2\right|\right].$$
(2)
$`\omega =`$Re$`A_0/`$Re$`A_2=22.2`$ is the ratio of the CP conserving $`K\pi \pi `$ isospin amplitudes; $`\mathrm{\Omega }_{\eta +\eta ^{}}`$ encodes the effect of the isospin breaking in the quark masses. $`\epsilon ^{}/\epsilon `$ is dominated by $`Q_6_0`$ and $`Q_8_2`$ which cannot be fixed from the CP conserving data. Beside the theoretical uncertainties coming from the calculation of the $`Q_i_I`$ and of $`\mathrm{\Omega }_{\eta +\eta ^{}}`$, the analysis of the CP ratio suffers from the uncertainties on the values of various input parameters, in particular of the CKM phase in Im$`\lambda _t`$, of $`\mathrm{\Lambda }_{\text{QCD}}\mathrm{\Lambda }_{\overline{\text{MS}}}^{(4)}`$, and of the strange quark mass.
To calculate the hadronic matrix elements we start from the effective chiral lagrangian for pseudoscalar mesons which involves an expansion in momenta where terms up to $`𝒪(p^4)`$ are included. The method we use is the $`1/N_c`$ expansion. In this approach, we expand the matrix elements in powers of the momenta and of $`1/N_c`$. For the $`1/N_c`$ corrections we calculated chiral loops as described in refs. . Especially important to this analysis are the non-factorizable corrections, which are UV divergent and must be matched to the short-distance part. They are regularized by a finite cutoff $`\mathrm{\Lambda }_c`$ which is identified with the short-distance renormalization scale. The definition of the momenta in the loop diagrams, which are not momentum translation invariant, is discussed in detail in ref. . Other recent work on matrix elements in the $`1/N_c`$ approach can be found in refs. .
For the Wilson coefficients we use the leading logarithmic and the next-to-leading logarithmic values. The absence of any reference to the renormalization scheme in the low-energy calculation, at this stage, prevents a complete matching at the next-to-leading order. Nevertheless, a comparison of the numerical results obtained from the LO and NLO coefficients is useful as regards estimating the uncertainties and testing the validity of perturbation theory.
## 2 Analysis of $`𝜺^{\mathbf{}}\mathbf{/}𝜺`$
Analytical formulas for all matrix elements, at next-to-leading order in the twofold expansion in powers of momenta and of $`1/N_c`$, are given in refs. . In the pseudoscalar approximation, the matching has to be done below 1 GeV. Varying $`\mathrm{\Lambda }_c`$ between 600 and 900 MeV, the bag factors $`B_1^{(1/2)}`$ and $`B_2^{(1/2)}`$ take the values $`8.214.2`$ and $`2.94.6`$; quadratic terms in $`Q_1_0`$ and $`Q_2_0`$ produce a large enhancement which brings the $`\mathrm{\Delta }I=1/2`$ amplitude in agreement with the data. Corrections beyond the chiral limit were found to be small.
For $`Q_6_0`$ and $`Q_8_2`$ the leading non-factorizable loop corrections, which are of $`𝒪(p^0/N_c)`$, are only logarithmically divergent. Including terms of $`𝒪(p^0)`$, $`𝒪(p^2)`$, and $`𝒪(p^0/N_c)`$, $`B_6^{(1/2)}`$ and $`B_8^{(3/2)}`$ take the values $`1.100.72`$ and $`0.640.42`$. As a result the experimental range for $`\epsilon ^{}/\epsilon `$ can be accommodated in the standard model only if there is a conspiracy of the input parameters.<sup>2</sup><sup>2</sup>2For supersymmetric contributions to $`\epsilon ^{}/\epsilon `$ see ref. and references therein. However, since the leading $`𝒪(p^0)`$ contribution vanishes for $`Q_6`$, corrections from higher order terms beyond the $`𝒪(p^2)`$ and $`𝒪(p^0/N_c)`$ are expected to be large. In ref. we investigated the $`𝒪(p^2/N_c)`$ contribution, i.e., the $`1/N_c`$ correction at the next order in the chiral expansion, because it brings about, for the first time, quadratic corrections on the cutoff. From counting arguments and more generally from the fact that the chiral limit is assumed to be reliable, the quadratic terms (which are not chirally suppressed) are expected to be dominant. It is still desirable to check that explicitly by calculating the corrections beyond the chiral limit, from logarithms and finite terms, as done for $`Q_1`$ and $`Q_2`$. Numerically, we observe a large positive correction from the quadratic term in $`Q_6_0`$. This point was already emphasized in ref. . The slope of the correction is qualitatively consistent and welcome since it compensates for the logarithmic decrease at $`𝒪(p^0/N_c)`$. Varying $`\mathrm{\Lambda }_c`$ between 600 and 900 MeV, the $`B_6^{(1/2)}`$ factor takes the values $`1.501.62`$. $`Q_6`$ is a $`\mathrm{\Delta }I=1/2`$ operator, and the enhancement of $`Q_6_0`$ indicates that at the level of the $`1/N_c`$ corrections the dynamics of the $`\mathrm{\Delta }I=1/2`$ rule applies to $`Q_6`$ as to $`Q_1`$ and $`Q_2`$.
Using the quoted values for $`B_6^{(1/2)}`$ together with the full leading plus next-to-leading order $`B`$ factors for the remaining operators we calculated $`\epsilon ^{}/\epsilon `$. The results for the three sets of Wilson coefficients LO, NDR, and HV and for $`\mathrm{\Lambda }_c`$ between 600 and $`900\text{MeV}`$ are given in Tab. 1. The numbers are close to the measured value for central values of the parameters (first column). They are obtained by assuming zero phases from final state interactions. This approximation is very close to the results we would get if we used the small imaginary part obtained at the one-loop level.
Performing a scanning of the parameters \[$`125\text{MeV}m_s(1\text{GeV})175`$ MeV, $`0.15\mathrm{\Omega }_{\eta +\eta ^{}}0.35`$, $`1.0410^4\text{Im}\lambda _t1.6310^4`$, and $`245\text{MeV}\mathrm{\Lambda }_{\text{QCD}}405\text{MeV}`$\] we obtain the numbers in the second column of Tab. 1. They can be compared with the results of refs. . The values of $`B_6^{(1/2)}`$ can also be compared with ref. and those of $`B_8^{(3/2)}`$ with refs. . The large ranges reported in the table can be traced back to the large ranges of the input parameters. This can be seen by comparing them with the relatively narrow ranges obtained for central values of the parameters. The parameters, to a large extent, act multiplicatively, and the large range for $`\epsilon ^{}/\epsilon `$ is due to the fact that the central value(s) for the ratio are enhanced roughly by a factor of two compared to the results obtained with $`B`$ factors for $`Q_6`$ and $`Q_8`$ close to the VSA. More accurate information on the parameters, from theory and experiment, will restrict the values for $`\epsilon ^{}/\epsilon `$.
To estimate the uncertainties due to higher order final state interactions we also calculated $`\epsilon ^{}/\epsilon `$ using the real part of the matrix elements and the phenomenological values of the phases , $`\delta _0=(34.2\pm 2.2)^{}`$ and $`\delta _2=(6.9\pm 0.2)^{}`$, i.e., we replaced $`|_iy_iQ_i_I|`$ in Eq. (2) by $`_iy_i\text{Re}Q_i_I/\mathrm{cos}\delta _I`$. The corresponding results are given in Tab. 2. They are enhanced by $`25\%`$ compared to the numbers in Tab. 1. We would like to emphasize that this $`25\%`$ uncertainty should be taken into account by any analysis which either does not include final state interactions or cannot reproduce the numerical values of the phases. To reduce the uncertainties in the $`1/N_c`$ approach it would be interesting to investigate the two-loop imaginary part and/or to combine our calculation with a dispersive calculation along the lines of refs. . In order to reduce the scheme dependence in the result, appropriate subtractions would be necessary (see refs. ). Finally, it is reasonable to assume that the effect of the pseudoscalar mesons is the most important one. Nevertheless, incorporating the vector mesons and higher resonances would be desirable in order to improve the treatment of the intermediate region around the rho mass and to show explicitly that the large enhancement we find at low energy in the treatment of the pseudoscalars remains valid up to the scale $`m_c`$, where the matching with the short-distance part can be done more safely.
## Acknowledgments
This work was supported by BMBF, 057D093P(7), Bonn, FRG, and DFG Antrag PA-10-1. T.H. acknowledges support from EEC, TMR-CT980169.
## References |
warning/0001/cond-mat0001004.html | ar5iv | text | # Characterization of ferrimagnetic Heisenberg chains according to the constituent spins
## I Introduction
One-dimensional quantum ferrimagnets are one of the hot topics and recent progress in the theoretical understanding of them deserves special mention. Such a vigorous argument a great deal originates in the pioneering efforts to synthesize binuclear magnetic materials including one-dimensional systems. The first ferrimagnetic chain compound , MnCu(dto)<sub>2</sub>(H<sub>2</sub>O)<sub>3</sub>$``$$`4.5`$H<sub>2</sub>O ($`\text{dto}=\text{dithiooxalato}=\text{S}_2\text{C}_2\text{O}_2`$), was synthesized by Gleizes and Verdaguer and stimulated the public interest in this potential subject. The following example of an ordered bimetallic chain , MnCu(pba)(H<sub>2</sub>O)<sub>3</sub>$``$$`2`$H<sub>2</sub>O ($`\text{pba}=1,3\text{-propylenebis(oxamato)}=\text{C}_7\text{H}_6\text{N}_2\text{O}_6`$), exhibiting more pronounced one dimensionality, further activated the physical , as well as chemical , investigations. There also appeared an idea that the alternating magnetic centers do not need to be metal ions but may be organic radicals.
A practical model of a ferrimagnetic chain is two kinds of spins $`S`$ and $`s`$ ($`S>s`$) alternating on a ring with antiferromagnetic exchange coupling between nearest neighbors, as described by the Hamiltonian,
$$=J\underset{j=1}{\overset{N}{}}\left(𝑺_j𝒔_j+𝒔_j𝑺_{j+1}\right).$$
(1)
Here $`N`$ is the number of unit cells and we set the unit-cell length equal to $`2a`$ in the following. The simplest case, $`(S,s)=(1,\frac{1}{2})`$, has so far been discussed fairly well. There lie both ferromagnetic and antiferromagnetic long-range orders in the ground state . The ground state, which is a multiplet of spin $`(Ss)N`$, shows elementary excitations of two distinct types . The excitations of ferromagnetic aspect, reducing the ground-state magnetization, form a gapless dispersion relation, whereas those of antiferromagnetic aspect, enhancing the ground-state magnetization, are gapped from the ground state. As a result of the low-energy structure of dual aspect, the specific heat shows a Shottky-like peak in spite of the ferromagnetic low-temperature behavior and the magnetic susceptibility times temperature exhibits a round minimum . When the exchange coupling turns anisotropic, the dispersion of the ferromagnetic excitations is no more quadratic and the plateau in the ground-state magnetization curve due to the gapped antiferromagnetic excitations vanishes via the Kosterlitz-Thouless transition .
The quantum behavior of the model with higher spins is not yet so clear as that in the $`(S,s)=(1,\frac{1}{2})`$ case. Though Drillon et al. made the first attempt to reveal the general behavior of the model, they fixed the smaller spin to $`\frac{1}{2}`$ in their argument with particular emphasis on the problem of the crystal engineering of a molecule-based ferromagnet$``$the assembly of the highly magnetic molecular entities within the crystal lattice in a ferromagnetic fashion. There also exists an extensive numerical study , but the leading attention was not necessarily directed to the consequences of the variation of the constituent spins. So far the generic behavior, rather than individual features, of ferrimagnetic mixed-spin chains has been accentuated or predicted, but there are few attempts to characterize or classify the typical one-dimensional ferrimagnetic behavior as a function of $`(S,s)`$. In such circumstances, we aim in this article at elucidating which property of the model (1) is universal and which one, if any, is variable with $`(S,s)`$.
It is true that numerical tools are quite useful in this context, but they are not almighty. Although the low-temperature ferromagnetic behavior is quite interesting from the experimental point of view, it is hardly feasible to numerically take grand-canonical averages at low enough temperatures. It is also unfortunate that with the increase of $`(S,s)`$, the available information is more and more reduced in both quality and quantity. Then we have an idea of describing the model in terms of the spin-wave theory. The conventional spin-wave treatment of low-dimensional magnets may discourage us. The Haldane conjecture revealed a qualitative breakdown of the spin-wave description of one-dimensional Heisenberg antiferromagnets. Neither quantum corrections nor additional constraints to spin fluctuations end up with an overall scenario for the low-energy physics applicable to general spins. However, Ivanov has recently reported that the spin-wave description of one-dimensional Heisenberg ferrimagnets is quite successful. Though his calculations were restricted to the ground-state energy and magnetization, the highly accurate estimates obtained there are surprising enough, considering the diverging ground-state magnetization in the one-dimensional antiferromagnetic spin-wave theory. For antiferromagnets, quantum fluctuations of domain-wall type, connecting the two degenerate Néel states, are important, whereas for ferrimagnets, domain-wall excitations lead to magnetization fluctuations and are thus of less significance. Therefore it is likely that the spin-wave approach is highly efficient for ferrimagnets. In an attempt to demonstrate such an idea, we try to describe not only the low-energy structure but also the thermal behavior of the Heisenberg ferrimagnetic spin chains (1) within the framework of the spin-wave theory. Extensive numerical calculations, supplemented by the spin-wave analysis, fully reveal the one-dimensional ferrimagnetic behavior as a function of the constituent spins.
## II Elementary Excitations
In order to investigate the low-energy structure, we employ a new quantum Monte Carlo technique as well as the conventional Lanczos diagonalization algorithm. The idea is in a word expressed as extracting the low-lying excitations from imaginary-time quantum Monte Carlo data at a low enough temperature. The imaginary-time correlation function $`S(q,\tau )`$ is generally defined as
$$S(q,\tau )=\mathrm{e}^\tau O_q\mathrm{e}^\tau O_q,$$
(2)
where $`O_q=N^1_{j=1}^NO_j\mathrm{e}^{2\mathrm{i}aqj}`$ is the Fourier transform of an arbitrary local operator $`O_j`$, which may be an effective combination of the spins $`𝑺`$ and $`𝒔`$, and the thermal average at a given temperature $`\beta ^1=k_\mathrm{B}T`$, $`A\mathrm{Tr}[\mathrm{e}^\beta A]/\mathrm{Tr}[\mathrm{e}^\beta ]`$, is taken in a certain subspace. $`S(q,\tau )`$ can be represented in terms of the eigenvectors and eigenvalues of the Hamiltonian, $`|l;k`$ $`(l=1,2,\mathrm{})`$ and $`E_l(k)`$ ($`E_1(k)E_2(k)\mathrm{}`$), and behaves like
$$S(q,\tau )\underset{l}{}\left|1;k_0|S_q^z|l;k_0+q\right|^2\mathrm{e}^{\tau \left[E_l(k_0+q)E_1(k_0)\right]},$$
(3)
at a sufficiently low temperature, where $`k_0`$ is the momentum at which the lowest-energy state in the subspace is located. Therefore we can reasonably approximate $`E_1(k_0+q)E_1(k_0)`$ by the slope $`\mathrm{ln}[S(q,\tau )]/\tau `$ in the large-$`\tau `$ region. When we take interest in the lower edge of the excitation spectrum, such a treatment is rather straightforward. The elementary excitations of the Haldane antiferromagnets were indeed revealed thus . Here, due to the two kinds of spins in a chain and the resultant dual aspect of the low-energy structure, the relevant subspace and operator $`O_j`$ are not uniquely defined. Since the total magnetization, $`M=_j(S_j^z+s_j^z)`$, is a conserved quantity in the present system, we consider calculating $`S(q,\tau )`$ independently in each subspace with a given $`M`$ . The elementary excitation energies for the ferromagnetic branch are obtained from a single calculation, $`S(q,\tau )`$ in the subspace of $`M=0`$, while those for the antiferromagnetic branch from a couple of calculations, $`S(q,\tau )`$ and the lowest energy in the subspace of $`M=(Ss)N+1`$. We have taken $`S_j^z\pm s_j^z`$ for $`O_j`$ and found that $`O_j=S_j^zs_j^z`$ extracts the eigenvalues of the bonding (lower-energy) states in both subspaces. The choice of the scattering matrices is a profound problem in itself and is fully discussed elsewhere .
Although the chain length we can reach with the exact-diagonalization method is strongly limited, the diagonalization results are still helpful in the present system whose correlation length is generally so small as to be comparable to the unit-cell length (see Fig. 5 bellow). Actually, the ground-state energies for $`N=10`$ coincide with their thermodynamic-limit values within the first several digits. The Lanczos algorithm gives the most precise estimate for the ground-state energy and the antiferromagnetic excitation gap, whereas the quantum Monte Carlo technique is necessary for the evaluation of the curvature of the small-momentum dispersion.
On the other hand, we consider a spin-wave description of the elementary excitations as well. We introduce the bosonic operators for the spin deviation in each sublattice via
$$\begin{array}{ccc}S_j^+=\sqrt{2Sa_j^{}a_j}a_j,\hfill & S_j^z=Sa_j^{}a_j,\hfill & \\ s_j^+=b_j^{}\sqrt{2sb_j^{}b_j},\hfill & s_j^z=s+b_j^{}b_j,\hfill & \end{array}$$
(4)
where we regard $`S`$ and $`s`$ as quantities of the same order. The Hamiltonian (1) is expressed in terms of the bosonic operators as
$$=E_{\mathrm{class}}+_0+_1+O(S^1),$$
(5)
where $`E_{\mathrm{class}}=2sSJN`$ is the classical ground-state energy, and $`_0`$ and $`_1`$ are the one-body and two-body terms of the order $`O(S^1)`$ and $`O(S^0)`$, respectively. We may consider the simultaneous diagonalization of $`_0`$ and $`_1`$ in the naivest attempt to go beyond the noninteracting spin-wave theory. Indeed, the higher-order terms we take into account, the better description of the ground-state properties we obtain . However, such an idea, as a whole, ends in failure, bringing a gap to the lowest-lying ferromagnetic excitation branch and thus qualitatively misreading the low-energy physics. Therefore, we take an alternative approach at the idea of first diagonalizing $`_0`$ and next extracting relevant corrections from $`_1`$. $`_0`$ is diagonalized as
$$_0=E_0+J\underset{k}{}\left(\omega _k^{}\alpha _k^{}\alpha _k+\omega _k^+\beta _k^{}\beta _k\right),$$
(6)
where
$$E_0=J\underset{k}{}\left[\omega _k(S+s)\right],$$
(7)
is the $`O(S^1)`$ quantum correction to the ground-state energy, and $`\alpha _k^{}`$ and $`\beta _k^{}`$ are the creation operators of the ferromagnetic and antiferromagnetic spin waves of momentum $`k`$ whose dispersion relations are given by
$$\omega _k^\pm =\omega _k\pm (Ss),$$
(8)
with
$$\omega _k=\sqrt{(Ss)^2+4Ss\mathrm{sin}^2(ak)}.$$
(9)
Using the Wick theorem, $`_1`$ is rewritten as
$`_1`$ $`=`$ $`E_1J{\displaystyle \underset{k}{}}\left(\delta \omega _k^{}\alpha _k^{}\alpha _k+\delta \omega _k^+\beta _k^{}\beta _k\right)`$ (10)
$`+`$ $`_{\mathrm{irrel}}+_{\mathrm{quart}},`$ (11)
where the $`O(S^0)`$ correction to the ground-state energy, $`E_1`$, and those to the dispersions, $`\delta \omega _k^\pm `$, are, respectively, given by
$`E_1=2JN\left[\mathrm{\Gamma }_1^2+\mathrm{\Gamma }_2^2+\left(\sqrt{S/s}+\sqrt{s/S}\right)\mathrm{\Gamma }_1\mathrm{\Gamma }_2\right],`$ (12)
$`\delta \omega _k^\pm =2(S+s)\mathrm{\Gamma }_1{\displaystyle \frac{\mathrm{sin}^2(ak)}{\omega _k}}+{\displaystyle \frac{\mathrm{\Gamma }_2}{\sqrt{Ss}}}\left[\omega _k\pm (Ss)\right],`$ (13)
with
$$\begin{array}{c}\mathrm{\Gamma }_1=\frac{1}{2N}\underset{k}{}\left(\frac{S+s}{\omega _k}1\right),\hfill \\ \mathrm{\Gamma }_2=\frac{1}{N}\underset{k}{}\frac{\sqrt{Ss}\mathrm{cos}^2(ak)}{\omega _k},\hfill \end{array}$$
(14)
while the irrelevant one-body terms
$$_{\mathrm{irrel}}=J\frac{(Ss)^2}{\sqrt{Ss}}\mathrm{\Gamma }_1\underset{k}{}\frac{\mathrm{cos}(ak)}{\omega _k}(\alpha _k\beta _k+\alpha _k^{}\beta _k^{}),$$
(15)
and the residual two-body interactions $`_{\mathrm{quart}}`$ are both neglected so as to keep the ferromagnetic branch gapless. The resultant Hamiltonian is compactly represented as
$$E_\mathrm{g}+J\underset{k}{}\left(\stackrel{~}{\omega }_k^{}\alpha _k^{}\alpha _k+\stackrel{~}{\omega }_k^+\beta _k^{}\beta _k\right),$$
(16)
with
$`\stackrel{~}{\omega }_k^\pm =\omega _k^\pm \delta \omega _k^\pm ,`$ (17)
$`E_\mathrm{g}=E_{\mathrm{class}}+E_0+E_1.`$ (18)
Now we compare all the calculations in Fig. 1. The diagonalization results fully demonstrate the steady applicability and good precision of our Monte Carlo treatment. The spin-wave approach generally gives a good description of the low-energy structure. Even the free spin waves allow us to have a qualitative view of the elementary excitations. The relatively poor description of the antiferromagnetic branch by the free spin waves implies that the quantum effect is more relevant in the antiferromagnetic branch. Here we may be reminded of the spin-wave treatment of mono-spin Heisenberg magnets. For the ferromagnetic chains, the spin-wave dispersions are nothing but the exact picture of the elementary excitations, while for the antiferromagnetic chains, those are generally no more than a qualitative view. The point is that in the present system the spin-wave picture is efficient for both elementary excitation branches and the spin-wave series potentially lead to an accurate description. More specifically, the divergence of the boson numbers, which plagues the antiferromagnetic spin-wave treatment in one dimension, does not occur in the present system. This viewpoint is further discussed in the final section. We conclude this section by listing in Table I the spin-wave estimates of a few interesting quantities in comparison with the numerical solutions, where we define the curvature $`v`$ as $`\stackrel{~}{\omega }_{k0}^{}=v(2ak)^2`$.
## III Thermodynamic Properties
The dual structure of the excitations leads to unique thermal behavior. In order to complement quantum Monte Carlo thermal calculations, especially at low temperatures, we consider describing the thermodynamics in terms of the spin-wave theory. Introducing the additional constraint of the total magnetization being zero into the conventional spin-wave theory, Takahashi succeeded in obtaining an excellent description of the low-temperature thermal behavior of one-dimensional Heisenberg ferromagnets. The present authors have recently applied the idea to the spin-$`(1,\frac{1}{2})`$ ferrimagnetic Heisenberg chain . Here we develop the method for general spin cases and make a detailed analysis of its validity as a function of $`(S,s)`$. The core idea of the so-called modified spin-wave theory can be summarized as controlling the boson numbers by imposing a certain constraint on the magnetization. From this point of view, the zero-magnetization constraint, which is quite reasonable for isotropic magnets, plays a relevant role in ferromagnets. The resultant low-temperature expansions of the specific heat and susceptibility of the spin-$`s`$ Heisenberg ferromagnetic chain,
$`{\displaystyle \frac{C}{Nk_\mathrm{B}}}={\displaystyle \frac{3}{8s^{\frac{1}{2}}}}{\displaystyle \frac{\zeta (\frac{3}{2})}{\sqrt{\pi }}}t^{\frac{1}{2}}{\displaystyle \frac{1}{2s^2}}t+O(t^{\frac{3}{2}}),`$ (19)
$`{\displaystyle \frac{\chi J}{N(g\mu _\mathrm{B})^2}}={\displaystyle \frac{2s^4}{3}}t^2s^{\frac{5}{2}}{\displaystyle \frac{\zeta (\frac{1}{2})}{\sqrt{\pi }}}t^{\frac{3}{2}}`$ (20)
$`+{\displaystyle \frac{s}{2}}\left[{\displaystyle \frac{\zeta (\frac{1}{2})}{\sqrt{\pi }}}\right]^2t^1+O(t^{\frac{1}{2}}),`$ (21)
with $`t=k_\mathrm{B}T/J`$ and Riemann’s zeta function $`\zeta (z)`$, indeed coincide with the thermodynamic Bethe-ansatz calculations for $`s=\frac{1}{2}`$ within the leading few terms.
In the present system, the zero-magnetization constraint is explicitly represented as
$$\underset{j}{}S_j^z+s_j^z=(Ss)N\underset{k}{}\underset{\sigma =\pm }{}\sigma \stackrel{~}{n}_k^\sigma =0,$$
(22)
where $`n_k^\pm =_{n^{},n^+}n^\pm P_k(n^{},n^+)`$ with $`P_k(n^{},n^+)`$ being the probability of $`n^{}`$ ferromagnetic and $`n^+`$ antiferromagnetic spin waves appearing in the $`k`$-momentum state. Equation (22) straightforwardly proposes that the thermal spin deviation should cancel the Néel-state magnetization. By minimizing the free energy
$`F=E_\mathrm{g}+{\displaystyle \underset{k}{}}(\stackrel{~}{n}_k^{}\stackrel{~}{\omega }_k^{}+\stackrel{~}{n}_k^+\stackrel{~}{\omega }_k^+)`$ (23)
$`+k_\mathrm{B}T{\displaystyle \underset{k}{}}{\displaystyle \underset{n^{},n^+}{}}P_k(n^{},n^+)\mathrm{ln}P_k(n^{},n^+),`$ (24)
with respect to $`P_k(n^{},n^+)`$ at each $`k`$ under the condition (22) as well as the trivial constraints $`_{n^{},n^+}P_k(n^{},n^+)=1`$, we obtain
$`{\displaystyle \frac{C}{Nk_\mathrm{B}}}={\displaystyle \frac{3}{4}}\left({\displaystyle \frac{Ss}{Ss}}\right)^{\frac{1}{2}}{\displaystyle \frac{\zeta (\frac{3}{2})}{\sqrt{2\pi }}}\stackrel{~}{t}^{\frac{1}{2}}{\displaystyle \frac{1}{Ss}}\stackrel{~}{t}+O(\stackrel{~}{t}^{\frac{3}{2}}),`$ (25)
$`{\displaystyle \frac{\chi J}{N(g\mu _\mathrm{B})^2}}={\displaystyle \frac{Ss(Ss)^2}{3}}\stackrel{~}{t}^2(Ss)^{\frac{1}{2}}(Ss)^{\frac{3}{2}}{\displaystyle \frac{\zeta (\frac{1}{2})}{\sqrt{2\pi }}}\stackrel{~}{t}^{\frac{3}{2}}`$ (26)
$`+(Ss)\left[{\displaystyle \frac{\zeta (\frac{1}{2})}{\sqrt{2\pi }}}\right]^2\stackrel{~}{t}^1+O(\stackrel{~}{t}^{\frac{1}{2}}),`$ (27)
where $`\stackrel{~}{t}=k_\mathrm{B}T/J\gamma `$ with $`\gamma =1\mathrm{\Gamma }_1(S+s)/Ss\mathrm{\Gamma }_2/\sqrt{Ss}`$. The specific heat $`C`$ has been obtained by differentiating the free energy $`F`$, whereas the susceptibility $`\chi `$ by calculating $`(M^2M^2)/3T`$, where we have set the $`g`$ factors of the spins $`𝑺`$ and $`𝒔`$ both equal to $`g`$, because the difference between them amounts to at most several percent of themselves in practice .
The analytic expressions (25) and (27) give us a bird’s-eye view of the one-dimensional ferrimagnetic behavior. The spin-$`(S,s)`$ ferrimagnet turns into the spin-$`s`$ antiferromagnet in the limit of $`Ss`$, whereas it looks like the spin-$`S`$ ferromagnet in the limit of $`S/s\mathrm{}`$. In this sense, the subtraction $`Ss`$ may be regarded as the ferromagnetic contribution, while the residual spin amplitude $`2s`$ as the antiferromagnetic one. No ferromagnetic aspect at $`S/s=1`$, while a hundred percent ferromagnetic aspect for $`S/s\mathrm{}`$. Another consideration also leads us to such a picture. Since the perturvation from the decoupled dimers qualitatively well describes the low-lying excitations of the system, we propose in Fig. 2 an idea of decomposing ferrimagnets into ferromagnets and antiferromagnets, where we let the decoupled dimers and the Affleck-Kennedy-Lieb-Tasaki valence-bond-solid states symbolize ferrimagnets and integer-spin gapped antiferromagnets, respectively. Now we expect spin-$`(S,s)`$ ferrimagnets to behave like combinations of spin-$`(2s)`$ antiferromagnets and spin-$`(Ss)`$ ferromagnets. Since the antiferromagnetic excitations of the ferrimagnetic ground state are gapped, the low-temperature behavior of ferrimagnets should be only of ferromagnetic aspect. At low temperatures there is indeed no effective contribution from the spin-$`(2s)`$ antiferromagnetic chain with an excitation gap immediately above the ground state. In this context, we are surprised but pleased to find that provided $`S=2s`$, the expressions (25) and (27) coincide with those for ferromagnets, (19) and (21), except for the quantum renormalizing factor $`\gamma `$. The low-temperature thermodynamics should be dominated by the small-momentum ferromagnetic excitations. Within the linearized spin-wave theory, the small-momentum dispersions of Heisenberg ferrimagnets and ferromagnets are characterized by the curvatures
$`v^{(S,s)\text{-}\mathrm{ferri}}={\displaystyle \frac{Ss}{2(Ss)}}J,`$ (28)
$`v^{(Ss)\text{-}\mathrm{ferro}}=(Ss)J,`$ (29)
respectively, where we find that they coincide with each other only when $`S=2s`$. The criterion, $`S=2s`$, is further convincing when we consider the high-temperature behavior. The paramagnetic behavior of the spin-$`(S,s)`$ ferrimagnet is given as
$`{\displaystyle \frac{F^{(S,s)\text{-}\mathrm{ferri}}}{Nk_\mathrm{B}T}}=\mathrm{ln}\left[(2S+1)(2s+1)\right],`$ (30)
$`{\displaystyle \frac{k_\mathrm{B}T\chi ^{(S,s)\text{-}\mathrm{ferri}}}{Ng^2\mu _\mathrm{B}^2}}={\displaystyle \frac{1}{3}}\left[S(S+1)+s(s+1)\right],`$ (31)
whereas those of the spin-$`(Ss)`$ ferromagnet and the spin-$`(2s)`$ antiferromagnet are as follows:
$`{\displaystyle \frac{F^{(Ss)\text{-}\mathrm{ferro}}+F^{(2s)\text{-}\mathrm{antiferro}}}{Nk_\mathrm{B}T}}`$ (32)
$`=\mathrm{ln}\left[(2S2s+1)(4s+1)\right],`$ (33)
$`{\displaystyle \frac{k_\mathrm{B}T(\chi ^{(Ss)\text{-}\mathrm{ferro}}+\chi ^{(2s)\text{-}\mathrm{antiferro}})}{Ng^2\mu _\mathrm{B}^2}}`$ (34)
$`={\displaystyle \frac{1}{3}}\left[(Ss)(Ss+1)+2s(2s+1)\right].`$ (35)
These asymptotic values agree with each other only when $`S=2s`$. This is simply the consequence of the degrees of freedom. In ferrimagnets of $`S>2s`$, the ferromagnetic spin degrees of freedom overbalance the antiferromagnetic ones, while in ferrimagnets of $`S<2s`$, vice versa. Only the balanced ferrimagnet with $`S=2s`$ is well approximated as the simple combination of the spin-$`(Ss)`$ ferromagnet and the spin-$`(2s)`$ antiferromagnet. However, we note that even in the case of $`S=2s`$, the similarity between the ferrimagnetic behavior, (25) and (27), and the ferromagnetic one, (19) and (21), does not go beyond the leading few terms shown here. The ferromagnetic features of ferrimagnets are thermally smeared out. On the other hand, the ferrimagnetic behavior further deviates from the purely ferromagnetic one due to the quantum effect characterized by $`\gamma `$. The low-temperature expansions (25) and (27) imply that as temperature goes to zero, the quantum effect is reduced for the specific heat $`C`$, whereas enhanced for the susceptibility $`\chi `$. In the limit of $`S/s\mathrm{}`$, the quantum corrections $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ both vanish.
Although the spin-wave theory combined with the additional constraint (22) is so useful, it should further be modified at higher temperatures so as to control the total number of the bosons $`_{\sigma =\pm }\stackrel{~}{n}_k^\sigma `$ rather than the subtraction $`_{\sigma =\pm }\sigma \stackrel{~}{n}_k^\sigma `$. In our naivest attempt to improve the theory, we replace the constraint (22) by
$`{\displaystyle \underset{j}{}}S_j^zs_j^z=(S+s2\mathrm{\Gamma }_1)N`$ (36)
$`(S+s){\displaystyle \underset{k}{}}{\displaystyle \underset{\sigma =\pm }{}}{\displaystyle \frac{\stackrel{~}{n}_k^\sigma }{\omega _k}}=0.`$ (37)
However, the alternative condition (37) changes the low-temperature description, (25) and (27), which should be kept unchanged under any artificial constraint, as well as considerably underestimates the Schottky-like characteristic peak of the specific heat. In an attempt to remove $`\mathrm{\Gamma }_1`$ from Eq. (37), we reach a phenomenological modification
$$\underset{j}{}:S_j^zs_j^z:=(S+s)N(S+s)\underset{k}{}\underset{\sigma =\pm }{}\frac{\stackrel{~}{n}_k^\sigma }{\omega _k}=0,$$
(38)
which proposes that the thermal fluctuation should cancel the Néel-state magnetization and results in the same low-temperature description as Eqs. (25) and (27), where the normal ordering is taken with respect to both operators $`\alpha `$ and $`\beta `$. Now the overall description of the thermal quantities is numerically obtained and is shown, together with the precise quantum Monte Carlo calculations, in Figs. 3 and 4. We stress that these are the spin-wave description of one-dimensional magnets. We are allowed to recognize that the spin-wave picture is qualitatively valid even at high temperatures. The spin-wave theory, which is to overestimate the spin degrees of freedom, inevitably overestimates the energy-derivative quantities and under estimate the magnetization-derivative ones at high temperatures. However, the location of the Schottky-like peak of the specific heat is generally reproduced well, which should be attributed to the fine description of the antiferromagnetic excitation gap by the interacting spin waves. One might say that the spin-wave description is a overlikely portrait of the actual behavior. For the ferromagnetic ferrimagnet of $`(S,s)=(\frac{3}{2},\frac{1}{2})`$, for instance, the low-temperature shoulder of $`C`$ is too pronounced and the high-temperature antiferromagnetic increase of $`\chi T`$ is too suppressed in the spin-wave theory. We note finally in this section that the double constraint does not improve the theory at all. At high temperatures, the constraint (38) dominates the thermal behavior but the constraint (22) works little. At low temperatures, where $`\stackrel{~}{n}_k^{}\stackrel{~}{n}_k^+`$, the two constraints (22) and (38) are almost degenerate and no numerical solution is found for the couple of Lagrange multipliers due to the double constraint. Further inclusion of any tricky constraint is likely to make us lose sight of the physical basis.
## IV Discussion
We have studied the low-energy structure and the thermal properties of Heisenberg ferrimagnetic spin chains featuring the spin-wave theory. Even with the modified spin-wave theory, it is still hard to describe the thermodynamics over the whole temperature range. However, considering the poor applicability of the spin-wave theory to one-dimensional antiferromagnets, the obtained description is quite successful. What is the difference between ferrimagnets and antiferromagnets, in spite of the antiferromagnetic coupling between nearest neighbors in common? The surprising efficiency of the present spin-wave treatment can be attributed to the ordered ground state of ferrimagnets and thus to the nondiverging sublattice magnetization. The key constant $`\mathrm{\Gamma }_1`$ is nothing but the quantum spin reduction
$$\frac{1}{N}\underset{j}{}a_j^{}a_j_\mathrm{g}=\frac{1}{N}\underset{j}{}b_j^{}b_j_\mathrm{g},$$
(39)
where $`A_\mathrm{g}`$ denotes the ground-state average of $`A`$. $`\mathrm{\Gamma }_1`$ monotonically decreases as $`S/s`$ increases, diverging at $`S/s=1`$ and vanishing for $`S/s\mathrm{}`$. Since the spin reduction $`\mathrm{\Gamma }_1`$ can be a measure for the validity of the spin-wave description, the spin-wave theory in general works better as $`S/s`$, as well as $`Ss`$, increases.
In this context, it is interesting to observe the ground-state spin correlations,
$`f_S(r)=\{\begin{array}{cc}S_j^zS_{j+l}^z\hfill & \text{for }r=2la,\hfill \\ S_j^zs_{j+l}^z\hfill & \text{for }r=(2l+1)a,\hfill \end{array}`$ (42)
$`f_s(r)=\{\begin{array}{cc}s_j^zs_{j+l}^z\hfill & \text{for }r=2la,\hfill \\ s_j^zS_{j+l+1}^z\hfill & \text{for }r=(2l+1)a.\hfill \end{array}`$ (45)
Figure 5 shows the small-$`r`$ initial behavior of $`f_S(r)`$ and $`f_s(r)`$, which demonstrates the considerably small correlation length and the existence of the long-range order in the present system. For the asymptotic correlation between the two larger spins far distant from each other, for example, the spin-wave estimate deviates from the actual value by $`23`$%, $`7`$%, $`18`$%, and $`7`$% for $`(S,s)=(1,\frac{1}{2}),(\frac{3}{2},\frac{1}{2}),(\frac{3}{2},1)`$, and $`(2,1)`$, respectively. We find that the above-defined criterion in terms of $`S/s`$ and $`Ss`$ works fairly well. The theory is not so bad even in the extreme quantum case.
Quite recently, the nuclear spin relaxation time in a one-dimensional ferrimagnetic Heisenberg model compound, NiCu(C<sub>7</sub>H<sub>6</sub>N<sub>2</sub>O<sub>6</sub>)(H<sub>2</sub>O)<sub>3</sub>$``$2H<sub>2</sub>O, has been measured and its temperature and field dependences have successfully been interpreted within the framework of the spin-wave theory . The theory must still be open to further applications to this fascinating system. We hope that the present study will motivate further explorations into low-dimensional ferrimagnets in both theoretical and experimental fields. The ferromagnetic and antiferromagnetic mixed nature may further be discussed from different points of view. For instance, the topological excitations such as instantons in ferromagnets and the topological terms in antiferromagnets might give further support to the present understanding of ferrimagnets.
## ACKNOWLEDGMENTS
This work was supported by the Japanese Ministry of Education, Science, and Culture through Grant-in-Aid No. 11740206 and by the Sanyo-Broadcasting Foundation for Science and Culture. T. F. was supported by the JSPS Postdoctoral Fellowship for Research Abroad. The numerical computation was done in part using the facility of the Supercomputer Center, Institute for Solid State Physics, University of Tokyo. |
warning/0001/hep-ph0001317.html | ar5iv | text | # DETERMINING THE CP VIOLATING PHASE 𝛾 11footnote 1To appear in Proceedings of The Third International Conference on B Physics and CP Violation, Taipei, December 3-7, 1999, H. -Y. Cheng and W. -S. Hou, eds. (World Scientific, 2000).
## 1 Introduction
$`B`$ meson decays open the window into new phenomena of CP nonconservation, providing useful information about CP violating phases . Phase measurements for $`\alpha ,\beta `$ and $`\gamma `$, which are not independent in the Kobayashi-Maskawa framework, are important for two reasons. First, they improve to a higher precision our knowledge of the CKM mixing matrix, consisting of fundamental quark couplings in our present theory. And second, they are sensitive probes for sources of CP violation outside the CKM matrix, and for new flavor changing interactions contributing to $`B^0\overline{B}^0`$ mixing and to rare $`B`$ decays.
The phase $`\beta \mathrm{Arg}(V_{cb}^{}V_{cd}/V_{tb}^{}V_{td})(=\mathrm{Arg}V_{td}^{}`$ in the standard phase convention) is measured by the time-dependent CP asymmetry in $`B^0(t)\psi K_S`$ . Theoretically, this measurement provides a very clean determination of $`\mathrm{sin}(2\beta )`$, since the single phase approximation holds in this case to better than 1$`\%`$ . The recent CDF measurement , $`\mathrm{sin}(2\beta )=0.79_{0.44}^{+0.41}`$, is an encouraging proof of the method, although it does not yet constitute definite evidence for CP violation in $`B`$ decays, as predicted in the CKM framework.
The determination of $`\mathrm{sin}2\alpha `$ ($`\alpha \mathrm{Arg}(V_{tb}^{}V_{td}/V_{ub}^{}V_{ud})`$) in $`B^0\pi ^+\pi ^{}`$ suffers from the appearance of a second amplitude ($`P`$) due to QCD penguin operators . The weak phase of $`P`$ differs from the phase of the dominant tree amplitude ($`T`$). The ratio $`|P|/|T|=0.3\pm 0.1`$, obtained by comparing within broken flavor SU(3) the measured rates of $`B\pi \pi `$ and $`BK\pi `$, implies a potentially large deviation of the measured asymmetry from $`\mathrm{sin}2\alpha `$. In fact, the time-dependent asymmetry in $`B^0(t)\pi ^+\pi ^{}`$ contains two terms
$$A(t)=a_{\mathrm{dir}}\mathrm{cos}(\mathrm{\Delta }mt)+\sqrt{1a_{\mathrm{dir}}^2}\mathrm{sin}2(\alpha +\theta )\mathrm{sin}(\mathrm{\Delta }mt),$$
(1)
where $`a_{\mathrm{dir}}`$ and $`\theta `$ are both proportinal to $`|P|/|T|`$ and are functions of the relative strong phase between $`P`$ and $`T`$. A precise knowledge of $`|P|/|T|`$ would permit an accurate determination of $`\alpha `$ from a measurement of the two terms in the asymmetry . One way of calculating $`|P|/|T|`$ in perurbative QCD is described by Deshpande at this meeting . However, the uncertainty in this calculation has not yet been shown to be under control.
As long as $`|P|/|T|`$ cannot be calculated reliably, a theoretically cleaner way of resolving the penguin correction will require combining the asymmetry in $`B^0\pi ^+\pi ^{}`$ with other measurements. A very early suggestion was to also measure the rates of isospin related processes, $`B^+\pi ^+\pi ^0`$ and $`B^0/\overline{B}^0\pi ^0\pi ^0`$. However, it is expected that the color-suppressed rates of tagged $`B^0`$ and $`\overline{B}^0`$ decays to neutral pions will involve large experimental errors, at least in the first round of experiments. Alternatively, $`\alpha `$ can be determined by applying the isospin technique to $`B\rho \pi `$ decays , or by relating the time-dependence in $`B^0(t)\pi ^+\pi ^{}`$ and in $`B_s(t)K^+K^{}`$ using flavor U-spin symmetry . Since none of these methods is expected to be both theoretically clean and experimentally accessible in the near future, some of these methods will have to be combined in order to reduce the error in $`\alpha `$.
The phase $`\gamma \mathrm{Arg}(V_{ub}^{}V_{ud}/V_{cb}^{}V_{cd})(=\mathrm{Arg}V_{ub}^{}`$ in the standard convention) is the relative weak phase between a CKM-favored ($`bc`$) and a CKM-suppressed ($`bu`$) decay amplitude. Therefore, it contributes to direct CP asymmetries, in which two such amplitudes interfere, without the need for neutral $`B`$ mixing. Asymmetry measurements in charged $`B`$ decays and in self-tagged neutral $`B`$ decays have the advantage of avoiding the price of flavor tagging. However their theoretical interpretation in terms of the weak phase $`\gamma `$ involves an unknown strong phase and an unknown ratio of two interferening hadronic matrix elements. A few methods were proposed to overcome this difficulty, often by measuring not only the asymmetry, but also the rates of certain related processes which provide information on these matrix elements. In the next few sections we describe several ways of measuring $`\gamma `$ in four classes of processes, discussing in each case both the theoretical uncertainties and the experimental limitations.
## 2 $`BDK^{()}`$
A simple idea for measuring $`\gamma `$ in $`B^\pm DK^\pm `$ was proposed some time ago . Neglecting very small CP violation in $`D^0\overline{D}^0`$ mixing, one can write a triangle amplitude relation
$$\sqrt{2}A(B^+D_1^0K^+)=A(B^+D^0K^+)+A(B^+\overline{D}^0K^+).$$
(2)
$`D_1^0`$ is an even CP-eigenstate, decaying, for instance, to $`\pi ^+\pi ^{}`$, while $`D^0`$ and $`\overline{D}^0`$ are two opposite flavor states. The two amplitudes on the right-hand-side contain CKM factors $`V_{ub}^{}V_{cs}`$ and $`V_{cb}^{}V_{us}`$, both of which are $`𝒪(\lambda ^3)`$ , and have a relative weak phase $`\gamma `$. If the two amplitudes were of about equal magnitudes, then the triangle relation and its charge-conjugate would permit a determination of $`\mathrm{sin}\gamma `$ within certain discrete ambiguities.
At a closer examination one observes, however, that the first amplitude is suppressed relative to the second one by a CKM factor, $`|V_{ub}^{}V_{cs}|/|V_{cb}^{}V_{us}|0.4`$, and by a color-suppression factor of about 0.25. \[This prediction comes from evidence for color-supression in $`B\overline{D}\pi `$ \]. An order of magnitude suppression of this amplitude relative to the measured amplitude of $`B^+\overline{D}^0K^+`$, causes serious experimental difficulties in tagging the flavor of $`D^0`$, through decays such as $`D^0K^{}\pi ^+`$. The very rare decay $`B^+D^0K^+(K^{}\pi ^+)K^+`$ interferes strongly with the doubly Cabibbo-suppressed decay of $`\overline{D}^0`$ in $`B^+\overline{D}^0K^+(K^{}\pi ^+)K^+`$.
Two ways were proposed for partially evading this difficulty by considering only very rare decays, typically with branching ratios of order $`10^7`$. Aside from the small rates, both methods have other limitations.
* Study only color-suppressed decays $`B^0D^0(\overline{D}^0)K^0`$, where the flavors of neutral $`D`$ and $`K^{}`$ are determined through $`D^0K^{}\pi ^+,K^0K^+\pi ^{}`$. Here the interference between Cabibbo-favored and Cabibbo-suppressed neutral $`D`$ decays is much weaker. Still, such interference occurs and prohibits a precise determination of $`\gamma `$.
* Study $`B^+D^0(\overline{D}^0)K^+fK^+`$, with two neutral $`D`$ decay modes, such as $`f=K^{}\pi ^+,K^{}\rho ^+`$. In this case the two interfering amplitudes have comparable magnitudes, one being color-suppressed in $`B`$ decay and the other being doubly-Cabibbo-suppressed in $`D`$ decay. Measurement of the rates of these two processes and their charge-conjugates would permit a determination of $`\gamma `$ , provided that the two doubly Cabibbo-suppressed $`D`$ decay branching ratios were known. Uncertainties in these branching ratios prohibit an accurate determination of $`\gamma `$. An intrinsic uncertainty follows from the difficulty of disentangling doubly Cabibbo-suppressed $`D^0`$ decays from $`D^0\overline{D}^0`$ mixing . The mixing may be larger than conventionally estimated , as anticipated recently from large resonance contributions in $`D^0`$ decays . Therefore, a precise measurement of $`\gamma `$ using this method requires knowledge of the $`D^0\overline{D}^0`$ mixing parameters.
There are also variants , which combine the above two schemes based on CP and flavor states. In all cases, the very small branching ratios of the color-suppressed processes imply that such measurements cannot be carried out effectively in the first round of experiments at $`e^+e^{}`$ $`B`$ factories, and would have to wait for hadron collider experiments providing at least $`10^9`$ $`B`$’s.
An interesting question remains : What can be learned by studying only the more abundant color-allowed $`B^\pm DK^\pm `$ decays which have larger branching ratios? Considers the charge-averaged ratio of rates for neutral $`D`$ mesons decaying to an even (odd) CP state and for a color-allowed flavor state
$$R_i\frac{2[\mathrm{\Gamma }(B^+D_iK^+)+\mathrm{\Gamma }(B^{}D_iK^{})]}{\mathrm{\Gamma }(B^+\overline{D}^0K^+)+\mathrm{\Gamma }(B^{}D^0K^{})},i=1,2.$$
(3)
One finds (neglecting the small $`D_1^0D_2^0`$ width-difference)
$$R_{1,2}=1+\overline{r}^2\pm 2\overline{r}\mathrm{cos}\overline{\delta }\mathrm{cos}\gamma ,$$
(4)
where $`A(B^+D^0K^+)/A(B^+\overline{D}^0K^+)\overline{r}\mathrm{exp}[i(\overline{\delta }+\gamma )]`$. This leads to two general inequalities
$$\mathrm{sin}^2\gamma R_{1,2},i=1,2,$$
(5)
one of which must imply a new constraint on $`\gamma `$, unless $`\gamma =\pi /2`$.
Assuming, for instance, $`\overline{r}=0.1,\overline{\delta }=0,\gamma =40^{}`$, one finds $`R_2=0.85`$. With $`10^8B^+B^{}`$ pairs, and using measured $`B`$ and $`D`$ decay branching ratios , one estimates an error $`R_2=0.85\pm 0.05`$. In this case, Eq.(5) excludes the range $`73^{}<\gamma <107^{}`$ with 90$`\%`$ confidence level. Including measurements of the CP asymmetries in $`B^\pm D_iK^\pm `$ permit, in principle, a determination of $`\gamma `$ (and not only bounds on the angle). This depends, of course, on the unknown strong phase $`\overline{\delta }`$. For further studies of related methods, see .
## 3 $`B_s^0(t)D_sK`$
The interference between the two $`\lambda ^3`$ subprocesses $`bc\overline{u}s`$ and $`bu\overline{c}s`$ operates also in the time-dependent decay $`B_s^0D_s^{}K^+`$, leading to a $`\mathrm{sin}(\mathrm{\Delta }m_st)`$ term
$`\mathrm{\Gamma }(B_s^0(t)D_s^{}K^+)`$ $`=`$ $`e^{\mathrm{\Gamma }_st}[|A|^2\mathrm{cos}^2({\displaystyle \frac{\mathrm{\Delta }m_st}{2}})+|\overline{A}|^2\mathrm{sin}^2({\displaystyle \frac{\mathrm{\Delta }m_st}{2}})`$ (6)
$``$ $`|A\overline{A}|\mathrm{sin}(\delta +\gamma )\mathrm{sin}(\mathrm{\Delta }m_st)].`$
The two color-allowed amplitudes, corresponding to $`bc\overline{u}s`$ and $`bu\overline{c}s`$, have magnitudes $`|A|`$ and $`|\overline{A}|`$, and involve relative strong and weak phases, $`\delta `$ and $`\gamma `$, respectively.
These four parameters describe in a similar way the time-dependence in the three processes in which the initial and/or final states are charge-conjugated, $`\overline{B}_s^0(t)D_s^{}K^+,B_s^0(t)D_s^+K^{},\overline{B}_s^0(t)D_s^+K^{}`$. Thus, measuring the time-dependent oscillations in these four processes, all of which require tagging the flavor of initial $`B_s^0`$, permits a determination of $`\gamma `$ . It is obvious that, for an accurate measurement, one would also need to know the width-difference between the two $`B_s`$ mass eigenstates, which was neglected in (6). Further studies and discussions of width-dependent effects can be found in .
## 4 $`BPP`$
We will consider $`B`$ decays to two light pseudoscalars, $`BPP`$ where $`P=\pi ,K`$, within the framework of flavor SU(3) symmetry . Final states involving $`\eta `$ and $`\eta ^{}`$ can be studied similarly . Occasional reference to SU(3) breaking effects will be made. A more ambitious approach, relying on generalized factorization , has not yet been justified quantitatively to a satisfactory level.
The low energy effective weak Hamiltonian governing $`BPP`$
$$=\frac{G_F}{\sqrt{2}}\underset{q=d,s}{}\left(\underset{q^{}=u,c}{}\lambda _q^{}^{(q)}[c_1Q_1^{(q^{})}+c_2Q_2^{(q^{})}]\lambda _t^{(q)}\underset{i=3}{\overset{10}{}}c_iQ_i^{(q)}\right),$$
(7)
where $`\lambda _q^{}^{(q)}=V_{q^{}b}^{}V_{q^{}q}`$, consists of the sum of three types of four quark operators: two $`(VA)(VA)`$ current-current operators ($`Q_{1,2}`$), four QCD penguin operators ($`Q_{3,4,5,6}`$), and four electroweak penguin (EWP) operators ($`Q_{7,8,9,10}`$) with different chiral structures. The EWP operators with dominant Wilson coefficients, $`Q_9`$ and $`Q_{10}(|c_{7,8}|/|c_9|0.04`$), have both a $`(VA)(VA)`$ structure. All the four-quark operators, $`(\overline{b}q_1)(\overline{q}_2q_3)`$, can be decomposed into a sum of $`\overline{\mathrm{𝟏𝟓}}`$, $`\mathrm{𝟔}`$ and $`\overline{\mathrm{𝟑}}`$ representations . The QCD penguin operators are pure $`\overline{\mathrm{𝟑}}`$.
A simple proportionality relation was noted recently to hold between the dominant EWP operators $`Q_9`$ and $`Q_{10}`$ and the current-current (“tree”) operators $`Q_1`$ and $`Q_2`$, both trasforming as given SU(3) representations
$`_{EWP}^{(q)}(\overline{\mathrm{𝟏𝟓}})`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{c_9+c_{10}}{c_1+c_2}}{\displaystyle \frac{\lambda _t^{(q)}}{\lambda _u^{(q)}}}_T^{(q)}(\overline{\mathrm{𝟏𝟓}}),`$ (8)
$`_{EWP}^{(q)}(\mathrm{𝟔})`$ $`=`$ $`{\displaystyle \frac{3}{2}}{\displaystyle \frac{c_9c_{10}}{c_1c_2}}{\displaystyle \frac{\lambda _t^{(q)}}{\lambda _u^{(q)}}}_T^{(q)}(\mathrm{𝟔}).`$ (9)
The superscripts $`q=d,s`$ denote strangeness-conserving and strangeness-changing transitions, respectively. The two ratios of Wilson coefficients are equal within 3$`\%`$
$$\frac{c_9+c_{10}}{c_1+c_2}\frac{c_9c_{10}}{c_1c_2}1.12\alpha .$$
(10)
As a consequence of (8) and (9), processes given by $`\overline{\mathrm{𝟏𝟓}}`$ and $`\mathrm{𝟔}`$ transitions contain EWP and current-current contributions which are proportional to each other. For instance, in the $`\mathrm{\Delta }S=1`$ pure $`\overline{\mathrm{𝟏𝟓}}`$ transition, $`B(K\pi )_{I=3/2}`$, where $`|I=3/2=|K^0\pi ^++\sqrt{2}|K^+\pi ^0`$, the ratio of these amplitudes
$$\delta _{EW}=\frac{3}{2}\frac{c_9+c_{10}}{c_1+c_2}\frac{|V_{tb}^{}V_{ts}|}{|V_{ub}^{}V_{us}|}=0.65\pm 0.15$$
(11)
is of order one , in spite of the fact that EWP amplitudes are higher order in electroweak couplings.
This simple result, obtained in the limit of flavor SU(3), was used in order to obtain a bound on $`\gamma `$ in $`B^\pm K\pi `$. Expressing $`B^+K\pi `$ in terms of reduced SU(3) amplitudes depicted in graphical form , one has
$`A(B^+K^0\pi ^+)`$ $`=`$ $`\lambda _t^{(s)}[P+EW]+|\lambda _u^{(s)}|e^{i\gamma }[P_{uc}+A],`$ (12)
$`\sqrt{2}A(B^+K^+\pi ^0)`$ $`=`$ $`\lambda _t^{(s)}[P+EW]|\lambda _u^{(s)}|[(T+C)(e^{i\gamma }\delta _{EW})+P_{uc}+A].`$
The dominant QCD penguin amplitudes $`P`$ in the two processes, carrying a weak phase $`\mathrm{Arg}\lambda _t^{(s)}=\pi `$, are equal because of isospin. Unequal rates of the two processes would be evidence for interference with smaller amplitudes involving a different weak phase. Defining the charge-averaged ratio of rates
$$R_{}^1\frac{2[B(B^+K^+\pi ^0)+B(B^{}K^{}\pi ^0)]}{B(B^+K^0\pi ^+)+B(B^{}\overline{K}^0\pi ^{})},$$
(13)
one finds, to leading order in small quantities
$$R_{}^1=12ϵ\mathrm{cos}\varphi (\mathrm{cos}\gamma \delta _{EW})+𝒪(ϵ^2)+𝒪(ϵϵ_A)+𝒪(ϵ_A^2).$$
(14)
$`ϵ`$ and $`ϵ_A`$ are ratios of tree-to-penguin and rescattering -to-penguin amplitudes, respectively
$$ϵe^{i\varphi }=\frac{|\lambda _u^{(s)}|}{|\lambda _t^{(s)}|}\frac{T+C}{P+EW},ϵ_A=\frac{|\lambda _u^{(s)}[P_{uc}+A]|}{|\lambda _t^{(s)}[P+EW]|}.$$
(15)
The first ratio is given by
$$ϵ=\sqrt{2}\frac{V_{us}}{V_{ud}}\frac{f_K}{f_\pi }\frac{|A(B^+\pi ^0\pi ^+)|}{|A(B^+K^0\pi ^+)|}=0.21\pm 0.05,$$
(16)
while the second one is roughly
$$ϵ_A\lambda \frac{|A(B^+\overline{K}^0K^+)|}{|A(B^+K^0\pi ^+)|}<0.12$$
(17)
Neglecting second order terms, one obtains a bound
$$|\mathrm{cos}\gamma \delta _{EW}|\frac{|1R_{}^1(K\pi )|}{2ϵ},$$
(18)
which would provide useful information about $`\gamma `$ if a value different from one were measured for $`R_{}^1`$. The present value, $`R_{\mathrm{exp}}^1=1.27\pm 0.48`$, is consistent with one. Further information about $`\gamma `$, applying also to the case $`R_{}^1=1`$, can be obtained by measuring separately $`B^+`$ and $`B^{}`$ decay rates . The solution obtained for $`\gamma `$ involves uncertainties due to SU(3) breaking in subdominant amplitudes and an uncertainty in $`|V_{ub}/V_{cb}|`$, both of which affect the value of $`\delta _{EW}`$. Combined with errors in $`ϵ|A(B^+\pi ^+\pi ^0)/A(B^+K^0\pi ^+)|`$, and in the rescattering parameter $`ϵ_A`$, the resulting uncertainty in $`\gamma `$ is unlikely to be smaller than $`20^{}`$ . For other ways of studying $`\gamma `$ in $`BK\pi `$, see .
## 5 $`BVP`$
$`B`$ mesons decays to a charmless vector meson ($`V`$) and a pseudoscalar meson ($`P`$) involve a larger number of SU(3) amplitudes than $`BPP`$. SU(3) relations between EWP and current-current amplitudes, following from (8) and (9), reduce considerably the number of independent amplitudes. In this section we describe briefly two applications .
### 5.1 $`B^\pm \rho K`$
Defining a charge-averaged ratio of rates for $`B^\pm \rho K`$
$$R_{}^1(\rho K)\frac{2[B(B^+\rho ^0K^+)+B(B^{}\rho ^0K^{})]}{B(B^+\rho ^+K^0)+B(B^{}\rho ^{}\overline{K}^0)},$$
(19)
one obtains, with some analogy to (18), the bound
$$|\mathrm{cos}\gamma |\frac{|1R_{}^1(\rho K)|}{2ϵ_V}\delta _{EW}\left(\frac{ϵ_P}{ϵ_V}\right),$$
(20)
where
$$ϵ_V=\sqrt{2}\frac{V_{us}}{V_{ud}}\frac{f_K}{f_\pi }\frac{|A(B^+\rho ^0\pi ^+)|}{|A(B^+\rho ^+K^0)|},ϵ_P=\sqrt{2}\frac{V_{us}}{V_{ud}}\frac{f_K^{}}{f_\rho }\frac{|A(B^+\rho ^+\pi ^0)|}{|A(B^+\rho ^+K^0)|}.$$
(21)
Although this constraint is weaker than (18), it shows that measuring charge-averaged ratios of rates, which differ from one, is of interest also for $`PV`$ final states.
### 5.2 $`B^0K^\pm \pi ^{}`$ vs. $`B^\pm \varphi K^\pm `$
Considering the amplitudes for $`B^0K^+\pi ^{}`$ and $`B^+\varphi K^+`$, both of which are expected to be QCD penguin dominated, and keeping only dominant and subdominant terms, one finds
$`A(B^0K^+\pi ^{})`$ $`=`$ $`\lambda _t^{(s)}P_P\lambda _u^{(s)}T_P,`$
$`A(B^+\varphi K^+)`$ $`=`$ $`\lambda _t^{(s)}[P_P+EW],`$ (22)
where the suffix $`P`$ denotes the pseudoscalar which includes the spectator in the graphic SU(3) amplitudes . The EWP contribution is related to the tree amplitude by
$$\lambda _t^{(s)}EW=\frac{1}{3}\delta _{EW}|\lambda _u^{(s)}|T_P.$$
(23)
Defining the ratio
$$R=\frac{|A(B^0K^+\pi ^{})|^2+|A(\overline{B}^0K^{}\pi ^+)|^2}{|A(B^+\varphi K^+)|^2+|A(B^{}\varphi K^{})|^2},$$
(24)
one finds
$$R=\frac{1+r^22r\mathrm{cos}\delta \mathrm{cos}\gamma }{1+(\delta _{EW}/3)^2r^2(2/3)\delta _{EW}r\mathrm{cos}\delta },$$
(25)
where
$$re^{i\delta }=\frac{|\lambda _u^{(s)}|}{|\lambda _t^{(s)}|}\frac{T_P}{P_P}.$$
(26)
Present 90$`\%`$ confidence level limits , $`(B^0K^\pm \pi ^{})>12\times 10^6`$ and $`(B^\pm \varphi K^\pm )<5.9\times 10^6`$, imply $`R>2`$, which is evidence for a nonzero contribution of $`T_P`$. In order to use this inequality for information about $`\gamma `$, one must include some input about $`r`$ and $`\delta `$. A reasonable assumption, supported both by perturbative and statistical calculations, is that $`\delta `$ does not exceed $`90^{}`$. A very conservative assumption about $`r`$ is $`r1`$. Making these two assumptions, one finds \[for $`\delta _{EW}`$ we will use the range (11)\]
$$\mathrm{cos}\gamma \frac{2}{3}\delta _{EW}<\frac{1+r^2[12(\delta _{EW}/3)^2]}{2r}.$$
(27)
This implies $`\gamma >62^{}`$ for $`r=1`$, and $`\gamma >105^{}`$ for $`r=0.5`$. Indirect evidence exists already for $`r<0.55`$. Direct information on $`r`$ will be obtained from future rate measurements of $`B^+K^0\pi ^+`$ and $`B^+\rho ^+\pi ^0`$ or $`B^0\rho ^+\pi ^{}`$.
These bounds neglect smaller terms in the amplitudes, primarily the color-suppressed tree amplitude. If this contribution is $`10\%`$ ($`20\%`$) of the color-favored tree amplitude $`T_P`$ , then the limits move up or down by about $`5^{}(10^{})`$. The bounds also assume \[by SU(3)\] equal QCD penguin contributions in the two processes. The constraint becomes stronger if the penguin amplitude in $`B^+\varphi K^+`$ is larger than in $`B^0K^+\pi ^{}`$, as predicted by factorization . However, the constraint may become weaker if SU(3) breaking in penguin amplitudes is not described by factorization.
## 6 Conclusion
We discussed several ways of determining the weak phase $`\gamma `$. The first two schemes, based on $`BDK`$ and $`B_sD_sK`$, seem at first sight to involve no hadronic uncertainties. However, at a closer look, they require taking care of smaller effects, such as $`D^0\overline{D}^0`$ mixing and the width-difference in the $`B_s^0`$ system. Methods based on charmless decays to two pseudoscalars, and decays to a pseudoscalar and a vector meson, involve dynamical hadronic parameters, such as those describing SU(3) breaking and rescattering amplitudes. Some of these quantities are under control, and others should be studied through a dialogue between theory and experiments.
## Acknowledgments
This work was supported in part by the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities, by the United States – Israel Binational Science Foundation under Research Grant Agreement 98-00237, by the Technion V.P.R. fund - Harry Werksman Research Fund, and by P. and E. Natahan Research Fund.
## References |
warning/0001/astro-ph0001409.html | ar5iv | text | # Evolution of D and 3He in the Galaxy
## 1. Introduction
In this review, I will try to describe what Galactic chemical evolution models tell us about the evolution and the primordial abundances of D and <sup>3</sup>He, and what, in turn, D and <sup>3</sup>He may tell us about stellar and Galactic evolution. In particular, I wish to emphasize that the light elements should not be treated separately, but should always be considered together with the other more diffuse elements, to better constrain their evolution.
The reason why Galactic chemical evolution models are required to derive the primordial D and <sup>3</sup>He abundances from the observed ones is that all the objects where the two elements are measurable are relatively young (the oldest being the sun with an age of 4.5 Gyr) and have therefore formed, with the only exception of high-redshift clouds, from an ISM whose chemical composition had been modified by the previous stellar generations. To infer the primordial abundances from these measurements, it is thus necessary to take into account the effects of the various cycles of gas astration and gas return, and the variations of the ISM chemical composition due to stellar nucleosynthesis and gas flows occurring up to the time when the observed objects have formed. This is accomplished by chemical evolution models.
D and <sup>3</sup>He are obviously related to each other, since all the D which enters a star is immediately burnt into <sup>3</sup>He (Reeves et al. 1973). However, the problems faced when studying their Galactic evolution are quite different, and I will thus treat them separately in this paper.
## 2. D evolution
Since D is completely destroyed inside stars already in pre-main sequence phase, if we consider the Big Bang nucleosynthesis as the only source of D, the amount of D which can be present in any galactic region and at any place is that contained in gas which has never been through stars. In other words, the fraction of primordial D surviving at any epoch and in any region is equal to the fraction of virgin gas there. Hence, in principle, to infer the primordial abundance of D from its present one it would be sufficient to know the current fraction of gas which has not entered a star yet (Steigman & Tosi 1995). Unfortunately, we do not know the fraction of pristine gas even in the most local medium and we must therefore rely on Galactic chemical evolution models to derive the D evolution.
Historically, there are two schools of thought on how to proceed in studying the D evolution, as sketched in Fig.1. The first one, in chronological order, can be referred to as the Cosmological School. The approach of this school is to start from standard Big Bang nucleosynthesis (SBBN) prescriptions, select the observational constraints on D which can be considered reliable, and infer from these sets of data what the D evolution must have been in the Galaxy. They then build models of Galaxy evolution able to reproduce the inferred trend of D vs time, and the predictions of such models on the other Galactic quantities are a by-product.
The other school, which I will call Galactic, follows the opposite approach. We start from chemical evolution models of our Galaxy, select only those which are able to reproduce the largest set of observational constraints, and take the predictions on D only from these selected models. The consequences of these predictions on the primordial D abundance and on cosmology are a by-product. Were we living in the best of all possible worlds, the two approaches should provide the same results. Instead, their predictions are quite different from each other.
Fig.2 shows the D abundances derived from all the available observations and plotted as a function of the supposed epoch of formation of the observed objects. All the error bars are $`2\sigma `$. The two vertical bars at t=0 represent the ranges of values derived by Songaila et al. (1994, hereinafter SCHR94) and Burles & Tytler (1998, B&T98) from high-redshift, low-metallicity, absorbers on the line of sight of distant QSO’s. The bar at 8.5 Gyr represents the value of the Protosolar Cloud inferred by Geiss & Gloeckler (1998, G&G98) from solar system data. The solid bar at 13 Gyr shows the range of abundances derived by Linsky (1998, L98) for the local ISM, while the length of the dotted bar shows the possible cloud-to-cloud variations suggested by Vidal-Madjar, Ferlet & Lemoine (1998, V-M98).
Since the Cosmological School was founded when the primordial <sup>4</sup>He was definitely supposed to be low (Y$`{}_{P}{}^{}`$0.23 in mass fraction), and since SBBN predicts the primordial D to be anti-correlated with Y<sub>P</sub>, members of this school obviously thought that the only reliable measures of D in almost primeval systems, like high-redshift, low-metallicity absorbers, were those leading to high D abundances. They thus thought that the natural evolution of deuterium with time is that connecting the SCHR94 value with the local ones and sketched by the solid line in Fig.2, i.e. a D destruction by one order of magnitude from the primordial to the present abundance.
To obtain such a high D destruction during the Galaxy evolution, one must invoke a high star formation rate (SFR), which is usually assumed to occur at the earliest epochs, because all the observational evidences are against high SFR at relatively recent times. These high SFRs (and their related metal enrichment), inevitably imply a large overproduction of the heavy elements with respect to the observed stellar abundances, unless compensated by mechanisms able to reduce the excess of metals, by diluting or removing them from the Galaxy. For this reason, models with high D destruction usually invoke infall of metal poor gas and galactic winds powered by supernovae explosions, sometimes coupled with variations in the initial mass function. There have been several attempts to find viable Galactic models with strong deuterium depletion, but no scenario consistent with all the Galactic data has been found. For instance, in their pioneering work, Vangioni-Flam & Audouze (1988) concluded that they excessively overproduced the metals, and Scully et al. (1997), in order to obtain the desired D without overproducing the metals, ended up with a present local SFR at least one order of magnitude lower than observed. Tosi et al. (1998) have tested all the possible combinations of the various parameters (SFR, infall, winds, etc.) and have always found significant inconsistencies in the models with high D destruction: metal overabundance with wrong galactocentric distribution, or metallicity distribution of the G-Dwarfs in the solar neighbourhood completely at odds with the observed one, or abundance ratios in halo or disk stars different from the observed ones (e.g. \[O/Fe\] vs \[Fe/H\]), or SFR inconsistent with the observed range. In no way have we been able to find a fairly self-consistent model with high D destruction.
The Galactic School works instead on chemical evolution models able to reproduce as well as possible the largest set of observed Galactic features. Thanks to the improvements both on the observational and on the theoretical sides, good chemical evolution models of the Milky Way nowadays can reproduce the average distribution of the following list of observed features (see e.g. Tosi 1996 and 2000, Boissier & Prantzos 1999 for references):
$``$ current distribution with Galactocentric distance of the SFR,
$``$ current distribution with Galactocentric distance of the gas and star densities,
$``$ current distribution with Galactocentric distance of element abundances as derived from HII regions and from B-stars,
$``$ distribution with Galactocentric distance of element abundances at slightly older epochs, as derived from PNe II,
$``$ age-metallicity relation not only in the solar neighbourhood but also at other distances from the center,
$``$ metallicity distribution of G-dwarfs in the solar neighbourhood,
$``$ local Present-Day-Mass-Function (PDMF),
$``$ relative abundance ratios (e.g. \[O/Fe\] vs \[Fe/H\]) in disk and halo stars.
When one compares with each other all the models in better agreement with these data (e.g. Tosi 1996), the striking result is that they all predict essentially the same deuterium evolution, in spite of the fact that they are computed by different people, with different assumptions on the input parameters and with different numerical codes. The bottom panel of Fig.2 shows an updated version of the comparison: the plotted models are from Galli et al. (1995a, short-dashed line), Dearborn, Steigman, & Tosi (1996, solid line), Chiappini & Matteucci (1996, long-dashed line) and Boissier & Prantzos (1999, dotted line). All the shown curves fit very well the average abundances derived for the local ISM, the pre-solar nebula and the high-redshift absorbers by B&T98. They all show a fairly moderate (a factor from 1.5 to 3, at most) D destruction during the Galaxy lifetime, and therefore suggest that the primordial D abundance should be low: 2$``$(D/H)$`{}_{P}{}^{}\times 10^54`$.
This homogeneity of predictions is not a chance effect, but the consequence of the circumstance that all these models fit equally well the observational data on the present SFR, gas and mass densities, and chemical abundances, which necessarily implies that they predict similar fractions of pristine gas and, therefore, of surviving primordial deuterium.
Our current knowledge on the Galactic evolution of D can thus be summarized as follows: Models predicting high deuterium destruction cannot account for all the observed Galactic properties; models able to reproduce the largest set of Galactic properties all predict low deuterium destruction and, hence, low primordial D.
## 3. <sup>3</sup>He evolution
<sup>3</sup>He has a more complex evolution than D, because it is produced not only during the Big Bang but also inside stars, during the main sequence phase. This early stellar production may be however largely compensated by further nuclear processing in subsequent phases. Standard stellar nucleosynthesis studies predict that, at the end of the star life, the <sup>3</sup>He present in the initial stellar composition is significantly destroyed in massive stars, but preserved or even strongly enhanced in lower mass stars, and that the <sup>3</sup>He net yield is a steeply decreasing function of the stellar initial mass, with a large net production in stars below 2–2.5 M (see e.g the monothonic curves in Fig.3). This behaviour was known since the late sixties (e.g. Iben 1967), and already in 1976 Rood, Steigman, & Tinsley noticed that it leads to overproduce the solar abundance. Only in the mid-nineties, however, with the advent of more detailed combinations of <sup>3</sup>He yields with galactic chemical evolution models (Vangioni-Flam, Olive & Prantzos 1994, Galli et al. 1995a, Dearborn, Steigman & Tosi 1996) it became apparent that the results on <sup>3</sup>He of standard nucleosynthesis studies are definitely inconsistent both with the solar and with the ISM observed abundances. This inconsistency is found with any type of Galactic evolution models, including those in agreement with all the other observational constraints (e.g. Tosi 1996) and was emphasized by several groups at the Elba meeting on the Light Element Abundances in 1994 (Cassé & Vangioni-Flam 1995, Galli et al. 1995b, Tosi, Steigman & Dearborn 1995). In that occasion, Michel Cassé concluded with what has been the most popular refrain on <sup>3</sup>He ever since: <sup>3</sup>He delendum est, like the city of Carthago for the ancient Roman M.P. Cato Censor.
The most probable solution to the <sup>3</sup>He problem is less drastic than that applied to Carthago by the Romans and was proposed already in 1995 (Charbonnel 1995, Hogan 1995). It consists in the further <sup>3</sup>He processing into heavier elements favoured by an extra-mixing occurring in the red giant phase of low-mass stars (see both Charbonnel and Sackman, this volume). When low-mass stars are assumed to experience this extra-mixing and the so-called Cool Bottom Processing (CBP), Galactic evolution models do not overproduce <sup>3</sup>He anymore and fit well the observed solar and HII region abundances (Tosi 1996). The question is: in what fraction of low-mass stars CBP should occur to best fit all the data, taking into account that Bania, Rood et al. (this volume) measure in a few PNe a high <sup>3</sup>He perfectly consistent with the predictions of standard stellar nucleosynthesis (Fig.3) ? Galli et al. (1997) showed that the fraction should be larger than 80$`\%`$ to fit the <sup>3</sup>He abundances observed in the solar system (Geiss & Gloeckler 1998), in PNe and in HII regions (Rood et al. 1995 and this volume).
Fig.4 shows the predictions of the best of models Tosi-1 (see Tosi 1988, Dearborn et al. 1996) when 0% (dotted line in both panels), 90$`\%`$ (solid lines), or 100$`\%`$ (short-dash-dotted lines) of stars with M $``$ 2.5 M are assumed to follow Sackman & Boothroyd’s (1999) prescriptions for CBP <sup>3</sup>He depletion. For the remaining low-mass stars, as well as for all the intermediate and high-mass ones, the <sup>3</sup>He yield is taken from Dearborn et al. (1996). The dotted, solid and short-dash-dotted lines correspond to models assuming as initial abundances (D/H)$`{}_{p}{}^{}=3\times 10^5`$ and (<sup>3</sup>He/H)$`{}_{p}{}^{}=1\times 10^5`$. The dashed lines show the predictions of the same model with 90$`\%`$ CBP, when only the initial D is changed to (D/H)$`{}_{p}{}^{}=10\times 10^5`$, while the long-dash-dotted lines correspond to (D/H)$`{}_{p}{}^{}=20\times 10^5`$.
The vertical bars in the left hand panel represent the ranges of <sup>3</sup>He abundances (at 2$`\sigma `$) derived by Geiss & Gloeckler (1998) and Gloeckler & Geiss (1998) for the Protosolar and the Local Interstellar Clouds, here assumed to be representative of the local ISM, 4.5 Gyr ago and now, respectively. The data points in the right hand panel show the <sup>3</sup>He abundances derived by Rood et al. (1995) from HII region radio observations. It is apparent that the models assuming 90$`\%`$ and 100$`\%`$ of low-mass stars with CBP fit quite well all the data when the initial D is sufficiently low. The CBP depletion is however insufficient to compensate the <sup>3</sup>He overproduction if the initial D/H, subsequently turned into <sup>3</sup>He, is higher than a few 10<sup>-5</sup>, in which case, first the observed protosolar abundance, and then also the local ISM one, cannot be reproduced any more. This is a further argument in favour of the low primordial deuterium resulting from the previous section and from Tytler’s (this volume) discussion of the observations at high redshift.
Hence, if (D/H)$`{}_{p}{}^{}3\times 10^5`$, the <sup>3</sup>He problem is solved if 90% of low-mass stars burn it during the extra-mixing occurring in their red giant phase. In fact, we can simultaneously reproduce the low <sup>3</sup>He abundances of the solar region and of HII regions at any Galactocentric distance, and the high abundance of NGC 3252 and the other PNe measured by Rood et al., which would consequently be associated to the remaining fraction (10%) of stars without deep mixing.
Since the deep mixing depletes not only <sup>3</sup>He, but also the <sup>12</sup>C/<sup>13</sup>C ratio (see Charbonnel and Sackman, this volume), it is important to check the self-consistency of the solution by comparing the model predictions with the carbon isotopic ratio. Charbonnel and do Nascimento (1998) find indeed that more than 90% of 191 field and cluster red giants present carbon ratios significantly lower than the <sup>12</sup>C/<sup>13</sup>C=25 predicted by standard nucleosynthesis. What we also want to check are the predictions of chemical evolution models. This has been done by Palla et al. (2000, hereinafter PBSTG) with a two-folding approach: a) we have compared the available observational data on the carbon isotopic ratio with the corresponding predictions of chemical evolution models assuming the deep mixing in various percentages of low-mass stars; b) we have observed <sup>12</sup>C and <sup>13</sup>C in 28 PNe in mm-waves and compared the derived ratios with those predicted by stellar nucleosynthesis.
Fig.5 shows what model Tosi-1 predicts for the carbon ratio when the <sup>12</sup>C and <sup>13</sup>C adopted yields are from Boothroyd & Sackman (1999) for low-mass stars with CBP, from Marigo (2000) for low and intermediate-mass stars without CBP, and from Limongi, Chieffi & Straniero (2000) for massive stars. Equivalent results are described by PBSTG for stellar yields from other sources. The dotted line shows that without extra-mixing in low-mass stars the <sup>12</sup>C/<sup>13</sup>C ratio is overpredicted with respect to both the abundances observed in the sun and in molecular clouds (assumed to be representative of the present disk abundances). Vice versa a good agreement is achieved if the fraction of stars with CBP is as high as possible (recall that one cannot assume 100% because of the few PNe with high <sup>3</sup>He). As discussed by PBSTG, the amount of predicted <sup>12</sup>C and <sup>13</sup>C strongly depends not only on the extra-mixing assumptions but also (mostly) on the assumptions for the nucleosynthesis in intermediate-mass stars, which are the major contributors to the ISM enrichment of the carbon isotopes. However, we can safely conclude that the observed carbon ratios are always better reproduced by models adopting high percentages of low-mass stars with CBP.
This is also supported by the comparison of the <sup>12</sup>C/<sup>13</sup>C derived by PBSTG for the PNe where <sup>13</sup>C was actually measurable with the carbon ratio predicted for stars right before the ejection of the PN by various nucleosynthesis studies. The left panel of Fig.6 shows that most of the data points (triangles with associated error on the progenitor mass estimate) present carbon ratios lower than those expected from standard nucleosynthesis. The right hand panel shows that the measured carbon ratios are consistent with the predictions of the CBP models at the end of the red giant phase (unfortunately, no nucleosynthesis models are available yet up to the pre-PN phase, and one cannot perform a more appropriate comparison with the PNe observed ratios).
Hence, with deep mixing in $``$90% of low-mass stars one can reproduce the abundances of <sup>3</sup>He observed in the sun, in the ISM and in PNe and the <sup>12</sup>C/<sup>13</sup>C measured in the sun, in red giants, in the ISM and in PNe. We can then conclude that this mechanism appears to be a very promising process, which needs to be further investigated, both to individuate its possible causes and to check its effects of later stellar evolution phases.
Our current knowledge on the Galactic evolution of <sup>3</sup>He can be summarized as follows: All its available observational abundances can be explained if a) its primordial abundance is low, (<sup>3</sup>He/H)$`{}_{p}{}^{}1\times 10^5`$, b) the deterium primordial abundance is also low, and c) deep mixing occurs in almost all low-mass stars.
### Acknowledgments.
Most of what has been described in this review results from the collaborations with C. Chiappini, D. Galli, F. Matteucci, F. Palla, L. Stanghellini and G. Steigman, and I thank them all for their help. Interesting conversations with C. Charbonnel and N. Prantzos have also been very useful. I thank the organizers for such an enjoyable and successful symposium. This work has been partially supported by the Italian COFIN98-MURST at Arcetri.
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warning/0001/astro-ph0001080.html | ar5iv | text | # X-ray and optical observations of three clusters of galaxies: Abell 901, Abell 1437, and Abell 3570Partly based on observations collected at the European Southern Observatory
## 1 Introduction
Clusters of galaxies are excellent probes for cosmological theories. Different cosmological models predict vastly different properties for the clusters as a function of redshift. Therefore relations of cluster properties, their distributions and evolution can be used to constrain cosmological parameters. Selection according to X-ray luminosity is a good way to find the highest mass concentrations because of a relatively well defined correlation between the X-ray luminosity and the total cluster mass (Reiprich & Böhringer 1999; Schindler 1999). A second advantage of high luminosity clusters is, of course, that they require shorter exposure times than their fainter counterparts to get the same numbers of photons. Therefore many X-ray studies have concentrated so far on luminous clusters.
But also low-luminosity clusters are interesting for different reasons. For example, low-luminosity clusters are the links between massive clusters and groups of galaxies. Groups differ in various properties from clusters. Their $`L_XT`$ relation is steeper than the one for clusters (compare Ponman et al. 1996 with Arnaud & Evrard 1999), their gas mass fraction is smaller (compare Pildis et al. 1995 and Ettori & Fabian 1999 or Mohr et al. 1999) and their silicon abundance is lower (Fukazawa et al. 1998). To investigate the physics causing these differences it is important to understand how these properties change with mass or X-ray luminosity.
Low-luminosity clusters are also needed to determine the slope of various relations well. With a long leverage the slope of a relation can be determined more accurately, therefore the luminosities should preferably span several orders of magnitude. For other studies, e.g. the determination of the correlation function or distribution of properties, large samples of clusters are desirable to get the most accurate results. Therefore, when large numbers are required also faint clusters must be taken into account.
Nearby low-luminosity clusters sometimes provide problems. Not only is the countrate determination more sensitive to the background subtraction, but also contamination by other fore- and background sources can become critical, in particular when the cluster shows asymmetries and substructure. The degree of source confusion affecting morphological studies and countrate determinations depends on the sensitivity and the resolution of the detector. Already the factor of 5 between the point spread function (PSF) of the ROSAT/PSPC and the ROSAT/HRI can make a big difference, in particular for low-flux clusters (compare the ROSAT/PSPC observation of Cl0939+4713 (Schindler & Wambsganss 1996) and the corresponding ROSAT/HRI observation (Schindler et al. 1998)). In the RASS (Voges et al. 1996) the spatial resolution is even more critical because the source is not always in the centre of the field-of-view, where the detector has the best spatial resolution, like in most pointed observations. But the source is observed at different off-axis angles as the detector scans the sky. The radius including 50% of the photons increases by about a factor of 13 when moving from the centre to the border of the ROSAT/PSPC. As the final image of a source is composed of many PSFs the final 50% radius in the RASS is about 80$`\mathrm{}`$, i.e. about six times larger than the on-axis PSPC resolution (depending also on the energy).
Clusters which have low fluxes because they are far away, are also very interesting, in particular for cosmological applications. But their extent in the sky is much smaller than for nearby clusters and thus they are less affected by source confusion. Therefore, we concentrate this investigation on nearby clusters with low to intermediate luminosity.
Three clusters – A901, A1437 and A3570 – were selected for a detailed look with the high resolution of the ROSAT/HRI and for optical observations. Each of the clusters shows irregular structure in the RASS. This investigation does of course not show representative RASS clusters, but it is a “worst-case” study. For RASS clusters with high flux and/or regular shape countrate determination and morphological analyses are straightforward and reliable.
The paper is organised as follows. After a description of the observations (Sect. 2) we present the analysis of the optical data in Sect. 3 and the analysis of the X-ray data in Sect. 4. Finally, Sect. 5 gives our summary and conclusions. Throughout this paper we use $`\mathrm{H}_0=50`$ km/s/Mpc and $`\mathrm{q}_0=0.5`$.
## 2 Observations
Three clusters were selected which have low to intermediate X-ray luminosity and show at the same time a very irregular structure in the RASS: A901, A1437, and A3570 (the RASS images will be published by Böhringer et al.). For this analysis X-ray and optical observations were conducted. The X-ray observations were carried out with the ROSAT/HRI (Trümper 1983) and analysed with EXSAS. The optical observations consist of spectra taken with MEFOS (Felenbok et al. 1997) at the ESO3.6m telescope. The analysis of the spectra was carried out with Figaro. The spectra determine the distance of the cluster in the case of A901 and for A3570 a velocity distribution. Details of the optical and X-ray observations are listed in Table 1.
## 3 Optical analysis
For two of the clusters – A901 and A3570 – optical spectra could be taken. In A901 redshifts of three galaxies could be well determined. These three galaxies are at similar redshifts (see Table 2). We derive a cluster redshift of $`z_{A901}=0.17`$. Unfortunately, the number of galaxies is not sufficient to estimate the velocity distribution.
In A3570 we found 17 cluster member galaxies (see Table 2) in addition to the three galaxies known previously (Postman & Lauer 1995). From these 20 galaxies we derive a cluster redshift $`z_{A3570}=0.0375`$ (very similar to the redshift 0.0372 by Abell et al. 1989). The velocity distribution is shown in Fig 1. The velocity dispersion based on these 20 galaxies $`\sigma =460`$km/s is relatively small. Such a small velocity dispersion is a first hint for a virialised state of the cluster. If merging of subclusters would take place and not all of the mergers would happen perpendicular to the line-of-sight a broad velocity distribution would be expected. Assuming virial equilibrium we estimate the cluster mass from the velocity dispersion and the estimate of the virial radius by Girardi et al. (1998) to $`M_{opt}(A3570)=2.7\times 10^{14}M_{}`$.
## 4 X-ray analysis
### 4.1 Abell 901
A ROSAT/HRI image of a 21$`\mathrm{}`$ region around A901 shows several sources (Fig. 2). A901 turns out to have a very compact structure contrary to the previous conclusion from the RASS. Ebeling et al. (1996) list A901 as a double cluster. Their so-called “brighter subcluster” turns out to be a number of point-like sources, while the true cluster emission is what they call the “fainter subcluster”. In Fig. 2a six X-ray point sources with a signal-to-noise ratio of at least 3 are indicated. For three of them optical counterparts can be found on ROE/NRL COSMOS finding charts (see Fig. 2b and Table 4). Unfortunately, for none of these counterparts optical spectra could be taken because the optical observations were carried out before the X-ray observations. At the position of A902 ($`2\mathrm{}`$ West of F) no X-ray emission can be detected.
The optical centre of A901 determined by Abell et al. (1989) is located between the X-ray position of A901 and source A. To test the extent of the X-ray emission of A we derive an X-ray profile (see Fig. 4a). The profile is compared with the on-axis PSF of the ROSAT/HRI. The profile of A is only slightly more extended than the on-axis PSF, which is expected for a point source 6$`\mathrm{}`$ away from the pointing position. Therefore, we conclude that most likely the X-ray emission from A is point-like and therefore not (sub-)cluster emission, but probably emission from an active nucleus in the centre of a galaxy. The most likely candidate for this AGN is a galaxy of 16.6<sup>m</sup> in B (see Table 4). This galaxy is located at a distance of 5$`\mathrm{}`$ from the cluster emission – a distance smaller than the pointing accuracy of ROSAT.
The only cluster emission is coming from the region indicated by “A901” in Fig. 2a ($`\alpha _{2000}=09^\mathrm{h}55^\mathrm{m}57.0^\mathrm{s}`$, $`\delta _{2000}=09\mathrm{°}58\mathrm{}59\mathrm{}`$). This emission is shown magnified in Fig. 3 on a scale of 1.4$`\mathrm{}`$. The emission is very compact, but not point-like, as can be clearly seen from the comparison of the cluster profile and PSF (see Fig. 4a). A $`\beta `$-model fit to the cluster profile (Cavaliere & Fusco-Femiano 1976; Jones & Forman 1984) reveals an extremely small core radius of $`0.10\pm 0.03\mathrm{}`$ or $`22\pm 5`$kpc (see also Table 5) reflecting the compactness of the emission.
The X-ray emission of A901 can be traced out to a radius of almost 2$`\mathrm{}`$, corresponding to 430 kpc. Within this radius a countrate of $`0.059\pm 0.002`$ cts/s is found. If the emission of A901 and the 6 point sources is summed up, the total countrate is at least a factor 1.8 larger than the cluster countrate, i.e. the cluster countrate would be largely overestimated if the point sources were not resolved. For the flux and luminosity shown in Table 5 only the cluster emission of A901 was used.
Estimating a temperature of 4 keV from $`L_XT`$ relations (Allen & Fabian 1998; Markevitch 1998; Arnaud & Evrard 1999) and assuming hydrostatic equilibrium we estimate the total mass at the outer radius $`M_{tot}(r<430`$kpc$`)9\times 10^{13}(\frac{T}{4\mathrm{k}\mathrm{e}\mathrm{V}})M_{}`$. The gas mass is $`M_{gas}(r<430`$kpc$`)=1.2\times 10^{13}M_{}`$, i.e. the gas mass fraction is about 13%.
Obviously, A901 with a flux of $`f_X(0.1`$-$`2.4`$keV$`)=3.0\times 10^{12}`$erg/s/cm<sup>2</sup> is falsely in the RASS X-ray brightest Abell cluster sample of Ebeling et al. (1996) as this sample has a flux limit of $`5\times 10^{12}`$erg/s/cm<sup>2</sup>. Ebeling et al. list a flux of $`5.2\times 10^{12}`$erg/s/cm<sup>2</sup> for the “brighter subcluster”, which is in reality not cluster emission. For the “fainter subcluster”, which is the true A901 emission, they list correctly $`3.0\times 10^{12}`$erg/s/cm<sup>2</sup>, but this value is far below their flux limit.
The compact (but not point-like) nature of the X-ray emission ($`r_c=22`$ kpc) could be interpreted as emission from a galaxy or from a group of galaxies. But a comparison of X-ray luminosity and blue luminosity of the central galaxy (16<sup>m</sup> in B) shows that A901 lies far above the $`L_XL_B`$ relation for early-type galaxies found by Eskridge et al. (1995) and Irwin & Sarazin (1998). A group of galaxies can also be excluded, not only because the X-ray luminosity is too high, but also from the gas mass fraction. The gas mass fraction of 13% is typical for a normal cluster (Ettori & Fabian 1999, Schindler 1999), and would be too high for a group of galaxies (e.g. Pildis et al. 1995). An estimate of the central cooling time yields about $`t_{cool}10^9`$ years. Therefore it is possible, that the compact X-ray emission is caused by a cooling flow.
### 4.2 Abell 1437
The cluster A1437 at a redshift z=0.1339 (Struble & Rood 1987) is the most X-ray luminous cluster of this sample. The cluster centre in X-rays (see Table 5) does not coincide with the optical position: Abell et al. (1989) determined a position 45$`\mathrm{}`$ in the SE of the X-ray maximum. The emission of the cluster is strongly elongated in SW-NE direction (see Fig. 5). This elongation can be seen as well in the RASS. The RASS distinguishes also easily the point source in the NE, for which an optical counterpart can be found on APM finding charts (see Table 4). Although the cluster shape is not exactly elliptical, we fit ellipses to the isophotes (Bender & Moellenhof 1987) to estimate the elongation. The position angle varies around $`55^{}`$ (N over E). The minimum axis ratio of 0.38 is reached at 0.01 cts/s/arcmin<sup>2</sup>. At this level the centre of the ellipse is shifted 35$`\mathrm{}`$ to the west and 36$`\mathrm{}`$ to the south with respect to the position of the X-ray maximum.
The fit parameters of the surface brightness profile are not well constrained (see Table 5) because of the non-spherical morphology of the cluster. Therefore, radial profiles of the cluster emission are determined in four different sectors using as centre the X-ray maximum listed in Table 3 and subsequently fitted with $`\beta `$ models (see Fig. 4b). The two central bins show some excess emission. This excess cannot come from a cooling flow because the central cooling time is about $`2\times 10^{10}`$ years. It is probably a small contamination by an AGN. Because of this excess we try to fit the overall profile with and without these two bins. The results are listed in Table 3. In both – the fit parameters and the fit curves – it is obvious, that the cluster is very asymmetric.
Such asymmetries can arise during a merger of subclusters. From combined N-body and hydrodynamic simulations it is known that such elongated morphologies are common shortly after the collision of two subclusters, when the intra-cluster gas is squeezed out perpendicular to the collision axis (Schindler & Müller 1993).
The X-ray emission can be traced out to about 9$`\mathrm{}`$. After excluding the point source in the NE a countrate of $`0.25\pm 0.01`$ cts/s is found. This corresponds to a flux of $`f_X(0.1`$-$`2.4`$keV$`)=(1.04\pm 0.03)\times 10^{11}`$erg/s/cm<sup>2</sup>. For A1437, which is the most luminous cluster of this sample, the flux determination from the RASS ($`f_X(0.1`$-$`2.4`$keV$`)=1.02\times 10^{11}`$erg/s/cm<sup>2</sup> (Ebeling et al. 1996) and $`f_X(0.1`$-$`2.4`$keV$`)=1.00\times 10^{11}`$erg/s/cm<sup>2</sup> (Ebeling et al. 1998), respectively) is reliable. Also the morphological determination from the RASS is good: the point source in the NE can be distinguished easily and the elongated shape of the cluster is visible in the RASS as well.
### 4.3 Abell 3570
The cluster A3570 is the nearest cluster of this sample ($`z=0.037`$). The X-ray emission is faint and the extent is of the same order as the field-of-view of the ROSAT/HRI. With small smoothing the cluster X-ray emission is hardly visible, because the region is dominated by discrete sources (see Fig. 6a and Table 4). One of the sources (D) is not point-like but has a small extent. This source can be identified with the galaxy ESO 325 - G016 – a cluster galaxy at redshift of $`z=0.03795`$ (Postman & Lauer 1995).
To make the cluster emission visible we remove all point sources, which have a signal-to-noise ratio of at least 3 above the surrounding cluster emission, by fitting a warped surface to the pixels surrounding the point source region and apply a much coarser smoothing (see Fig. 6b). The cluster emission is extended and regular. There is no significant sign of subclustering or merging, i.e. the complex structure seen in the RASS disappears on resolving the discrete sources. Therefore, A3570 is very likely a relaxed cluster. Fitting the profile for this cluster is not possible because the profile is so shallow.
Because of the large extent of the cluster the image had to be vignetting corrected for the countrate determination (Snowdon 1998). The countrate determination is difficult, because the cluster emission fills probably the whole field-of-view of the HRI. We estimate the countrate to be $`1.0_{0.7}^{+0.4}`$ cts/s by counting all the photons within $`r=15\mathrm{}`$ (corresponding to 250 kpc), excluding the discrete sources and using a standard ROSAT/HRI background. Out to this radius we can clearly trace the X-ray emission, but probably the cluster extends further beyond the field-of-view. Therefore, we estimate the upper limit of the countrate by adding the photons found beyond this radius north and east of the cluster and assume the same number in the south and west, which is not covered by the detector. The lower limit is obtained by using the background at the border of the field-of-view, which is a very conservative estimate. For flux and luminosity see Table 5.
The discrete sources change the cluster morphology drastically by feigning substructure in the RASS image. But they do not contribute significantly to the countrate. We estimate the countrate of the discrete sources by fitting a warped surface to the pixels surrounding the point source region and subtract these fitted counts from the original counts. The discrete sources together have a very small count rate of about 0.01 cts/s, which is negligible compared to the cluster emission.
For the X-ray luminosity $`L_X(0.1`$-$`2.4`$keV$`)=(3.2_{2.3}^{+1.1})\times 10^{44}`$erg/s the velocity dispersion $`\sigma =46`$0 km/s is relatively low. While a temperature of 4 keV is consistent with the $`L_XT`$ relations, the $`\sigma T`$ relations predict only 2 keV (White et al. 1997; Mushotzky & Scharf 1997, Wu et al. 1999). The small velocity dispersion confirms the conclusion from the ROSAT/HRI observation, that A3570 is a regular, non-merger cluster.
## 5 Conclusions
We analysed three clusters of galaxies with low to intermediate X-ray luminosities which show an irregular appearance in the ROSAT All-Sky Survey. As the confusion with fore- and background sources is increasingly critical with decreasing flux and increasing substructure of the cluster, we test the limitations of the RASS. With only three clusters, of course, this cannot be a statistical study, but it is meant to be a “worst-case” study. We would like to stress that for higher fluxes source confusion is less important and that for clusters with these fluxes the RASS results are very reliable.
The results of the follow-up observations – X-ray observations with the ROSAT/HRI and optical spectroscopic observations – are summarized in Table 5. The ROSAT/HRI observations, which have a much better spatial resolution than the RASS, revealed in two of the three cases a different morphology than seen in the RASS. The reasons are point sources which could not be resolved in RASS and were therefore confused with the cluster emission. In one of the three clusters the countrate measurement is affected by the point sources, in the two other cases the countrate measurements of the RASS are reliable. The results for the three clusters are the following:
* Abell 901: With a flux of $`f_X(0.1`$-$`2.4`$keV$`)=3.0\pm 0.1\times 10^{12}`$erg/cm<sup>2</sup>/s this cluster is the faintest one in the sample. Both the measurement of the countrate and the determination of the morphology provide difficulties in the RASS. Several discrete sources cannot be distinguished from the cluster emission in the RASS. With the RASS information only it was concluded that this cluster consists of two subclumps (Ebeling et al. 1996), while the true cluster X-ray emission is very compact ($`r_c=22`$kpc). The countrate would be overestimated by at least 80% if the discrete sources cannot be separated from the cluster.
* Abell 1437: This cluster has the highest flux of all the clusters in the sample $`f_X(0.1`$-$`2.4`$keV$`)=1.04\pm 0.03\times 10^{11}`$erg/cm<sup>2</sup>/s. With such a high flux the RASS results are very reliable. There is no problem with the countrate or morphology determination. The point source in the NE of the cluster can easily be distinguished in the RASS.
* Abell 3570: This cluster has several point sources superposed on the cluster emission. Although the countrate determination is not affected by these sources (their countrate is only 1% of the cluster countrate) the true morphology is very regular in contrast to the impression from the RASS. Optical measurements of the cluster galaxies give a very small velocity dispersion confirming the picture of a relaxed cluster.
We conclude that for clusters with fluxes less than a few times $`10^{12}`$erg/cm<sup>2</sup>/s, which have at the same time an irregular morphology, the confusion with fore- and background sources can be a problem in the RASS. So morphological analyses of RASS clusters tend to overestimate the fraction of clusters with substructure.
A new X-ray all-sky survey with a potential second ABRIXAS (Friedrich et al. 1996) mission would have several advantages to push this source confusion limit down.
* The spatial resolution of ABRIXAS is about 2.5 times better than the resolution of the RASS, which makes it easier to distinguish between point and extended sources.
* The energy range of ABRIXAS (0.5-12 keV) is much better suited for clusters than the ROSAT range (0.1-2.4 keV). Although both surveys have about the same sensitivity at 1 keV, ABRIXAS with its harder energy range would detect 3-5 times more photons from a standard cluster.
* The spectral resolution of ABRIXAS is better than that of the ROSAT/PSPC. Together with the wider energy range of ABRIXAS this provides an improved possibility to separate hard cluster emission from soft foreground sources.
###### Acknowledgements.
I thank Chris Collins for introducing me into the secrets of optical spectra, Hans Böhringer for making the RASS images of the three clusters available to me, Peter Friedrich for providing the ABRIXAS numbers, and Phil James and Joachim Wambsganss for carefully reading the manuscript. It is a pleasure to thank Carlo Izzo for his most helpful EXSAS support. I acknowledge gratefully the hospitality of the Institut d’Estudis Espacials de Catalunya in Barcelona. During the stay there I was supported by the TMR grant ERB-FMGE CT95 0062 by CESCA-CEPBA. |
warning/0001/astro-ph0001075.html | ar5iv | text | # The first galaxies: clues from element abundances
## 1 Introduction
The aim of this talk is to consider the information on the first episodes of star formation in the universe provided by studies of element abundances at high redshift. This is very much a growth area at present. Thanks largely to the new opportunities offered by the Keck telescopes and their VLT counterparts in the southern hemisphere, we find ourselves in the exciting position of being able, for the first time, to measure the abundances of a wide range of chemical elements directly in stars, H II regions, cool interstellar gas and hot intergalactic medium, all observed when the universe was only $`1/10`$ of its present age. Our simple-minded hope is that, by moving back to a time when the universe was young, clues to the nature, location, and epoch of the first generations of stars may be easier to interpret than in the relics left today, some 11 Gyrs later. Furthermore, as we shall see, the metallicities of different structures in the universe and their evolution with redshift are key factors to be considered in our attempts to track the progress of galaxy formation through the cosmic ages. As can be readily appreciated from inspection of Figure 1, our knowledge in this field is still very sketchy. Given the limitations of space, this review focuses primarily on results obtained in the last year on the three components of the high $`z`$ universe shown in Figure 1.
## 2 Lyman Break Galaxies
Undoubtedly, one of the turning points in extragalactic astronomy in the 1990s has been the realisation that high redshift galaxies can be found in large numbers using a highly efficient photometric selection technique based on the passage of the Lyman edge—at the rest wavelength of 912 Å—through the $`U`$-band. After many years of fruitless searches (targeted mainly to Ly$`\alpha `$ emission which turned out to be a less reliable marker than anticipated), we have witnessed a veritable explosion of data from the Hubble Deep Fields and ground-based surveys. Galaxies with measured redshifts in excess of $`z2.5`$ now number in the many hundreds (the 1000 mark is just around the corner); such large samples have made it possible to trace the star formation history of the universe over most of the Hubble time and to measure large-scale properties of the population, most notably their clustering and luminosity functions (Madau et al. 1996; Steidel et al. 1998, 1999 and references therein).
However, for a quantitative study of many of the physical properties of Lyman break galaxies, even the light-gathering power of the world’s largest telescopes is not enough and we have to rely on gravitational lensing to boost the flux to levels where their spectra can be recorded with sufficiently high resolution and signal-to-noise ratio. This is the case for the $`z=2.73`$ galaxy MS 1512-cB58, an $`L^{}`$ Lyman break galaxy fortuitously magnified by a factor of $`30`$ by the foreground cluster MS 1512+36 at $`z=0.37`$ (Yee et al. 1996; Seitz et al. 1998). Somewhat ironically, our Keck spectrum of cB58 (Pettini et al. 2000b) is one of the best examples of the ultraviolet spectrum of a starburst galaxy at any redshift, thanks to the combined effects of gravitational lensing, redshift, and collecting area of the Keck telescopes.
At $`z3`$ optical wavelengths correspond to the rest-frame far-UV, where we see the integrated light of short-lived O and early B stars. Such spectra are most effectively analysed with population synthesis models, the most sophisticated of which is Starburst99 developed by the Baltimore group (Leitherer et al. 1999). In Figure 2 we compare Starburst99 model predictions for different IMFs with our data in the region of the C IV $`\lambda \lambda 1548,1550`$ doublet. It is important to realise that the comparison only refers to stellar spectral features and does not include the interstellar lines, readily recognisable by their narrower widths (these IS lines are much stronger in cB58, where we sample the whole ISM of the galaxy, than in the models which are based on libraries of nearby Galactic O and B stars). With this clarification, it is evident from Figure 2 that the spectral properties of at least this Lyman break galaxy are remarkably similar to those of present-day starbursts—a continuous star formation model with a Salpeter IMF provides a very good fit to the observations. In particular, the P-Cygni profiles of C IV, Si IV and N V are sensitive to the slope and upper mass limit of the IMF; the best fit in cB58 is obtained with a standard Salpeter IMF with slope $`\alpha =2.35`$ and $`M_{\mathrm{up}}=100M_{}`$ (top panel of Figure 2). IMFs either lacking in the most massive stars or, conversely, top-heavy seem to be excluded by the data (middle and bottom panels of Figure 2 respectively).
The only significant difference between the observed and synthesised spectrum is in the optical depth of the P-Cygni absorption trough which is lower than predicted (top panel of Figure 2). This is likely to be an abundance effect, since an analogous weakening of the absorption is seen in OB stars with mass loss in the Magellanic Clouds (e.g. Lennon 1999) and is also predicted by stellar wind theory (e.g. Kudritzki 1998). In future, when the libraries of stellar spectra in Starburst99 are expanded to include Magellanic Cloud stars (a project which is already underway), it may be possible to calibrate the optical depth of C IV absorption with Carbon abundance and use this feature to deduce the metallicity of high redshift star-forming galaxies. For the moment, we conclude on the basis of a qualitative comparison that the metallicity of the young stellar population in cB58 is similar to that in the LMC, where \[C/H\]$`0.6`$ . Weak interstellar lines of Sulphur, Silicon, and Nickel are consistent with this abundance estimate.
### 2.1 Moving to the Infrared
Very recently, the successful commissioning of NIRSPEC on Keck II and ISAAC on VLT1 have made it possible to extend spectroscopic studies of Lyman break galaxies to the near infrared which, at $`z3`$, includes the familiar optical emission lines from H II regions on which much of our knowledge of local star-forming galaxies is based. As indicated by exploratory observations with UKIRT (Pettini et al. 1998), detecting these lines is a challenging task even with large telescopes, so that we may be restricted to studying the brightest examples of Lyman break galaxies, with $`L\stackrel{>}{}L^{}`$. Figure 3 shows an example of such data. The relative strengths of \[O III\] and H$`\beta `$ in Q0201–C6 are typical of the dozen or so objects observed so far; we find that generally $`I_{\mathrm{H}\beta }\stackrel{<}{}I_{4959}`$ and $`R_{23}\stackrel{>}{}+0.7`$, where $`R_{23}`$ is the familiar strong line ratio index of Pagel et al. (1979). This in turn implies abundances of $`1/31/6`$ solar; as shown by Teplitz et al. (2000), cB58 conforms to this pattern with \[O/H\] $`0.5`$, in good agreement with the UV analysis discussed above.
### 2.2 Kinematics
The combination of (rest-frame) optical and UV observations gives insights into several other properties of LBGs, apart from chemical abundances. The widths of the emission lines are likely to be better indicators of the underlying masses than the interstellar absorption lines which, being sensitive to very low column densities, can be broadened by gas accelerated to high velocities by supernovae and stellar winds associated with the star-formation activity. A preliminary analysis of the dozen objects in our sample indicates velocity dispersions $`\sigma 60120`$ km s<sup>-1</sup>; $`\sigma 80`$ km s<sup>-1</sup> seems to be typical.
In Figure 4 we compare these values with analogous data for nearby galaxies. $`W_{20}`$ is full width at 20% of the peak intensity ($`W_{20}=\sigma \times 3.62`$ for a Gaussian profile) which in the Lyman break galaxies we measure most accurately from \[O III\] $`\lambda 5007`$, and $`M_B`$ is the absolute magnitude in the rest-frame $`B`$-band which at $`z3`$ can be deduced directly from the observed $`K`$-band magnitude without the need for a substantial $`k`$-correction. The horizontal bar in Figure 4 shows the range of values of $`W_{20}`$ for the LBGs observed so far, which are mostly at the bright end of the luminosity function, with $`M_B=22`$ to $`23`$ (as indicated by the vertical bars). The most appropriate comparison is probably with the compilation of \[O II\] widths in local star-forming galaxies (filled dots) by Kobulnicky & Gebhardt (2000) who mimicked the conditions under which these measurements are conducted at high redshift, by obtaining global values of $`W_{20}`$ which refer to a whole galaxy, and did not correct for inclination and internal extinction. As can be seen from Figure 4, the widths of the emission lines in the brighter Lyman break galaxies are at the upper end of the range of values observed locally, and are significantly larger than those of H II galaxies (triangles). Thus the typical $`\sigma 80`$ km s<sup>-1</sup> of LBGs is the value which one may expect from a disk galaxy (viewed at a random inclination) rotating at $`150`$ km s<sup>-1</sup>. Indeed there are hints in two of the objects observed that we may be seeing the rotation curve directly in spatially resolved \[O III\]$`\lambda 5007`$ emission lines. Kinematical masses in excess of a few times $`10^{10}M_{}`$ are indicated.
Also shown in Figure 4 is the Tully-Fisher relation for local spirals. This comparison is less straightforward, because the Tully-Fisher relation is derived from spatially resolved rotation curves corrected for inclination and internal extinction—we have tried to take these factors into account in a statistical sense when reproducing the mean relation by Pierce & Tully (1992) in Figure 4. Vogt et al. (1997) found a mild brightening, by $`\stackrel{<}{}0.4`$ mag, of the relation in a sample of 16 galaxies at $`0.15<z<0.75`$ and interpreted it as being due to luminosity evolution in the field galaxy population (the Vogt et al. data are, like ours, based on observations of \[O II\] and \[O III\] emission lines from H II regions). Taken at face value, the very preliminary comparison in Figure 4 suggests a much more significant luminosity evolution when we look back to $`z3`$, perhaps amounting to as much as 2 magnitudes in the $`B`$-band.
Finally, we find that there are systematic velocity offsets between nebular emission lines, interstellar absorption lines, and Ly$`\alpha `$ in most Lyman break galaxies observed; these offsets can be explained as resulting from large scale outflows with velocities of up to several hundred km s<sup>-1</sup>. In cB58 Pettini et al. (2000b) deduced a mass outflow rate $`\dot{M}60M_{}\mathrm{yr}^1`$, comparable to the rate at which gas is being turned into stars. Such galactic ‘superwinds’ seem to be a common feature of starburst galaxies at all redshifts (see Tim Heckman’s article in this volume), and may well be the mechanism which self-regulates star formation, distributes metals over large volumes and allows the escape of ionizing photons into the intergalactic medium.
In summary, all the available information is consistent with the notion that Lyman break galaxies are already well developed systems at $`z3`$, with stellar populations, chemical abundances and kinematics very much in line with those of the more massive star-forming galaxies in the local universe. As explained above, the best studied examples so far are all at the bright end of the luminosity function; thus, perhaps we should not be surprised to find that their properties are relatively uniform. What is interesting is that galaxies at such an advanced stage of evolution were already in place at the relatively early epochs corresponding to $`z=34`$; it therefore seems most natural to associate these objects with the progenitors of today’s elliptical galaxies and bulges of spirals, as proposed by Steidel et al. (1996).
## 3 Damped Lyman alpha Systems
These are the absorption systems with the highest column density of neutral hydrogen, $`N`$(H I)$`2\times 10^{20}`$ cm<sup>-2</sup>, seen in the spectra of QSOs and they provide us with the best opportunity to measure accurately the abundances of a wide range of elements at high redshift. The reason is simple: QSOs can be several hundred times brighter than Lyman break galaxies at the same redshift. Surveys with HIRES on Keck I have produced data of exquisite quality—a 10% accuracy in the determination of interstellar gas-phase abundances is achievable with only modest efforts (Lu et al. 1996; Prochaska & Wolfe 1999).
This makes it all the more frustrating that a clear connection between DLAs and galaxies has yet to be established. In principle, selecting galaxies by their H I absorption cross-section should provide a more representative sampling of the field population at a given redshift than conventional magnitude limited surveys, either in the continuum or emission lines. Thus one may conjecture that DLAs pick out galaxies from the whole luminosity function, particularly if H I cross section has only a mild dependence on galaxy luminosity (Steidel, Dickinson, & Persson 1994), and that Lyman break galaxies may just be the most luminous DLAs at $`z3`$ (Pettini et al. 2000b; Steidel, Pettini, & Hamilton 1995; Djorgovski et al. 1996). While this interpretation is attractive in its simplicity, we must face up to the fact that it cannot readily account for the most recent observations of DLAs (reproduced in Figure 5) which find no significant evolution of either the gas mass or the metallicity (Rao & Turnshek 2000; Pettini et al. 1999; Prochaska & Wolfe 2000) over a redshift interval ($`z0.54`$) during which most of today’s stars were apparently formed (see Figure 7 of Pettini 1999). Possibly, existing samples of damped Ly$`\alpha `$ systems are subject to subtle selection effects of their own and may preferentially trace a particular stage in the evolution of galaxies, when the gas has an extended distribution and only moderate surface density, and the metal—and therefore dust—content is low. There is both theoretical (Mo, Mao, & White 1999) and observational (Le Brun et al. 1997; Rao & Turnshek 1998) evidence in support of this picture.
These latest developments do not detract from our interest in damped Ly$`\alpha `$ systems. First, as emphasised repeatedly by Fall and collaborators (Fall 1996), the column density-weighted mean metallicity of DLAs is the closest measure we have of the degree of metal enrichment reached by the gaseous component of galaxies at a given epoch, irrespectively of the precise nature of the absorbers. Thus, values of \[$``$Z$`{}_{\mathrm{DLA}}{}^{}`$\] at different redshifts (so far most effectively deduced from the abundance of Zn—Pettini et al. 1999) are essential reference points for models of global chemical evolution (Prantzos & Silk 1998; Pei, Fall, & Hauser 1999). The only uncertainty which remains to be resolved for a full use of this information is the degree to which existing samples of DLAs are biased against sight-lines sufficiently dusty to obscure the background QSOs; this is a question which we are in the process of exploring by examining the statistics of damped systems in radio selected QSOs. At redshifts $`z23`$ the metallicity distribution of known DLAs is intermediate between those of stars in the halo and thick disk of the Milky Way; at this epoch most of the galaxies giving rise to damped systems were clearly less evolved chemically than the stellar population forming the thin disk of our Galaxy (see Figure 6).
Second, DLAs present us with the opportunity to extend local studies of the relative abundances of different elements to unexplored regimes and to earlier epochs. Potentially, DLAs have an important role to play here in complementing the information so far obtained from observations of Galactic stars and nearby H II regions and providing fresh clues both to the origin of different stellar populations and to the stellar yields.
For example, Pettini, Lipman, & Hunstead (1995) and Lu, Sargent, & Barlow (1998) showed that in DLAs it is possible to follow the behaviour of the (N/O) ratio to lower metallicities than those probed up to now (IZw18 still remains the most metal-poor star-forming region known in our vicinity). Their results appear to lend support to the idea of a delayed production of primary nitrogen by intermediate mass stars, although this interpretation has been challenged more recently (Centurión et al. 1998; Izotov & Thuan 1999—see also Pilyugin 1999) and more observations are clearly required in order to settle the issue.
The latest application of this technique involves Silicon (an $`\alpha `$-capture element) and Manganese. In Figure 7 (reproduced from Pettini et al. 2000a) the abundances of these two elements in damped Ly$`\alpha `$ systems of different metallicities are compared with analogous data for stars in the disk and halo of our Galaxy. The DLAs considered are those where less than 50% of Si, Mn, and Fe is locked up in dust grains, so that the total (gas+dust) abundances can be recovered with minimum uncertainty. The first-order conclusion is that the DLA values roughly follow the local trends, but there are notable differences too, as we now discuss.
The rise of \[Si/Fe\] from to solar to between $`+0.3`$ and $`+0.4`$ as the metallicity drops to \[Fe/H\]$`=2`$ (top panel of Figure 7) is the well-known overabundance of the $`\alpha `$-elements commonly attributed to the delayed production of additional Fe by Type Ia supernovae. While some DLAs do show enhanced \[Si/Fe\], we also find counter-examples of near-solar abundance of Si at metallicities in the range \[Fe/H\]$`1`$ to $`2`$. Current wisdom would interpret such cases as arising in galaxies where star formation has proceeded slowly, or in bursts, so that there has been sufficient time for Fe to build up to solar abundance relative to Si, while the overall metallicity remained low. Corroborating evidence for this interpretation may be provided by deep imaging of the absorbers, if they are found to be low surface brightness or dwarf galaxies.
Turning to Mn (lower panel of Figure 7), the strong decrease in \[Mn/Fe\] towards low metallicities is now well documented, but its origin is unclear. Two possibilities have been proposed and both have problems with the DLA data, at least at face value. If the stellar trend is due to a metallicity dependent yield of Mn in massive stars, it is difficult to explain the one DLA with \[Mn/Fe\] $`0.3`$ at solar \[Fe/H\]. On the other hand, enhanced (relative to Fe) production of Mn in Type Ia supernovae cannot explain DLAs with low \[Mn/Fe\] and near-solar \[Si/Fe\] of which there are at least two examples.
In concluding this section, it is important to stress the preliminary nature of the above conclusions which are based on the comparison of very few measurements in DLAs with a much larger body of stellar data. One of the lessons from stellar work is that there is considerable scatter, both observational and intrinsic, in the relative abundances of different elements so that most trends only become apparent when a large set of measurements has been assembled. Thus Figure 7 should be taken as no more than an illustration of the issues which can be addressed with surveys of element abundances in damped systems. Although work on element ratios in high redshift galaxies is still a long way behind its counterpart in Galactic stars, it may well play a decisive role in resolving some outstanding questions on the origin of elements.
## 4 The Lyman alpha Forest
The final component of the high redshift universe considered in this review is the all-pervading intergalactic medium which manifests itself as fluctuating absorption bluewards of the Ly$`\alpha `$ emission line of every QSO. Observationally, the term Ly$`\alpha `$ forest is used to indicate the bulk of discrete Ly$`\alpha `$ absorption lines with column densities in the range $`10^{16}\stackrel{>}{}N(\mathrm{HI})\stackrel{>}{}10^{12}`$ cm<sup>-2</sup>; since this gas is highly ionised, it may account for most of the baryons at $`z3`$ (Rauch 1998). Hydrodynamical simulations have shown that the Ly$`\alpha `$ forest is a natural consequence of structure formation in a universe dominated by cold dark matter and bathed in a diffuse ionising background (e.g. Weinberg, Katz, & Hernquist 1998). In this picture, the physics of the absorbing gas is relatively simple and the run of optical depth $`\tau `$(Ly$`\alpha `$) with redshift can be thought of as a ‘map’ of the density structure of the IGM along a given line of sight. At low densities, where the temperature of the gas is determined by the balance between photoionisation heating and adiabatic cooling, $`\tau `$(Ly$`\alpha `$)$`(1+\delta )^{1.5}`$, where the $`\delta `$ is the overdensity of baryons $`\delta (\rho _\mathrm{b}/\rho _\mathrm{b}1)`$. At $`z=3`$ $`\tau `$(Ly$`\alpha `$)$`=1`$ corresponds to a region of the IGM which is just above the average density of the universe at that time ($`\delta 0.6`$).
The lack of associated metal lines was originally one of the defining characteristic of the Ly$`\alpha `$ forest and was interpreted as evidence for a primordial origin of the clouds (Sargent et al. 1980). As it is often the case, subsequent improvements in the observations have shown this to be an oversimplification and in reality weak metal absorption, principally by C IV, is present at the redshift of most Ly$`\alpha `$ clouds down to the detection limit of the data (Songaila & Cowie 1996). The degree of metal enrichment implied is relatively high, (\[C/H\]$`2.5`$ with a scatter of perhaps a factor of $`3`$—Davé et al. 1998), in the sense that stars with significantly lower metallicities are known to exist in the halo of our Galaxy.
It is not easy to understand how the low density IGM came to be polluted so uniformly by the products of stellar nucleosynthesis at such an early epoch. While, as explained above, we see directly the outflow of metal-enriched gas in ‘superwinds’ from Lyman break galaxies at the same redshift, most of this gas is not expected to travel far from the production sites, because it is either trapped by the gravitational potential of the galaxies, if they are sufficiently massive, or is confined by the pressure of the hot IGM (Ferrara, Pettini, & Shchekinov 2000). Whether an early episode of pre-galactic star formation is required depends on whether C IV lines continue to be seen in Ly$`\alpha `$ clouds of diminishing H I column density. Current limits are for $`N`$(H I)$`\stackrel{>}{}3\times 10^{14}`$ cm<sup>-2</sup> (some 75% of such Ly$`\alpha `$ clouds have associated C IV absorption—Songaila & Cowie 1996) corresponding to moderately overdense gas ($`\delta \stackrel{>}{}10`$) which in the simulations is preferentially found in the vicinity of collapsing structures and may thus reflect local, rather than universal, metal pollution.
The detection of C IV lines when $`N`$(H I)$`\stackrel{<}{}1\times 10^{14}`$ is a challenging task, even with a 10-m telescope, because we are dealing with observed equivalent widths $`W_\lambda (1550)\stackrel{<}{}2.5`$ mÅ. A possible approach in these circumstances is to try and recover such a weak signal from a statistical treatment of many lines which individually are below the detection limit. Unfortunately, different analyses have reached conflicting conclusions (Lu et al. 1998; Cowie & Songaila 1998). Furthermore, a recent reappraisal of the techniques with the help of extensive simulations of the spectra has indicated that many subtle effects, such as small random differences between the redshifts of Ly$`\alpha `$ and C IV absorption, make the interpretation of the results far from straightforward (Ellison et al. 1999).
A more direct way to tackle the problem is to push the detection limit further by securing spectra of exceptionally high signal-to-noise ratio; as for Lyman break galaxies this is most effectively achieved with the aid of gravitational lensing. In this way Ellison et al. (2000) were able to reach S/N$`=200300`$ in the C IV region between $`z=2.91`$ and 3.54 of the gravitationally lensed QSO Q1422+231 after adding together data recorded over several nights with HIRES on Keck I.
As can be seen from Figure 8, the number of weak C IV lines continues to rise as the signal-to-noise ratio of the spectra increases and any levelling off in the column density distribution presumably occurs at $`N`$(C IV)$`<5\times 10^{11}`$ cm<sup>-2</sup>. This limit is one order of magnitude more sensitive than those reached previously. In other words, we have yet to find any evidence in the Ly$`\alpha `$ forest for regions of the IGM which are truly of primordial composition or have abundances as low as those of the most metal-poor stars in the Milky Way halo. Pushing the sensitivity of this search even further is certainly one of the goals for the future.
## 5 Conclusions
This review has covered a lot of ground reflecting the fast pace of progress in the study of element abundances at high redshift, now set to accelerate further with the forthcoming availability of new, efficient spectrographs on Keck, VLT, Subaru and Gemini. The picture which is emerging is that of a universe at $`z3`$ with many of today’s characteristics already in place. At this epoch, Lyman break galaxies resembled closely today’s star-forming galaxies dominated by Population I stars, damped Ly$`\alpha `$ systems exhibited mostly Population II chemical abundances, and the low density Ly$`\alpha `$ forest may well have been the repository of the first heavy elements synthesised by Population III stars. This does not necessarily imply a one-to-one correspondence between these objects then and now, given the substantial time interval available for evolution. It is likely that Lyman break galaxies, which at $`z3`$ trace the highest peaks in the underlying mass distribution, have through subsequent mergers evolved into today’s massive ellipticals and bulges to be found preferentially in rich clusters. It is also plausible that the gas giving rise to at least some high redshift damped Ly$`\alpha `$ systems has turned into the stars of today’s spiral galaxies. And the heavy elements ejected into the IGM by the first stars which formed in low mass collapsed structures have by now presumably been augmented by the much more substantial metal-enriched outflows from successive generations of stars in more massive galaxies. Thus the lack of a clear age-metallicity relationship in our own Galaxy is reflected on a much larger scale by the universe as a whole—old (and high redshift) do not necessarily mean metal-poor.
What is clear is that much work still needs to be done before we have a full picture of the chemical enrichment of the universe at $`z3`$. The very substantial gaps in our knowledge evident in Figure 1 are reflected by the results of a simple accounting exercise. As discussed by Pettini (1999) and more recently Pagel (2000), the comoving density of metals so far detected in the Ly$`\alpha `$ forest, damped Ly$`\alpha `$ systems and Lyman break galaxies accounts for only about 10% of the total metal production associated with the star formation activity we see directly at $`z\stackrel{>}{}3`$ . Presumably then, as now, at least some if not most of the ‘missing metals’ are to be found in hot gas—in galactic halos and (proto)clusters—which has not yet been fully accounted for, mainly because we remain ignorant of both its metallicity and baryon content. Food for thought, as we enter the new millennium.
I should like to acknowledge my collaborators in the various projects described in this talk, particularly Chuck Steidel, Sara Ellison, Kurt Adelberger, David Bowen, Len Cowie, Jean-Gabriel Cuby, Mark Dickinson, Mauro Giavalisco, Alan Moorwood, Joop Schaye, Alice Shapley and Toni Songaila. I am grateful to the Royal Society and the organisers for inviting me to take part in such a stimulating discussion meeting. |
warning/0001/math0001170.html | ar5iv | text | # Untitled Document
Deformations of chiral algebras
and quantum cohomology of toric varieties
Fyodor Malikov<sup>1</sup><sup>1</sup>1partially supported by an NSF grant and Vadim Schechtman
Let $`X`$ be a smooth complex variety. It was shown in \[MSV\] that the complex cohomology algebra $`H^{}(X)`$ may be obtained as a cohomology of a certain vertex algebra $`H^{ch}(X)`$ canonically associated with $`X`$. By definition, $`H^{ch}(X)=H^{}(X;\mathrm{\Omega }_X^{ch})`$, where $`\mathrm{\Omega }_X^{ch}`$ is a sheaf of vertex superalgebras constructed in \[MSV\]. (If $`X`$ is compact, then $`H^{ch}(X)`$ may be called the chiral Hodge cohomology algebra of $`X`$.) The algebra $`H^{ch}(X)`$ is equipped with a canonical odd derivation $`Q`$ of square zero, and the cohomology of $`H^{ch}(X)`$ with respect to $`Q`$ is equal to $`H^{}(X)`$.
In the very interesting paper \[B\] Borisov defined for a toric complete intersection $`X`$ a certain vertex superalgebra $`V(X)`$ equipped with an odd derivation of square zero so that $`H^{ch}(X)`$ equals the cohomology of $`V(X)`$ with respect to this derivation. It follows that $`H^{}(X)`$ may also be represented as the cohomology of $`V(X)`$ with respect to another odd derivation $`d`$.
Let $`X`$ be a smooth complete toric variety. In the present note we include Borisov’s algebra $`V(X)`$ and its derivation $`d`$ in a family $`(V_q(X),d_q)`$ of vertex superalgebras with derivation, parametrized by $`qH^2(X)`$, so that the cohomology of $`V_q(X)`$ with respect to $`d_q`$ is equal to the quantum cohomology algebra of $`X`$.
In sect. 2.5 we present a simpler version of this construction in the case of $`^N`$ and apply the deformation technique to compute $`H^{}(^N;\mathrm{\Omega }_^N^{ch})`$
We also get similar (partial) results for Fano hypersurfaces in $`P^N`$.
§1. Borisov’s construction
1.1. Lattice vertex algebras. Let $`L`$ be a free abelian group on $`2N`$ generators $`A^i,B^i,\mathrm{\hspace{0.33em}1}iN`$. Give $`L`$ an integral lattice structure by defining a bilinear symmetric $``$-valued form
$$(.,.):L\times L$$
so that
$$(A^i,B^j)=\delta _{ij},(A^i,A^j)=(B^i,B^j)=0.$$
Introduce the complexification of $`L`$:
$$𝔥_L=L_{}.$$
The bilinear form $`(.,.)`$ carries over to $`𝔥_L`$ by bilinearity. Let
$$\widehat{𝔥_L}=𝔥_L[t,t^1]K$$
be a Lie algebra with bracket
$$[xt^i,yt^j]=i(x,y)\delta _{i+j}K,[xt^i,K]=0.$$
Associated with $`L`$ there is a group algebra $`[L]`$ with basis $`e^\alpha ,\alpha L,`$ and multipliciation
$$e^\alpha e^\beta =e^{\alpha +\beta },e^0=1,\alpha ,\beta L.$$
Denote by $`S_{𝔥_L}`$ the symmetric algebra of the space $`𝔥_Lt^1[t^1]`$. The space $`S_{𝔥_L}[L]`$ carries the well-known vertex algebra structure, see for example \[K\]. Borisov proposes to enlarge this lattice vertex algebra by fermions as follows.
We tacitly assumed that $`𝔥_L`$ is a purely even vector space: $`𝔥_L^{(0)}=𝔥_L,𝔥_L^{(1)}=0`$. Let $`\mathrm{\Pi }𝔥_L`$ satisfy the relations $`\mathrm{\Pi }𝔥_L^{(1)}=𝔥_L,\mathrm{\Pi }𝔥_L^{(0)}=0`$. Thus $`\mathrm{\Pi }𝔥_L`$ is a purely odd vector space with basis $`\mathrm{\Psi }^i,\mathrm{\Phi }^i`$ carrying the following odd bilinear form:
$$(.,.):\mathrm{\Pi }𝔥_L\times \mathrm{\Pi }𝔥_L,$$
$$(\mathrm{\Psi }^i,\mathrm{\Phi }^j)=\delta _{ij},(\mathrm{\Psi }^i,\mathrm{\Psi }^j)=(\mathrm{\Phi }^i,\mathrm{\Phi }^j)=0.$$
Given all this, one defines the Clifford algebra, $`Cl_{𝔥_L}`$, to be the vector superspace
$$Cl_{𝔥_L}=\mathrm{\Pi }𝔥_L[t,t^1]K^{},Cl_{𝔥_L}^{(1)}=\mathrm{\Pi }𝔥_L[t,t^1],Cl_{𝔥_L}^{(0)}=K^{},$$
with (super)bracket $`[xt^i,yt^j]=(x,y)\delta _{i+j}K^{}`$.
Let $`\mathrm{\Lambda }_{𝔥_L}`$ be the symmetric algebra of the superspace
$$_{i=1}^N(\mathrm{\Phi }^i[t^1]\mathrm{\Psi }^it^1[t^1].$$
(If we had been allowed to forget about the parity, we would have equivalently defined $`\mathrm{\Lambda }_{𝔥_L}`$ to be the exterior algebra of the indicated space.) The space $`\mathrm{\Lambda }_{𝔥_L}`$ carries the well-known vertex algebra structure, see for example \[K\].
Finally let
$$V_L=\mathrm{\Lambda }_{𝔥_L}S_{𝔥_L}[L].$$
Being a tensor product of vertex algebras, $`V_L`$ is also a vertex algebra.
1.2. Explicit description of the vertex algebra structure on $`V_L`$. To simplify the notation, we identify $`[L]`$ with the subspace $`11[L]`$. As an $`\widehat{𝔥_L}Cl_{𝔥_L}`$-module, $`V_L`$ is a direct sum of irreducibles and there is one irreducible module, $`V_L(\alpha )`$, for each $`\alpha L`$. $`V_L(\alpha )`$ is freely generated by the supercommutative associative algebra $`S_{𝔥_L}\mathrm{\Lambda }_{𝔥_L}`$ from the highest weight vector $`e^\alpha `$. The words “highest weight vector” mean that the following relations hold:
$$A_n^ie^\alpha =\mathrm{\Psi }_n^ie^\alpha =B_n^ie^\alpha =\mathrm{\Phi }_{n+1}^ie^\alpha =0,n0,$$
$$Ke^\alpha =K^{}e^\alpha =e^\alpha ,xe^\alpha =(x,\alpha )e^\alpha ,x𝔥_L.$$
Thus, $`V_L(\alpha ),\alpha L,`$ are different as $`\widehat{𝔥_L}Cl_{𝔥_L}`$-modules, but isomorphic as $`\widehat{𝔥_L}_1Cl_{𝔥_L}`$-modules, where $`\widehat{𝔥_L}_1\widehat{𝔥_L}`$ is the subalgebra linearly spanned by $`xt^i,i0,x𝔥_L`$. In fact, the multiplication by $`e^\beta `$ provides an isomorphism of $`\widehat{𝔥_L}_1Cl_{𝔥_L}`$-modules:
$$e^\beta :V_L(\alpha )V_L(\alpha +\beta ),xe^\alpha xe^{\alpha +\beta }.$$
Let us now define the state-field correspondence, that is, attach a field $`x(z)\text{End}(V_L)((z,z^1))`$ to each state $`xV_L`$. As has become customary, we shall write $`x_i`$ for $`xt^i`$ ($`x𝔥_L\text{ or }\mathrm{\Pi }𝔥_L`$). We have:
$$(x_{n1}e^0)(z)=\frac{1}{n!}x(z)^{(n)},x𝔥_L,$$
where
$$x(z)=\underset{j}{}x_jz^{j1}.$$
In particular, $`(x_1e^0)(z)=x(z)`$.
We continue in the same vein:
$$(\mathrm{\Phi }_n^ie^0)(z)=\frac{1}{n!}\mathrm{\Phi }^i(z)^{(n)},$$
where
$$\mathrm{\Phi }^i(z)=\underset{j}{}\mathrm{\Phi }_j^iz^j;$$
$$(\mathrm{\Psi }_{n1}^ie^0)(z)=\frac{1}{n!}\mathrm{\Psi }^i(z)^{(n)},$$
where
$$\mathrm{\Psi }^i(z)=\underset{j}{}\mathrm{\Psi }_j^iz^{j1};$$
$$e^\alpha (z)=e^\alpha \mathrm{exp}(\underset{n<0}{}\frac{\alpha _n}{n}z^n)\mathrm{exp}(\underset{n>0}{}\frac{\alpha _n}{n}z^n)z^{\alpha _0}.$$
Finally,
$$x_{n_1}^{(1)}x_{n_2}^{(2)}\mathrm{}x_{n_k}^{(k)}e^\alpha (z)=:x_{n_1}^{(1)}(z)x_{n_2}^{(2)}(z)\mathrm{}x_{n_k}^{(k)}(z)e^\alpha (z):.$$
The vertex algebra structure on $`V_L`$ is equivalently described by the following family of $`n`$-th products ($`n`$):
$${}_{(n)}{}^{}:V_LV_LV_L,xyx_{(n)}y\stackrel{\text{def}}{=}(x(z)z^n)(y),$$
where $`x(z)z^n`$ stands for the linear transformation of $`V_L`$ equal to the coefficient of $`z^{n1}`$ in the series $`x(z)`$.
1.3. Degeneration of $`V_L`$. Denote by $`L_A`$ the subgroup of $`L`$ generated by $`A^i,i=1,\mathrm{},N`$. Any smooth toric variety $`X`$ can be defined via a fan, $`\mathrm{\Sigma }`$, that is, a collection of “cones” lying in $`L_A`$. Borisov uses such $`\mathrm{\Sigma }`$ to define a certain degeneration, $`V_L^\mathrm{\Sigma }`$, of the vertex algebra structure on $`V_L`$. He further shows that the cohomology of $`V_L^\mathrm{\Sigma }`$ with respect to a certain differential $`D^\mathrm{\Sigma }:V_L^\mathrm{\Sigma }V_L^\mathrm{\Sigma }`$ equals $`H^{}(X,\mathrm{\Omega }_X^{ch})`$, where $`\mathrm{\Omega }_X^{ch}`$ is the chiral de Rham complex of \[MSV\]. Let us describe the outcome of this construction in the case when $`X=^N`$.
Consider the following set of $`N+1`$ elements of $`L_A`$: $`\xi _1=A^1,\xi _2=A^2,\mathrm{},\xi _N=A^N,\xi _{N+1}=A^1A^2\mathrm{}A^N`$. Define the cone $`\mathrm{\Delta }_iL_A`$ to be the set of all non-negative integral linear combinations of the elements $`\xi _1,\mathrm{},\xi _{i1},\xi _{i+1},\mathrm{},\xi _{N+1}`$. It is easy to see that $`L_A=_i\mathrm{\Delta }_i`$ and the intersection $`\mathrm{\Delta }_i\mathrm{\Delta }_j`$ is a face of both $`\mathrm{\Delta }_i`$ and $`\mathrm{\Delta }_j`$. The fan $`\mathrm{\Sigma }`$ in this case is the set consisting of $`\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_{N+1}`$ and their faces.
We now define $`V_L^\mathrm{\Sigma }`$ to be a vertex algebra equal to $`V_L`$ as a vector space with $`n`$-th product <sub>(n),Σ</sub> as follows:
if $`\{_in_iA^i,_in_i^{}A^i\}\mathrm{\Delta }_j`$ for some $`j`$, then
$$(xe^{_im_iB^i+_in_iA^i})_{(n),\mathrm{\Sigma }}(ye^{_im_i^{}B^i+_in_i^{}A^i})$$
$$=(xe^{_im_iB^i+_in_iA^i})_{(n)}(ye^{_im_i^{}B^i+_in_i^{}A^i});$$
otherwise
$$(xe^{_im_iB^i+_in_iA^i})_{(n),\mathrm{\Sigma }}(ye^{_im_i^{}B^i+_in_i^{}A^i})=0,$$
where <sub>(n)</sub> stands for the n-th product on $`V_L`$. The fact that these new operations satisfy the Borcherds identities can be proved by including both $`V_L`$ and $`V_L^\mathrm{\Sigma }`$ in a 1-parameter family of vertex algebras; this will be done in 2.1.
Let
$$D=\underset{i=1}{\overset{N}{}}\mathrm{\Psi }^i(z)(e^{A^i}e^{_jA^j})(z).$$
$`(1.1)`$
It is obvious that $`D\text{End}(V_L^\mathrm{\Sigma })`$ and $`D^2=0`$; therefore, the cohomology $`H_D(V_L^\mathrm{\Sigma })`$ arises.
Theorem 1.3. (\[B\])
$$H_D(V_L^\mathrm{\Sigma })=H^{}(^N,\mathrm{\Omega }_^N^{ch}).$$
§2. Deforming $`H^{}(^N)`$
2.1. The family $`V_{L,q}`$.
Here we exhibit a family of vertex algebras, $`V_{L,q},q`$, so that $`V_{L,q}`$ is isomorphic to $`V_L`$ if $`q0`$ and $`V_{L,0}`$ is isomorphic to $`V_L^\mathrm{\Sigma }`$; cf. the end of sect.8 \[B\].
Define the height function
$$ht:L_A_>$$
as follows. It is easy to see that each $`\alpha L_A`$ is uniquely represented in the form
$$\alpha =\underset{i=1}{\overset{N+1}{}}n_i\xi _i$$
$`(2.1)`$
so that all $`n_i0`$ and $`\mathrm{\#}\{i:n_i>0\}N`$. Let
$$ht(\alpha )=\underset{i}{}n_i,$$
where $`n_1,\mathrm{},n_N`$ are as in (2.1).
Define the linear automorphism
$$t_q:V_LV_L,q\{0\}$$
by the formula
$$t_q(xe^{_im_iB^i+_in_iA^i})=q^{ht(_in_iA^i)}xe^{_im_iB^i+_in_iA^i}.$$
Define $`V_{L,q}`$ to be the vertex algebra equal to $`V_L`$ as a vector space with the following n-th product:
$$(xe^{_im_iB^i+_in_iA^i})_{(n),q}(ye^{_im_i^{}B^i+_in_i^{}A^i})$$
$$=t_q^1(t_q(xe^{_im_iB^i+_in_iA^i})_{(n)}t_q(ye^{_im_i^{}B^i+_in_i^{}A^i})).$$
By definition,
$$t_q:V_{L,q}V_L,q\{0\},$$
is a vertex algebra isomorphism. It is also easy to see that if $`_in_iA^i`$ and $`_in_i^{}A^i`$ belong to the same cone from $`\mathrm{\Sigma }`$, then
$$(xe^{_im_iB^i+_in_iA^i})_{(n),q}(ye^{_im_i^{}B^i+_in_i^{}A^i})$$
$$=(xe^{_im_iB^i+_in_iA^i})_{(n)}(ye^{_im_i^{}B^i+_in_i^{}A^i});$$
otherwise
$$(xe^{_im_iB^i+_in_iA^i})_{(n),q}(ye^{_im_i^{}B^i+_in_i^{}A^i})$$
$$q[q](xe^{_im_iB^i+_in_iA^i})_{(n)}(ye^{_im_i^{}B^i+_in_i^{}A^i}).$$
Two things follow at once: first, the operations
$${}_{(n),0}{}^{}=\underset{q0}{lim}{}_{(n),q}{}^{},n$$
are well defined and satisfy the Borcherds identities; second, the vertex algebra, $`V_{L,0}`$, obtained in this way is isomorphic to $`V_L^\mathrm{\Sigma }`$. By the way, this remark proves that $`V_L^\mathrm{\Sigma }`$ is indeed a vertex algebra.
To get a better feel for this kind of deformation, and for the future use, let us consider the subspace $`[L_A]V_{L,q}`$ with basis $`e^\alpha ,\alpha L_A`$. The $`(1)`$-st product makes this space a commutative algebra. The subspace $`[\mathrm{\Delta }_j]`$ defined to be the linear span of $`e^\alpha ,\alpha \mathrm{\Delta }_j,`$ is a polynomial ring on generators $`e^{\xi _1},\mathrm{},e^{\xi _{j1}},e^{\xi _{j+1}},\mathrm{},e^{\xi _{N+1}}`$. For example, if we denote $`x_i=e^{A^i}`$, then $`[\mathrm{\Delta }_{N+1}]=[x_1,\mathrm{},x_N]`$ and this isomorphism identifies $`e^{_jn_jA^j}`$ with the monomial $`x_1^{n_1}\mathrm{}x_N^{n_N}`$.
The entire $`[L_A]`$ is not a polynomial ring. For example, as follows from the definition of the deformation, there is a relation
$$(e^{A^1\mathrm{}A^N})_{(1)}(e^{A^1+\mathrm{}+A^N})=q^{N+1}e^0,$$
because $`ht(0)=0`$, $`ht(A^1+\mathrm{}+A^N)=N`$, $`ht(A^1\mathrm{}A^N)=1`$. If we let $`T=e^{A^1\mathrm{}A^N}`$, then the last equality rewrites as follows:
$$Tx_1x_2\mathrm{}x_N=q^{N+1},$$
and a moment’s thought shows that in fact
$$[L_A]=[x_1,\mathrm{},x_N,T]/(Tx_1x_2\mathrm{}x_Nq^{N+1}).$$
Being a group algebra, $`[L_A]`$ carries another algebra structure, a priori different from the one we just described and independent of $`q`$. We see that the two structures are isomorphic if $`q0`$; at $`q=0`$, however, the one we just described degenerates in an algebra with zero divizors.
2.2. The algebra $`H^{}(^N)`$.
Let
$$Q(z)=A^i(z)\mathrm{\Phi }^i(z)\underset{j}{}\mathrm{\Phi }^j(z)^{},$$
$$G(z)=B^i(z)\mathrm{\Psi }^i(z),$$
$$J(z)=:\mathrm{\Phi }^i(z)\mathrm{\Psi }^i(z):+\underset{j}{}B^j(z)^{},$$
$$L(z)=:B^i(z)A^i(z):+:\mathrm{\Phi }^i(z)^{}\mathrm{\Psi }^i(z):,$$
where the summation with respect to repeated indices is assumed.
One checks that the Fourier components of these 4 fields satsify the commutation relations of the $`N=2`$ algebra. It is also easy to see that the fields $`G(z),L(z)`$ commute with Borisov’s differential $`D`$, see (1.1), and therefore define the fields, to be also denoted $`G(z),L(z)`$, acting on $`H_D(V_L^\mathrm{\Sigma })`$.
The fields $`Q(z),J(z)`$ do not commute with $`D`$, but their Fourier components $`Q_0=Q(z)`$ and $`J_0=J(z)`$ do:
$$[Q_0,D]=[J_0,D]=0.$$
Thus we get 2 operators, to be also denoted $`Q_0,J_0`$, acting on $`H_D(V_L^\mathrm{\Sigma })`$. All this is summarized by saying that $`H_D(V_L^\mathrm{\Sigma })`$ is a topological vertex algebra.
A glance at the formulas on p. 17 of \[B\] shows that the isomorphism $`H_D(V_L^\mathrm{\Sigma })=H^{}(^N,\mathrm{\Omega }_^N^{ch})`$ (see Theorem 1.3) identifies these $`G(z),L(z),Q_0,J_0`$ with the fields (operators) constructed in \[MSV\] and denoted in the same way. One of the main results of \[MSV\] then gives
$$H^{}(^N)=H_{Q_0}(H_D(V_L^\mathrm{\Sigma })).$$
$`(2.2)`$
Further, the algebra structure of $`H^{}(^N)`$ is restored from the $`(1)`$-st product on $`H_D(V_L^\mathrm{\Sigma })`$.
2.3. Deformation of the algebra structure.
It follows from the proof of Theorem 2.3 below that the cohomology (2.2) can be calculated in the reversed order:
$$H^{}(^N)=H_D(H_{Q_0}(V_L^\mathrm{\Sigma })).$$
$`(2.3)`$
Note that $`D`$ and $`Q_0`$ can also be regarded as well-defined operators acting on the deformed algebra:
$$D=\underset{i=1}{\overset{N}{}}\mathrm{\Psi }^i(z)(e^{A^i}e^{_jA^j})(z),Q_0=A^i(z)\mathrm{\Phi }^i(z):V_{L,q}V_{L,q}.$$
It is immediate to see that $`D^2=Q_0^2=0`$ on $`V_{L,q}`$ as well. Moreover,
$$[D,Q_0]=0.$$
$`(2.4)`$
Indeed, the formulas of 1.2 imply the following OPE:
$$\underset{i=1}{\overset{N}{}}\mathrm{\Psi }^i(z)(e^{A^i}e^{_jA^j})(z)A^j(w)\mathrm{\Phi }^j(w)=\frac{_je^{A^j}(w)^{}e^{_jA^j}(w)^{}}{zw}.$$
Therefore,
$$[D,Q_0]=\{\underset{j}{}e^{A^j}(w)e^{_jA^j}(w)\}^{}=0.$$
Thus it is natural to take the space $`H_D(H_{Q_0}(V_{L,q}))`$ for a deformation of $`H^{}(^N)`$.
Theorem 2.3.
$$H_D(H_{Q_0}(V_{L,q}))=[T]/(T^{N+1}q^{N+1}).$$
Proof.
1) Computation of $`H_{Q_0}(V_{L,q})`$. By definition
$$Q_0=\underset{n}{}A_n^i\mathrm{\Phi }_n^i$$
$`(2.5)`$
Therefore,
$$[Q_0,\mathrm{\Psi }_0^j]=A_0^j,[Q_0,G_0]=L_0.$$
These relations imply that
$$H_{Q_0}(V_{L,q})=H_{Q_0}(_jKerA_0^jKerL_0).$$
$`(2.6)`$
It follows from 1.2 that the space $`_jKerA_0^jKerL_0`$ is a linear span of elements of the form:
$$\mathrm{\Phi }_0^{i_1}\mathrm{}\mathrm{\Phi }_0^{i_m}e^{_in_iA^i}.$$
Formula (2.5) shows that the restriction of $`Q_0`$ to this subspace is 0. Thus
$$H_{Q_0}(_jKerA_0^jKerL_0)=_jKerA_0^jKerL_0.$$
The $`(1)`$-st product makes this subspace a supercommutative algebra. In the same way as in 2.1 we get an isomorphism
$$_jKerA_0^jKerL_0=[x_1,\mathrm{},x_N,T;\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N]/(Tx_1\mathrm{}x_Nq^{N+1}),$$
where $`\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N`$ are understood as grassman variables ($`[x_i,\mathrm{\Phi }_j]=[T,\mathrm{\Phi }_j]=0`$, $`\mathrm{\Phi }_i\mathrm{\Phi }_j+\mathrm{\Phi }_j\mathrm{\Phi }_i=0`$) and $`(Tx_1\mathrm{}x_Nq^{N+1})`$ stands for the ideal generated by $`Tx_1\mathrm{}x_Nq^{N+1}`$.
2) Computation of $`H_D(H_{Q_0}(V_{L,q}))`$. In view of Step 1), we have to restrict $`D`$ to
$$_jKerA_0^jKerL_0.$$
The isomorphism
$$_jKerA_0^jKerL_0=[x_1,\mathrm{},x_N,T;\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N]/(Tx_1\mathrm{}x_Nq^{N+1}),$$
identifies $`D`$ with $`_i(x_iT)/(\mathrm{\Phi }_i)`$. Therefore, the complex
$$([x_1,\mathrm{},x_N,T;\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N]/(Tx_1\mathrm{}x_Nq^{N+1}),D)$$
is simply the Koszul resolution of the algebra
$$\{[x_1,\mathrm{},x_N,T]/(Tx_1\mathrm{}x_Nq^{N+1})\}/(x_1T,x_2T,\mathrm{},x_NT)$$
associated with the sequence $`x_1T,x_2T,\mathrm{},x_NT`$. This sequence is regular and we get at once
$$H_D(H_{Q_0}(V_{L,q}))=[T]/(T^{N+1}q^{N+1}).\mathit{}$$
It is easy to infer from Borisov’s proof of Theorem 1.3 that the element $`T=e^{A^1\mathrm{}A^N}V_{L,q}`$ is a cocycle representing the cohomology class proportional to that of a hyperplane in $`^N`$. This means that the deformation of $`H^{}(^N)`$ we obtained coincides with the standard one, except that for some reason $`q`$ happened to be raised to the power of $`N`$.
2.4. Reduction to a single differential.
Of course it would be nicer to get $`H^{}(^N)`$, or its deformation, as the cohomology of this or that vertex algebra with respect to a single differential rather than to compute a repeated cohomology.
Theorem 2.4.
$$H_{D+Q_0}(V_{L,q}))=[T]/(T^{N+1}q^{N+1}).$$
It is no wonder, in view of Theorem 2.3, that this assertion is a result of computation of a certain spectral sequence. We shall use several spectral sequences arising in the following situation, which is slightly different from the standard one. Let
$$W=_{n=\mathrm{}}^+\mathrm{}W^n$$
be a graded vector space with two commuting differentials
$$d_1:W^nW^{n+1},d_2:W^nW^{n1}.$$
$`(2.7)`$
There arise the total differential $`d=d_1+d_2`$ and the cohomology $`H_{d_1+d_2}(W)`$. Note that this cohomology group is not graded since $`d_1`$ and $`d_2`$ map in opposite directions. We can, however, introduce the filtration
$$W=_nW^n,W^n=_{m=\mathrm{}}^nW^m.$$
Then
$$(d_1+d_2)(W^n)W^{(n+1)}$$
and there arises a filtration $`H_{d_1+d_2}(W)^n`$ on the cohomology and the graded object $`GrH_{d_1+d_2}(W)`$.
It is straightforward to define a spectral sequence
$$\{E(W)_r^n,d^{(r)}:E(W)_r^nE(W)_r^{nr+1}\},E(W)_{r+1}^n=H_{d^{(r)}}(E(W)_r^n),$$
$`(2.8)`$
the first three terms being as follows:
$$E(W)_0^n=W^n,E(W)_1^n=H_{d_1}(W^n),E(W)_2^n=H_{d_2}(H_{d_1}(W^n)),$$
$`(2.9)`$
where
$$H_{d_1}(W^n)=\frac{\text{Ker}\{d_1:W^nW^{n+1}\}}{\text{Im}\{d_1:W^{n1}W^n\}},$$
$$H_{d_2}(H_{d_1}(W^n))=\frac{\text{Ker}\{d_2:H_{d_1}(W^n)H_{d_1}(W^{n1})\}}{\text{Im}\{d_2:H_{d_1}(W^{n+1})H_{d_1}(W^n)\}}.$$
In the situation pertaining Theorem 2.4 we take $`V_{L,q}`$ for $`W`$, $`Q_0`$ for $`d_1`$, and $`D`$ for $`d_2`$. The space $`V_{L,q}`$ is graded by fermionic charge; this grading is defined by letting the degree of $`\mathrm{\Psi }_j^i`$ be equal $`1`$, the degree of $`\mathrm{\Phi }_j^i`$ be equal $`1`$, the degree of $`A_j^i,B_j^i,e^\alpha `$ be equal $`0`$. By definition,
$$Q_0(V_{L,q}^n)V_{L,q}^{n+1},D(V_{L,q}^n)V_{L,q}^{n1}$$
and we get a spectral sequence $`\{E(V_{L,q})_r^n,d^{(r)}\}`$.
Observe that the grading by fermionic charge and the corresponding filtration are infinite in both directions. Therefore, the standard finiteness conditions that guarantee convergence of spectral sequences fail. Nevertheless the following lemma holds true.
Lemma 2.4. The spectral sequence $`\{E(V_{L,q})_r^n,d^{(r)}\}`$ converges to $`H_{Q_0+D}(V_{L,q})`$ and collapses:
$$H_D(H_{Q_0}(V_{L,q}))=H_{Q+D}(V_{L,q}).$$
Lemma 2.4 combined with Theorem 2.3 gives Theorem 2.4 at once and it remains to prove Lemma 2.4.
Proof of Lemma 2.4. Introduce yet another grading of the space $`V_{L,q}^n`$ as follows. Let $`\alpha =(\alpha _1,\mathrm{},\alpha _{N+1})`$ be an element of the group $`^{N+1}`$. Let
$$V_{L,q}^n[\alpha ]=(_{i=1}^NKer(A_0^i\alpha _iId))Ker(L_0\alpha _{N+1}Id).$$
Of course
$$V_{L,q}^n=_{\alpha ^{N+1}}V_{L,q}^n[\alpha ]$$
and both the differentials preserve this grading. Therefore all calculations can be carried out inside $`V_{L,q}^n[\alpha ]`$ with a fixed $`\alpha `$. Consider the following two cases.
1) $`\alpha 0`$. In this case, as was observed in the beginnning of the proof of Theorem 2.3 (see e.g. (2.6)), $`H_{Q_0}(V_{L,q}[\alpha ])=0`$ and, therefore, $`E[\alpha ]_1=0`$. It remains to show that $`H_{Q_0+D}(V_{L,q}[\alpha ])=0`$. Let $`xV_{L,q}[\alpha ]^n`$ be a cocycle. This means that there is a “chain” of elements $`x_iV_{L,q}[\alpha ]^{n2i},i=0,1,\mathrm{},k`$ so that
$$x=\underset{i=0}{\overset{k}{}}x_i,$$
and the following holds
$$Q_0(x_0)=0,Q_0(x_{i+1})+D(x_i)=0,D(x_k)=0,i=0,\mathrm{},k1.$$
$`(2.10)`$
We now repeatedly use the condition $`H_{Q_0}(V_{L,q}[\alpha ])=0`$ and (2.10) to construct another chain $`y_iV_{L,q}[\alpha ]^{n2i1},i0`$, satisfying
$$Q_0(y_0)=x_0,Q_0(y_{i+1})+D(y_i)=x_{i+1}.$$
$`(2.11)`$
Indeed, since $`Q_0(x_0)=0`$, there is $`y_0V_{L,q}[\alpha ]^{n1}`$ so that $`Q_0(y_0)=x_0`$.
Since
$$Q_0(D(y_0)+x_1)=DQ_0(y_0)+Q_0(x_1)=D(x_0)+Q_0(x_1)=0,$$
there is $`y_1V_{L,q}[\alpha ]^{n3}`$ so that $`Q_0(y_1)+D(y_0)=x_1`$.
In general, having found $`y_iV_{L,q}[\alpha ]^{n2i1},y_{i1}V_{L,q}[\alpha ]^{n2i+1}`$ so that $`Q_0(y_i)+D(y_{i1})=x_i`$, we calculate as follows:
$$0=D(0)=D(Q_0(y_i)+D(y_{i1})x_i)=DQ_0(y_i)D(x_i).$$
Due to (2.10), the last expression rewrites as $`DQ_0(y_i)+Q_0(x_{i+1})`$ and we get
$$Q_0D(y_i)+Q_0(x_{i+1})=0.$$
Therefore, $`Q_0(D(y_i)x_{i+1})=0`$ and there is $`y_{i+1}`$ so that $`Q_0(y_{i+1})=D(y_i)+x_{i+1}`$ as desired.
Formally, (2.11) means that
$$(D+Q_0)(\underset{i=0}{\overset{\mathrm{}}{}}y_i)=x$$
and what does not allow us to conclude immediately that $`x=_{i=0}^{\mathrm{}}x_i`$ is a coboundary is that the sum $`_{i=0}^{\mathrm{}}y_i`$ looks infinite. To complete case 1) it remains to show that $`y_i=0`$ for all sufficiently large $`i`$. This is achieved by the following dimensional argument. Note that by construction
$$y_i_{|m_j|<ki}(S_{𝔥_L}\mathrm{\Lambda }_{𝔥_L}^{n2i1}e^{_jm_jA^j+_j\alpha _jB^j}),$$
$`(2.12)`$
where $`k`$ is a number independent of $`i`$. Indeed, each application of $`D`$ changes $`m_j`$ by at most 1, $`Q_0`$ preserves $`m_j`$, and the linear estimate of $`m_j`$ follows. On the other hand we have an explicit formula for $`L_0`$ (see the beginning of 2.2), and this formula implies that the smallest eigenvalue of $`L_0`$ restricted to $`\mathrm{\Lambda }_{𝔥_L}^{n2i1}`$ is nonnegative and grows faster than a polynomial of degree 2, say $`q(i)`$, as $`i+\mathrm{}`$. The same formula gives
$$L_0e^{_jm_jA^j+_j\alpha _jB^j}=\underset{j}{}m_j\alpha _je^{_jm_jA^j+_j\alpha _jB^j}.$$
Therefore, if $`y_i0`$, then it is a sum of eigenvectors associated to eigenvalues of $`L_0`$ greater or equal $`q(i)(\alpha _1+\mathrm{}\alpha _n)ki`$. Since this number tends to $`+\mathrm{}`$ as $`i+\mathrm{}`$, we arrive at contradiction with the assumption $`L_0y_i=\alpha _{N+1}`$ if $`i`$ is sufficiently large. Hence, $`y_i=0`$ for all sufficiently large $`i`$, each cocycle is a coboundary, and case 1) is accomplished.
2) $`\alpha =0`$. As we saw in the beginning of the proof of Theorem 2.3, the restriction of $`Q_0`$ to $`V_{L,q}[0]`$ is 0 and, by definition, the complex $`(V_{L,q}[0],D+Q_0)`$ is equal to $`(E(V_{L,q})[0]_1,d^{(1)})`$.
2.5. The vertex algebra $`H_D(V_{L,q})`$ and a computation of $`H^{}(^N,\mathrm{\Omega }_^N^{ch})`$.
In this section we prove the following two theorems.
Theorem 2.5A If $`q0`$, then $`H_D(V_{L,q})`$ equals the quantum cohomology of $`^N`$.
Theorem 2.5B The natural embedding of sheaves (\[MSV\], see also (2.20) below)
$$\mathrm{\Omega }_^N^{}\mathrm{\Omega }_^N^{ch}$$
provides an isomorphism
$$H^i(^N,\mathrm{\Omega }_^N^{})\stackrel{}{}H^i(^N,\mathrm{\Omega }_^N^{ch}),\mathrm{\hspace{0.33em}0}<i<N,$$
where $`\mathrm{\Omega }_^N^{}`$ is the sheaf of all differential forms.
Recall the previously known results on the cohomology of $`\mathrm{\Omega }_^N^{ch}`$. $`\mathrm{\Omega }_^N^{ch}`$ is a sheaf of $`\widehat{sl}_{N+1}`$-modules \[MS1\], see also 2.5.2. In particular, if $`U_0=^N^N`$ is a big cell, then $`\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})`$ is a generalized Wakimoto module over $`\widehat{sl}_{N+1}`$ introduced in \[FF\]. We proved in \[MS1\] that
$$H^0(^N,\mathrm{\Omega }_^N^{ch})=\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})^{int},$$
$`(2.13)`$
where $`\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})^{int}`$ stands for the maximal $`sl_{N+1}`$-integrable submodule of $`\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})`$.
On the other hand, it follows from the chiral Serre duality \[MS2\] that
$$H^N(^N,\mathrm{\Omega }_^N^{ch})=H^0(^N,\mathrm{\Omega }_^N^{ch})^d,$$
$`(2.14)`$
where <sup>d</sup> stands for the restricted dual.
Unfortunately, little is known about the structure of $`\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})`$ and $`\mathrm{\Gamma }(U_0,\mathrm{\Omega }_^N^{ch})^{int}`$, if $`N>1`$; see, however, \[MS1\] for the case of $`N=1`$. Otherwise, Theorem 2.5 and (2.13-14) give a complete description of $`H^{}(^N,\mathrm{\Omega }_^N^{ch})`$.
The proofs of Theorems 2.5A and B are contained in 2.5.2. In 2.5.1 we collect some well-known material in order to place these results in the proper context and to formulate (2.18-19), two well-known assertions needed in 2.5.2.
2.5.1 A vertex algebra structure on a vector space $`V`$ comprises a countable family of multiplications:
$${}_{(n)}{}^{}:VVV,xyx_{(n)}y,n,$$
a map
$$T:VV,$$
and a vacuum vector
$$\text{1}V.$$
These data satisfy the Borcherds identities which imply, in particular, that $`T`$ and $`x_{(0)},xV`$, are derivations of the $`n`$-th product for all $`n`$. Thus,
$$[T,y_{(j)}]=(Ty)_{(j)},[x_{(0)},y_{(j)}]=(x_{(0)}y)_{(j)}.$$
$`(2.15)`$
In the case of the vertex algebra $`V_{L,q}`$, the $`n`$-th multiplication was defined in the end of 1.2, 1 equals $`e^0`$, and $`T`$ will be defined below.
Call $`V`$ commutative (or holomorphic, see \[K\] 1.4) if $`{}_{(n)}{}^{}=0`$ for all $`n0`$. If $`V`$ is commutative, then the (-1)-st multiplication gives it the structure of a commutative superalgebra with derivation $`T`$, and the functor arising in this way is an equivalence of the category of commutative vertex algebras and the category of commutative superalgebras with derivation, see again \[K\] 1.4.
If $`d_x:VV`$ is a differential ($`d^2=0`$), then the cohomology $`H_{d_x}(V)`$ arises. We assert that
$$d_x=x_{(0)}\text{ for some }xVH_{d_x}(V)\text{ is a vertex algebra},$$
$`(2.16)`$
since all products on $`V`$ descend to $`H_{d_X}(V)`$ due to (2.15).
All vertex algebras we are concerned with are conformal. This means that there is a Virasoro field $`L(z)=_iL_iz^{i2}`$, $`L_iEnd(V)`$, such that $`L_i`$ satisfy the Virasoro commutation relations, $`T=L_1`$, $`L_0`$ is diagonalizable, and $`L(z)`$ is the field attached to the state $`L_2\text{1}V`$. The formula at the beginning of 2.2 shows that $`V_{L,q}`$ is a conformal vertex algebra, the state $`L_2\text{1}`$ being equal to $`_i(B_1^iA_1^i+\mathrm{\Phi }_1^i\mathrm{\Psi }_1^i)e^0`$.
The eigenvalues of $`L_0`$ are called conformal weights. Hence a conformal vertex algebra $`V`$ is graded by conformal weights, $`V=_nV_n`$, and in the case of $`V=V_{L,q}`$ this grading (but not the name) has already been used in the proofs of Theorems 2.3 and 2.4.
Returning to the cohomology vertex algebra $`H_{d_x}(V)`$ in the case when $`V`$ is conformal and $`x`$ is an eigenvector of $`L_0`$, we see that
$$L_2\text{1}Kerd_xH_{d_x}(V)\text{ is conformal,}$$
$`(2.17)`$
$$L_2\text{1}Imd_xH_{d_x}(V)\text{ is commutative.}$$
$`(2.18)`$
Indeed, if $`L_2\text{1}Kerd_X`$, then the operators $`L_iEnd(V)`$ descend to $`H_{d_X}(V)`$ due to (2.15). If, in addition, $`L_2\text{1}=d_x(y)`$, then all $`L_i`$’s act on $`H_{d_x}(V)`$ trivially again due to (2.15). Hence $`L_0`$ acts on $`H_{d_X}(V)`$ trivially, each element of $`H_{d_X}(V)`$ is represented by a cocycle of conformal weight 0, and the $`n`$-th product on $`H_{d_X}(V)`$ vanishes unless $`n=1`$.
If $`xV_1`$, then $`d_x(V_n)V_n`$ for all $`n`$, and (2.18) can be sharpened as follows:
$$L_2\text{1}Imd_x\text{ and }xV_1H_{d_X}(V)=H_{d_X}(V_0).$$
$`(2.19)`$
In our previous work (\[MSV\], \[MS1\], \[MS2\]) we have dealt with conformal vertex algebras having the following properties: all conformal weights are nonnegative; the conformal weight 0 component is a finitely generated supercommutative ring and the corresponding multiplication coincides with the restriction of the (-1)st multiplication. For example, $`\mathrm{\Omega }_X^{ch}`$ is a sheaf of such vertex algebras over a smooth manifold $`X`$: the conformal weight 0 component of $`\mathrm{\Gamma }(U,\mathrm{\Omega }_X^{ch})`$ is the algebra of differential forms over $`UX`$. In other words, there is a natural embedding
$$\mathrm{\Omega }_X^{}\stackrel{}{}\mathrm{\Omega }_{X,0}^{ch}\mathrm{\Omega }_X^{ch},$$
$`(2.20)`$
and it is this embedding that was invoked in Theorem 2.5B.
$`H^{}(X,\mathrm{\Omega }_^N^{ch})`$ is also a vertex algebra of this kind because its conformal weight 0 component equals the cohomology algebra $`H^{}(X)`$. It is, therefore, natural to ask if there is a conformal vertex algebra with nonnegative conformal weights so that the (-1)-st multiplication identifies its conformal weight 0 component with the quantum cohomology of $`X`$.
The quantum cohomology itself is one such vertex algebra due to the equivalence of categories reviewed above. A more appealing possibility seems to be provided by $`H_D(V_{L,q})`$: it is a vertex algebra due to (2.16) because (1.1) is equivalent to
$$D=\underset{i=1}{\overset{N}{}}(\mathrm{\Psi }_1^i(e^{A^i}e^{_jA^j}))_{(0)},$$
$`(2.21)`$
and it is conformal because, as one easily checks, $`D(L_2e^0)=0`$.
Even though Theorem 2.5A says that in this way we do not get anything new either, it allows us to observe a curious phenomenon: $`H_D(V_{L,q})`$, $`q`$, is a family of vertex algebras over $``$ with fiber that equals $`H^{}(^N)`$ over any non-zero point and blows up to the non-commutative infinite dimensional vertex algebra $`H^{}(^N,\mathrm{\Omega }_^N^{ch})`$ over $`0`$.
Rather unexpectedly, Theorem 2.5B turns out to be a by-product of the proof of Theorem 2.5A.
2.5.2 Proof Theorems 2.5A and B.
By definition, the complex $`(V_{L,q},D)`$ is the constant vector space $`V_L`$ with differential $`D`$ polynomially depending on $`q`$. To make this more precise, observe that $`V_L`$ is graded by the function $`ht`$ defined in 2.1:
$$V_L=_{n0}V_L^n,$$
$`(2.22)`$
where $`V_L^n`$ is a linear span of $`xe^{_im_iB^i+_in_iA^i}`$ with $`ht(_in_iA^i)=n`$. The differential $`D`$ then breaks in a sum
$$D=d_++q^Nd_{},$$
$`(2.23a)`$
so that
$$d_+(V_L^n)V_L^{n+1},$$
$`(2.23b)`$
$$d_{}(V_L^n)V_L^{nN},$$
$`(2.23c)`$
and
$$(d_+)^2=(d_{})^2=[d_+,d_{}]=0.$$
$`(2.23d)`$
Again by definition, the complex $`(V_L,d_+)`$ coincides with Borisov’s complex $`(V_L^\mathrm{\Sigma },D)`$. It follows from formulas (2.23a-d) and Theorem 1.3 that there is a spectral sequence of the same type as (2.8), the 1st term and the 1st differential being as follows
$$E_1=H^{}(^N,\mathrm{\Omega }_^N^{ch})$$
$`(2.24)`$
$$d_1=q^Nd_{}:H^{}(^N,\mathrm{\Omega }_^N^{ch})H^{}(^N,\mathrm{\Omega }_^N^{ch}),$$
$$d_{}(H^n(^N,\mathrm{\Omega }_^N^{ch}))H^{nN}(^N,\mathrm{\Omega }_^N^{ch}).$$
$`(2.25)`$
Simply because $`\text{dim}^N=N`$, the 2nd term equals
$$\frac{H^0(^N,\mathrm{\Omega }_^N^{ch})}{\text{Im}\{d_{}:H^N(^N,\mathrm{\Omega }_^N^{ch})H^0(^N,\mathrm{\Omega }_^N^{ch})\}}$$
$$\text{Ker}\{d_{}:H^N(^N,\mathrm{\Omega }_^N^{ch})H^0(^N,\mathrm{\Omega }_^N^{ch})\}_{i=1}^{N1}H^i(^N,\mathrm{\Omega }_^N^{ch}),$$
and all higher differentials vanish. An argument similar to (and simpler than) the one used in the proof of Lemma 2.4 shows that this spectral sequence converges to $`H_D(V_{L,q})`$. Therefore
$$H_D(V_{L,q})$$
$$=\frac{H^0(^N,\mathrm{\Omega }_^N^{ch})}{\text{Im}\{d_{}:H^N(^N,\mathrm{\Omega }_^N^{ch})H^0(^N,\mathrm{\Omega }_^N^{ch})\}}$$
$$\text{Ker}\{d_{}:H^N(^N,\mathrm{\Omega }_^N^{ch})H^0(^N,\mathrm{\Omega }_^N^{ch})\}_{i=1}^{N1}H^i(^N,\mathrm{\Omega }_^N^{ch}).$$
$`(2.26)`$
Lemma 2.6. There is $`yH^N(^N,\mathrm{\Omega }_^N^{ch})`$ such that $`d_{}(y)H^0(^N,\mathrm{\Omega }_^N^{ch})`$ equals the Virasoro element $`L_2e^0`$.
This lemma allows us to complete the proof of Theorems 2.5A and B instantaneously. Our differentials come from elements of $`V_L`$ of conformal weight 1, see (2.21); hence, due to Lemma 2.6, (2.18) and (2.19) apply: $`H_D(V_{L,q})`$ equals $`H_D((V_{L,q})_0)`$, which is known (Theorem 2.4) to be equal to the quantum cohomology. In particular, as follows from (2.26),
$$H^i(^N,\mathrm{\Omega }_^N^{ch})=H^i(^N,(\mathrm{\Omega }_^N^{ch})_0),\mathrm{\hspace{0.33em}0}<i<N,$$
the latter space being canoncally isomorphic to $`H^i(^N,\mathrm{\Omega }_^N^{})`$ due to (2.20). Thus it remains to prove Lemma 2.6.
Proof of Lemma 2.6 To find an appropriate $`yH^N(^N,\mathrm{\Omega }_^N^{ch})`$ and calculate $`d_{}(y)`$ we need to take a plunge in \[MSV,B\].
Let $`x^0:x^1:\mathrm{}:x^N`$ be homogeneous coordinates on $`^N`$ and $`b^i=x^i/x^0`$. Consider the $`N`$-dimensional torus $`𝕋^N=\text{Spec}[(b^1)^{\pm 1},\mathrm{},(b^N)^{\pm 1}]^N`$.
We shall need the following facts about the sheaf $`\mathrm{\Omega }_^N^{ch}`$.
First,
$$\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})=[(b_0^i)^{\pm 1},b_{j1}^i,a_{j1}^i;\varphi _j^i,\psi _{j1}^i;1iN,j0],$$
$`(2.27)`$
where $`b_j^i,a_{j1}^i`$ are even, $`\varphi _j^i,\psi _{j1}^i`$ odd.
By letting $`\text{deg}x_j^i=j`$, $`x=b,a,\varphi `$ or $`\psi `$, we recover the grading by conformal weight. By letting $`\text{deg}b_j^i=\text{deg}a_j^i=0`$, $`\text{deg}\varphi _j^i=1`$, $`\text{deg}\psi _j^i=1`$ we get another grading, that by fermionic charge. Therefore, $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})`$ is bigraded and this bigrading extends to the entire sheaf:
$$\mathrm{\Omega }_^N^{ch}=_{m=\mathrm{}}^+\mathrm{}_{n=0}^+\mathrm{}\mathrm{\Omega }_{^N,n}^{ch,m}.$$
$`(2.28)`$
Next, we discuss “tensor” properties of $`\mathrm{\Omega }_^N^{ch}`$. We identify $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{})`$ with $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})_0`$ by identifying $`b^i`$ with $`b_0^i`$ and $`db^i`$ with $`\varphi _0^i`$. This identification extends to the isomorphism (2.20).
The structure of higher conformal weight components is more complicated, but here is what we can say about the component of conformal weight 1. Consider the following elements of $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_{^N,1}^{ch,0})`$:
$$e_{ij}=b_0^{i1}a_1^{j1}+\varphi _0^{i1}\psi _1^{j1},i,j1,$$
$`(2.29a)`$
$$e_{1j}=a_1^{j1},j1$$
$`(2.29b)`$
$$e_{i1}=\underset{l=1}{\overset{N}{}}b_0^{i1}b_0^la_1^l\underset{l=1}{\overset{N}{}}b_0^{i1}\varphi _0^l\psi _1^l$$
$$\underset{l=1}{\overset{N}{}}b_0^l\varphi _0^{i1}\psi _1^l,i1.$$
$`(2.29c)`$
It was checked in \[MS1\] III that these elements come from $`H^0(^N,\mathrm{\Omega }_{^N,1}^{ch,0})\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_{^N,1}^{ch,0})`$ and that the Fourier components of the corresponding fields span a Lie subalgebra of $`\text{End}(\mathrm{\Omega }_^N^{ch})`$ isomorphic to the loop algebra $`Lsl_{N+1}=sl_{N+1}[t,t^1]`$. Therefore,
$$Lsl_{N+1}\text{End}(\mathrm{\Omega }_^N^{ch}),$$
$`(2.30a)`$
so that
$$sl_{N+1}H^0(^N,\mathrm{\Omega }_{^N,1}^{ch,0}),E_{ij}e_{ij}.$$
$`(2.30b)`$
is a morphism of $`sl_{N+1}`$-modules, where $`E_{ij}`$, $`ij`$, $`1i,jN+1`$ are the standard generators of $`sl_{N+1}`$, and $`sl_{N+1}`$ operates on $`H^0(^N,\mathrm{\Omega }_{^N,1}^{ch,0})`$ by means of the composite map $`sl_{N+1}\stackrel{}{}sl_{N+1}1\widehat{sl}_{N+1}`$
Elements (2.29a-c) have fermionic charge 0. For the fermionic charge $`N+1`$ component there is an isomorphism:
$$\mathrm{\Omega }_^N^1\mathrm{\Omega }_^N^N\stackrel{}{}\mathrm{\Omega }_{^N,1}^{ch,N+1}.$$
$`(2.31)`$
Over $`𝕋^N`$ it is defined by the assignment
$$f_i(b^1,\mathrm{},b^N)db^i(db^1db^2\mathrm{}db^N)f_i(b_0^1,\mathrm{},b_0^N)\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N,$$
$$f_i(b^1,\mathrm{},b^N)[(b^1)^{\pm 1},\mathrm{},(b^N)^{\pm 1}].$$
Isomorphism (2.31) induces the isomorphism
$$H^N(^N,\mathrm{\Omega }_^N^1\mathrm{\Omega }_^N^N)\stackrel{}{}H^N(^N,\mathrm{\Omega }_{^N,1}^{ch,N+1}).$$
$`(2.32)`$
By the Serre duality,
$$H^N(^N,\mathrm{\Omega }_^N^1\mathrm{\Omega }_^N^N)\stackrel{}{}H^0(^N,𝒯)^{},$$
$`(2.33)`$
where $`𝒯`$ is the tangent sheaf. The Lie algebra $`sl_{N+1}`$ operates on $`^N`$, therefore there arises the map $`sl_{N+1}H^0(^N,𝒯)^{}`$, which is well known to be an isomorphism. Hence, (2.33) combined with (2.32) rewrites as follows
$$H^N(^N,\mathrm{\Omega }_{^N,1}^{ch,N+1})\stackrel{}{}sl_{N+1}.$$
$`(2.34)`$
This map is an isomorphism of $`sl_{N+1}`$-modules, and it is not hard to find a Cech cochain representing a highest weight vector of $`H^N(^N,\mathrm{\Omega }_^N^{ch,N+1})_1`$, that is, a non-zero vector $`v`$ satisfying
$$E_{ij}v=0,i<j.$$
$`(2.35)`$
If we denote by $`U_i`$ the open subset of $`^N`$ satisfying $`x_i0`$, then $`\{U_0,\mathrm{},U_N\}`$ is an affine cover of $`^N`$, so that $`𝕋^N=U_0U_1\mathrm{}U_N`$. The $`N`$-th term of the Cech complex equals, therefore, $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})`$, and it is an exercise to check that
$$(b_0^1)^1(b_0^2)^1\mathrm{}(b_0^{N1})^1(b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N$$
$`(2.36)`$
represents a highest weight vector of $`H^N(^N,\mathrm{\Omega }_{^N,1}^{ch,N+1})`$.
Observe that another copy of $`sl_{N+1}`$ we have discovered earlier has $`e_{1N+1}`$ for its highest weight vector, see (2.29b,2.30). The assertion crucial for our proof is that $`d_{}`$ sends one highest weight vector to another:
$$d_{}((b_0^1)^1b_0^2)^1\mathrm{}b_0^{N1})^1b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N)=e_{1N+1}H^0(^N,\mathrm{\Omega }_{^N,1}^{ch,0}).$$
$`(2.37)`$
Lemma 2.6 follows from (2.37) easily. To explain this implication we have to digress on elementary representation theory of $`Lsl_{N+1}`$.
Consider the decomposition
$$Lsl_{N+1}=L_{}sl_{N+1}sl_{N+1}L_+sl_{N+1},$$
where
$$L_\pm sl_{N+1}=t^{\pm 1}[t^{\pm 1}].$$
Let $`L_{}sl_{N+1}=sl_{N+1}L_+sl_{N+1}.`$ Any $`sl_{N+1}`$-module becomes an $`L_{}sl_{N+1}`$-module if the action of $`sl_{N+1}`$ is extended to the entire $`L_{}sl_{N+1}`$ by the requirement $`L_+sl_{N+1}0`$. Therefore for any $`sl_{N+1}`$-module $`U`$ there arises the Weyl module, denoted $`𝕎_U`$ and defined as follows:
$$𝕎_U=\text{Ind}_{L_{}sl_{N+1}}^{Lsl_{N+1}}U.$$
The Weyl module induced from the trivial representation, $`𝕎_{}`$, is well-known to be a conformal vertex algebra due to \[FZ\], see also \[K\] 4.7. Therefore, it has vacuum vector, 1, and Virasoro element, $`L_2^{aff}\text{1}`$. Other Weyl modules are modules over $`𝕎_{}`$. This means, in particular, that Fourier components $`L_i^{aff}`$ act on Weyl modules. The action of $`L_0^{aff}`$ is diagonalizable and defines a grading on each Weyl module also called the grading by conformal weight. The aim of this digression was to formulate the following well-known (and easily derived from the Kac-Kazhdan equations) assertion:
$$I𝕎_{sl_{N+1}}\text{ is a }\text{proper}\text{ }Lsl_{N+1}\text{-submodule }I(𝕎_{sl_{N+1}})_2=\{0\},$$
$`(2.38)`$
where $`𝕎_{sl_{N+1}}`$ stands for the Weyl module induced from the adjoint representation, and $`(𝕎_{sl_{N+1}})_2`$ is its conformal weight 2 component.
Return to the proof of Lemma 2.6. Due to (2.30a), $`H^i(^N,\mathrm{\Omega }_^N^{ch,m})`$ is an $`Lsl_{N+1}`$-module for all $`i`$ and $`m`$. The component $`H^N(^N,\mathrm{\Omega }_{^N,1}^{ch,N+1})`$ is an $`L_{}sl_{N+1}`$-module isomorphic to $`sl_{N+1}`$, see (2.34), on which $`L_+sl_{N+1}`$ acts trivially because $`H^N(^N,\mathrm{\Omega }_{^N,m}^{ch,N+1})=0`$ for all $`m<1`$. By the universality property of induced modules, $`𝕎_{sl_{N+1}}`$ maps onto the $`Lsl_{N+1}`$-submodule of $`H^N(^N,\mathrm{\Omega }_^N^{ch,N+1})`$ generated by $`H^N(^N,\mathrm{\Omega }_{^N,1}^{ch,N+1})`$. Denote this submodule $`\widehat{𝕎}_{sl_{N+1}}`$.
Similarly, $`H^0(^N,\mathrm{\Omega }_^N^{ch,0})`$ is an $`Lsl_{N+1}`$-module, and the $`Lsl_{N+1}`$-submodule generated by 1 is a quotient of $`𝕎_{}`$. This quotient contains yet another submodule, the one generated by $`sl_{N+1}H^0(^N,\mathrm{\Omega }_^N^{ch})_1`$, see (2.30b), to be denoted $`\widehat{\widehat{𝕎}}_{sl_{N+1}}`$. This submodule, again for the same reason, is a quotient of $`𝕎_{sl_{N+1}}`$. By definition, the above mentioned Virasoro element $`L_2^{aff}\text{1}`$ belongs to $`(\widehat{\widehat{𝕎}}_{sl_{N+1}})_2`$
We are practically done. It is easy to derive from \[B\] that
$$d_{}:H^N(^N,\mathrm{\Omega }_^N^{ch,N+1})H^0(^N,\mathrm{\Omega }_^N^{ch,0})$$
is an $`Lsl_{N+1}`$-morphism. Equality (2.37) then means that $`d_{}(\widehat{𝕎}_{sl_{N+1}})\widehat{\widehat{𝕎}}_{sl_{N+1}}`$ is non-zero, and is therefore a quotient of $`𝕎_{sl_{N+1}}`$ by a proper submodule. Due to (2.38)
$$(d_{}(\widehat{𝕎}_{sl_{N+1}}))_2=(𝕎_{sl_{N+1}})_2=(\widehat{\widehat{𝕎}}_{sl_{N+1}})_2.$$
Hence $`L_2^{aff}\text{1}d_{}(\widehat{𝕎}_{sl_{N+1}})`$. To complete the proof of Lemma 2.6 it remains to check that the affine Virasoro element, $`L_2^{aff}\text{1}`$, coincides with $`L_2\text{1}`$ and this is easy.
Finally we have to prove (2.37). The difficulty with computation of
$$d_{}((b_0^1)^1(b_0^2)^1\mathrm{}(b_0^{N1})^1(b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N)$$
lies in that the operator $`d_{}`$ is defined in terms of the vertex algebra $`V_L`$, while
$$(b_0^1)^1(b_0^2)^1\mathrm{}(b_0^{N1})^1(b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N$$
is an element of $`\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})`$. The vertex algebra embedding
$$\mathrm{\Gamma }(𝕋^N,\mathrm{\Omega }_^N^{ch})V_L,$$
an important ingredient of Borisov’s proof of Theorem 1.3, is determined by the rules
$$(b_0^i)^\pm e^{\pm B^i},\varphi _0^i\mathrm{\Phi }_0^ie^{B^i},\psi _1^i\mathrm{\Psi }_1^ie^{B^i},$$
$`(2.39a)`$
$$a_1^iA_1^ie^{B^i}\mathrm{\Phi }_0^i\mathrm{\Psi }_1^ie^{B^i},$$
$`(2.39b)`$
$$xXL_1xL_1X,$$
$`(2.39c)`$
$$xX,yYx_{(1)}yX_{(1)}Y.$$
$`(2.39d)`$
These rules imply
$$(b_0^1)^1(b_0^2)^1\mathrm{}(b_0^{N1})^1(b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^Ne^{B^N}\mathrm{\Phi }_1^N\mathrm{\Phi }_0^1\mathrm{\Phi }_0^2\mathrm{}\mathrm{\Phi }_0^N.$$
It follows from Borisov’s proof of Theorem 1.3 that an element of $`V_L`$ representing the class of
$$(b_0^1)^1(b_0^2)^1\mathrm{}(b_0^{N1})^1(b_0^N)^3\varphi _1^i\varphi _0^1\varphi _0^2\mathrm{}\varphi _0^N$$
can be chosen to be equal to
$$(\mathrm{\Psi }_1^Ne^{A^N})_{(0)}(\mathrm{\Psi }_1^{N1}e^{A^{N1}})_{(0)}\mathrm{}(\mathrm{\Psi }_1^1e^{A^1})_{(0)}e^{B^N}\mathrm{\Phi }_1^N\mathrm{\Phi }_0^1\mathrm{\Phi }_0^2\mathrm{}\mathrm{\Phi }_0^N.$$
The formulas of 1.1-2 imply that $`(\mathrm{\Psi }_1^ie^{A^i})_{(0)}`$, $`1iN1`$, simply erases $`\mathrm{\Phi }_0^i`$. Hence
$$(\mathrm{\Psi }_1^{N1}e^{A^N})_{(0)}(\mathrm{\Psi }_1^{N2}e^{A^{N1}})_{(0)}\mathrm{}(\mathrm{\Psi }_1^1e^{A^1})_{(0)}e^{B^N}\mathrm{\Phi }_1^N\mathrm{\Phi }_0^1\mathrm{\Phi }_0^2\mathrm{}\mathrm{\Phi }_0^N$$
$$=\mathrm{\Phi }_1^N\mathrm{\Phi }_0^Ne^{B^N+A^1+A^2+\mathrm{}A^{N1}}.$$
The calculation of the last operation is a little more tedious, but also straightforward; the result is this:
$$(\mathrm{\Psi }_1^Ne^{A^N})_{(0)}(\mathrm{\Phi }_1^N\mathrm{\Phi }_0^Ne^{B^N+A^1+A^2+\mathrm{}A^{N1}})$$
$$(\mathrm{\Psi }_1^N\mathrm{\Phi }_1^N\mathrm{\Phi }_0^N\mathrm{\Phi }_1^NA_1^N+\frac{1}{2}\mathrm{\Phi }_0^NA_2^N+\frac{1}{2}\mathrm{\Phi }_0^N(A_1^N)^2)e^{B^N+A^1+A^2+\mathrm{}A^N}.$$
$`(2.40)`$
To complete our calculation we have to apply $`d_{}`$ to this element. Observe that this element comes from the interior of the cone spanned by $`A^1,\mathrm{},A^N`$ and has height $`N`$. It follows from the definition of the spectral sequence and (1.1) or (2.21) that on this element $`d_{}`$ equals
$$((\mathrm{\Psi }_1^1+\mathrm{\Psi }_1^2+\mathrm{}+\mathrm{\Psi }_1^N)e^{A^1A^2\mathrm{}A^N})_{(0)}.$$
Indeed, it is precisely the component of Borisov’s differential (1.1) that decreases the height of the element (2.40). (By the way, it decreases it by $`N`$, which explains the assertion (2.23c).) Another calculation similar to those performed shows that
$$((\mathrm{\Psi }_1^1+\mathrm{\Psi }_1^2+\mathrm{}+\mathrm{\Psi }_1^N)e^{A^1A^2\mathrm{}A^N})_{(0)}.$$
sends the element (2.40) to
$$A_1^Ne^{B^N}\mathrm{\Phi }_0^N\mathrm{\Psi }_1^Ne^{B^N}.$$
According to (2.39b), the latter element corresponds to $`a_1^N`$ and hence to $`e_{1N+1}`$, see (2.29b), as desired.
§3. Deforming cohomology algebras of hypersurfaces in projective spaces
Let $`^N`$ be a degree $`n<0`$ line bundle, $`^{}^N`$ its dual, $`s\mathrm{\Gamma }(^N,^{})`$ a global section so that its zero locus $`Z(s)^N`$ is a smooth hypersurface. The way Borisov calculates the cohomology of the chiral de Rham complex over $`Z(s)`$ is as follows.
Extend the lattice $`(L,(.,.))`$ introduced in 1.1 to the lattice $`(\widehat{L},(.,.))`$ so that
$$\widehat{L}=LA^uB^u,(A^u,B^u)=1,(A^u,L)=0,(B^u,L)=0.$$
There arises the corresponding lattice vertex algebra $`V_{\widehat{L}}`$. Observe that any subset $`L^{}\widehat{L}`$ closed under addition gives rise to the vertex subalgebra $`V_L^{}V_L`$ generated by $`𝔥_{\widehat{L}}`$ and $`Cl_{\widehat{L}}`$ from the highest weight vectors $`e^\beta ,\beta L^{}`$; see 1.1-1.2. In our geometric situation let $`\widehat{L}_n`$ be the span of $`B^i`$ ($`i=1,\mathrm{},N`$), $`B^u`$ with arbitrary integral coefficients and $`A^i`$ ($`i=1,\mathrm{},N`$), $`A^u`$, $`nA^uA^1\mathrm{}A^N`$ with nonnegative integral coefficients.
The vertex algebra $`V_{\widehat{L}_n}`$ affords a degeneration, $`V_{\widehat{L}_n}^\mathrm{\Sigma }`$, and includes in a family, $`V_{\widehat{L}_n,q}`$, $`q`$, in the same way the algebra $`V_L`$ did, see 1.3, 2.1. To construct $`V_{\widehat{L}_n}^\mathrm{\Sigma }`$, consider the following $`N+1`$ elements of $`\widehat{L}`$ : $`\xi _1=A^1,\xi _2=A^2,\mathrm{},\xi _N=A^N,\xi _{N+1}=nA^uA^1A^2\mathrm{}A^N`$. Define the cone $`\mathrm{\Delta }_i`$ to be the set of all non-negative integral linear combinations of the elements $`\xi _1,\mathrm{},\xi _{i1},\xi _{i+1},\mathrm{},\xi _{N+1},A^u`$ and let $`\mathrm{\Sigma }=\{\mathrm{\Delta }_1,\mathrm{},\mathrm{\Delta }_{N+1}\}`$. The vertex algebra $`V_{\widehat{L}_n}^\mathrm{\Sigma }`$ is now defined by repeating word for word the definition of $`V_L^\mathrm{\Sigma }`$ in 1.3.
Similarly, the family $`V_{\widehat{L}_n,q}`$, $`q0`$, is defined by repeating word for word the definition of $`V_{L,q}`$ in 2.1. This family extends “analytically” to $`q=0`$ if $`nN+1`$ and we again obtain an isomorphism
$$V_{L,0}=V_{\widehat{L}_n}^\mathrm{\Sigma }\text{ if }n<N+1.$$
$`(3.1)`$
(The condition $`n<N+1`$ will be clarified below.)
Borisov’s differential is as follows:
$$D=\{\underset{i=1}{\overset{N}{}}\mathrm{\Psi }^i(z)(e^{A^i}e^{nA^u_jA^j})(z)+\mathrm{\Psi }^u(z)(ne^{nA^u_jA^j}e^{A^u})(z)\}.$$
$`(3.2a)`$
(For the future use let us note that the right hand side of this equality can be rewritten as a sum over lattice points:
$$D=\{\underset{i=1}{\overset{N+1}{}}\mathrm{\Psi }^{\xi _i}(z)e^{\xi _i}(z)\mathrm{\Psi }^u(z)e^{A^u}(z)\},$$
$`(3.2b)`$
where $`\mathrm{\Psi }^{\xi _i}=\mathrm{\Psi }^i`$ ($`iN`$) and $`\mathrm{\Psi }^{\xi _{N+1}}=n\mathrm{\Psi }^u_j\mathrm{\Psi }^j`$.)
It is obvious that $`D\text{End}(V_{\widehat{L}_n,q})`$ and $`D^2=0`$; therefore there arise the cohomology groups $`H_D(V_{\widehat{L}_n,q})`$ and $`H_D(V_{\widehat{L}_n}^\mathrm{\Sigma })=H_D(V_{\widehat{L}_n,0})`$.
Theorem 3.1. (\[B\])
$$H_D(V_{\widehat{L}_n}^\mathrm{\Sigma })=H^{}(,\mathrm{\Omega }_{}^{ch}).$$
Borisov proposes to calculate the chiral de Rham complex over the hypersurface $`Z(s)^N`$ by means of a certain Koszul-type resolution of the complex $`\mathrm{\Omega }_{}^{ch}`$. The combinatorial data that determine $`s\mathrm{\Gamma }(^N,^{})`$ consists of the finite set
$$\mathrm{\Delta }^{}=\{\beta =B^u+\underset{j=1}{\overset{N}{}}n_jB^j\text{ s.t. }(\beta ,\xi _i)0,i=1,2,\mathrm{},N+1\},$$
$`(3.3)`$
and a function
$$g:\mathrm{\Delta }^{}_{}.$$
$`(3.4)`$
Define
$$K_g=\underset{\beta \mathrm{\Delta }^{}}{}g(\beta )\mathrm{\Phi }^\beta (z)e^\beta (z),$$
$`(3.5)`$
where $`\mathrm{\Phi }^\beta =\mathrm{\Phi }^u+_jn_j\mathrm{\Phi }^j`$ provided $`\beta =B^u+_jn_jB^j`$. It is easy to see that
$$K_g\text{End}(V_{\widehat{L}_n,q}),K_g^2=0,[K_g,D]=0.$$
Therefore, there arise the cohomology groups $`H_{D+K_g}(V_{\widehat{L}_n,q})`$ and $`H_{D+K_g}(V_{\widehat{L}_n}^\mathrm{\Sigma })=H_{D+K_g}(V_{\widehat{L}_n,0})`$.
Theorem 3.2. (\[B\])
$$H_{D+K_g}(V_{\widehat{L}_n}^\mathrm{\Sigma })=H^{}(Z(s),\mathrm{\Omega }_{Z(s)}^{ch}).$$
All the vertex algebras in sight being topological (see the beginning of 2.2), Theorem 3.2 and the main result of \[MSV\] give
$$H^{}(Z(s))=H_{Q_0}(H_{D+K_g}(V_{\widehat{L}_n}^\mathrm{\Sigma })),$$
$`(3.6a)`$
or, equivalently,
$$H^{}(Z(s))=H_{D+K_g}(V_{\widehat{L}_n}^\mathrm{\Sigma })_0,$$
$`(3.6b)`$
where $`H_{D+K_g}(V_{\widehat{L}_n}^\mathrm{\Sigma })_0`$ stands for the kernel of $`L_0`$.
This prompts the following
Conjecture 3.3. If $`n<N+1`$, then the algebra $`H_{D+K_g}(V_{\widehat{L}_n,q})_0`$ is isomorphic to the quantum cohomology algebra of $`Z(s)`$.
Unfortunately we do not have a proof of this conjecture; we cannot even prove that $`H_{D+K_g}(V_{\widehat{L}_n,q})_0`$ is a deformation of $`H^{}(Z(s))`$. What we know is collected in the following
Proposition 3.4. (i) The element $`e^{nA^u_jA^j}`$ satisfies
$$(D+K_g)(e^{nA^u_jA^j})=0,$$
and, therefore, determines an element of $`H_{D+K_g}(V_{\widehat{L}_n,q})_0`$ for all $`q`$. If $`q=0`$, then this element, considered as an element of $`H^{}(Z(s))`$ (see (3.6b)), is proportional to the cohomology class of a hyperplane section.
(ii) Due to (i), $`e^{nA^u_jA^j}`$ generates a subalgebra of $`H_{D+K_g}(V_{\widehat{L}_n,q})_0`$ to be denoted $`𝒜_q`$. This subalgebra is a deformation of $`𝒜_0`$.
(iii) If $`Z(s)`$ is a hyperplane (i.e. $`n=1`$), then Conjecture 3.3 is correct.
(iv) If $`Z(s)`$ is a non-degenerate quadric in $`^3`$, then $`H_{D+K_g}(V_{\widehat{L}_n,q})_0`$ is isomorphic to $`[x,y]/(x^21,y^21).`$. Hence Conjecture 3.3 is true in this case.
Since these results are by no means complete, we shall confine ourselves to sketching a proof of Proposition 3.4. The first part of assertion (i) is a result of the obvious calculation using the formulas of 1.1-1.2. The fact that at $`q=0`$ the element $`e^{nA^u_jA^j}`$ is proportional to the cohomology class of a hyperplane section follows from Borisov’s proof of Theorem 3.2; this observation is completely analogous to the one made in the end of 2.3.
To prove (ii) observe that we have a constant family of vector spaces $`V_{\widehat{L}_n,q}`$, $`q`$, with differential $`D+K_g`$ depending on $`q`$. At $`q=0`$ the complex $`(V_{\widehat{L}_n,q},D+K_g)`$ degenerates in Borisov’s complex $`(V_{\widehat{L}_n}^\mathrm{\Sigma },D+K_g)`$. As it always happens in situations of this kind, the differential $`d=D+K_g`$ breaks in a sum $`d=d_{}(q)+d_+`$ so that $`[d_{}(q),d_+]=0`$, $`d_{}(0)=0`$, and $`d_+`$ equals Borisov’s differential on $`V_{\widehat{L}_n}^\mathrm{\Sigma }`$. There arises a spectral sequence converging to $`H_{D+K_g}(V_{\widehat{L}_n,q})`$ with 1-st term equal to $`H^{}(Z(s),\mathrm{\Omega }_{Z(s)}^{ch})`$. The 2-nd term equals the cohomology of the complex $`(H^{}(Z(s),\mathrm{\Omega }_{Z(s)}^{ch}),d^{(1)})`$ with $`d^{(1)}=d_{}(q))`$. It remains to show that
$$𝒜_0\text{Ker}d^{(r)},𝒜_0\text{Im}d^{(r)}=0,r1.$$
$`(3.7)`$
All these spaces are subquotients of the subalgebra of $`V_{\widehat{L}_n,0}`$ generated by $`e^0`$, $`e^{A^i}`$ ($`i=1,\mathrm{},N`$), $`e^{A^u}`$, and $`\mathrm{\Phi }_0^i`$ ($`i=1,\mathrm{},N`$), $`\mathrm{\Phi }_0^u`$, the product being equal to <sub>(-1)</sub>. This is a supercommutative algebra isomorphic to
$$[x_1,\mathrm{},x_N,T,u;\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N,\mathrm{\Phi }_u]/(x_1x_2\mathrm{}x_NT),$$
where we let $`x_i=e^{A^i}`$, $`u=e^{A^u}`$, $`T=e^{nA^u_jA^j}`$, $`\mathrm{\Phi }_i,\mathrm{\Phi }_u`$ being the corresponding grassman variables. (All this is completely analogous to our discussion in the end of 2.1.) Formula (3.2a) says that when restricted to this space Borisov’s differential $`D`$ coincides with the Koszul differential associated with the sequence $`x_iT,unT`$ ($`i=1,\mathrm{},N`$) and our space quickly shrinks to
$$[T]/(T^{N+1}),$$
on which (3.7) is obviously true at least when $`r=1`$. If $`r2`$, then the first part of (3.7) is obviously true and the second follows from a simple dimensional argument.
Before turning to (iii) let us note that a quantum version of this argument gives:
$$𝒜_q\text{ is a quotient of }[T]/(T^{N+1}q^{N+1n}n^nT^n).$$
$`(3.8)`$
Indeed, again by definition (as in the end of 2.1), the subalgebra of $`V_{\widehat{L}_n,q}`$ generated by $`e^{A^i}`$ ($`i=1,\mathrm{},N`$), $`e^{A^u}`$, and $`\mathrm{\Phi }_0^i`$ ($`i=1,\mathrm{},N`$), $`\mathrm{\Phi }_0^u`$ is isomorphic to
$$[x_1,\mathrm{},x_N,T,u;\mathrm{\Phi }_1,\mathrm{},\mathrm{\Phi }_N,\mathrm{\Phi }_u]/(x_1x_2\mathrm{}x_NTq^{N+1n}u^n),$$
and the restriction of $`D`$ to this supercommutative algebra coincides with the Koszul differential associated with the regular sequence $`x_iT,unT`$ ($`i=1,\mathrm{},N`$). The relation (3.8) follows at once. By the way, the appearance of $`N+1n`$ as a power of $`q`$ in (3.8) explains why the condition $`n<N+1`$ was imposed in (3.1).
Return to the proof of (iii). In this case the quantum cohomology algebra is isomorphic to the algebra of functions on an $`N`$-point set. Because of (ii), $`𝒜_q`$ is isomorphic to $`[T]/p(T)`$, $`\text{deg}p(T)=N`$, and, because of (3.8), $`p(T)`$ divides $`T^{N+1}q^NT`$. The latter has no multiple roots. Hence $`𝒜_q`$ is also the algebra of functions on an $`N`$-point set.
(iv) follows from the same spectral sequence that was used for the proof of (ii): due to (3.6b) $`H^{}(Z(s),\mathrm{\Omega }_{Z(s)}^{ch})_0=[x,y]/(x^2,y^2)`$ and the elements $`x,y`$ are annihilated by all higher differentials because on the one hand $`x,yH^1(Z(s),\mathrm{\Omega }_{Z(s)}^{ch})`$ and on the other hand it is true in general that all $`d^{(r)}`$, $`r1`$, send $`H^1(Z(s),\mathrm{\Omega }_{Z(s)}^{ch})`$ to 0. The rest follows from (3.8), which in this case reads as follows:
$$T^44q^2T^2=0.$$
Remarks. (i) By Corollary 9.3 of \[G\], the cohomology class $`p`$ of a hyperplane section satsifies in the quantum cohomology of $`Z(s)`$ the relation
$$p^N=qn^np^{n1}.$$
The amusing similarity between this equality and (3.8) suggests that $`𝒜_q`$ might be equal to $`[T]/(T^Nq^{N+1n}n^nT^{n1})`$.
(ii) Borisov’s suggestion to treat the mirror symmetry as a flip interchanging $`A`$’s and $`B`$’s seems to be working in our “quantized” situation as well. Compare (3.5) with (3.2b) to note that $`D`$ and $`K_g`$ are sums over two sets of lattice points defined by self-dual condition (3.3). Hence the $`AB`$ flip changes $`D`$ to a similar differential to be associated with the mirror partner of $`Z(s)`$ lying in another toric manifold, see the next section. Of course the vertex algebra $`V_{\widehat{L}_n}`$ bears a certain asymmetry, since not all elements of the type $`e^{_jn_jA^j+n_uA^u}`$ are allowed, but Borisov’s “transition to the whole lattice” (see Theorem 8.3 in \[B\]) and the above spectral sequence seem to straighten things out.
§4. Quantum cohomology of toric varieties
Let us briefly explain how the constructions and results of section 2 carry over to an arbitrary smooth compact toric variety of dimension $`N`$. Each such variety is determined by a complete regular fan in $`L_A`$. This and other relevant concepts can be defined as follows (see \[D, Bat\] for details).
4.1 Let $`IL_A`$. The cone generated by $`I`$ is said to be the set of all non-negative integral combinations of elements of $`I`$ and is denoted $`\mathrm{\Delta }_I`$.
A cone generated by part of a basis of $`L_A`$ is called regular.
A complete regular fan $`\mathrm{\Sigma }`$ is defined to be a collection of regular cones $`\{\sigma _1,\mathrm{},\sigma _s\}`$ so that the following conditions hold:
(i) If $`\sigma ^{}`$ is a face of $`\sigma \mathrm{\Sigma }`$, then $`\sigma ^{}\mathrm{\Sigma }`$;
(ii) If $`\sigma ,\sigma ^{}\mathrm{\Sigma }`$, then $`\sigma \sigma ^{}`$ is a face of $`\sigma `$;
(iii) (the completeness condition) $`L_A=\sigma _1\mathrm{}\sigma _s`$.
We skip the construction of the smooth compact toric manifold $`X_\mathrm{\Sigma }`$ attached to a regular complete fan $`\mathrm{\Sigma }`$ referring the reader to \[D\], but formulate Batyrev’s result on $`H^2(X_\mathrm{\Sigma },)`$, see \[Bat\].
A function $`\varphi :L_A`$ is called piecewise linear if its restriction to any cone in $`\mathrm{\Sigma }`$ is a morphism of abelian groups. Denote by $`PL(\mathrm{\Sigma })`$ the space of all piecewise linear functions.
Let $`G(\mathrm{\Sigma })=\{\xi _1,\mathrm{},\xi _n\}`$ be the set of the generators of all 1-dimensional cones in $`\mathrm{\Sigma }`$. Since each piecewise linear function is determined by its values on $`\xi _i`$ ($`i=1,\mathrm{},n`$), $`PL(\mathrm{\Sigma })`$ is an $`n`$-dimensional real vector space. It contains the $`N`$-dimensional subspace of globally linear functions; the latter is naturally isomorphic to $`L_B_{}`$.
Theorem 4.1 (\[Bat\])
$$H^2(X_\mathrm{\Sigma },)=PL(\mathrm{\Sigma })/L_B_{}.$$
4.2 Let us return to the vertex algebra $`V_L`$. Having fixed an arbitrary $``$-valued function $`\varphi `$ on $`L_A`$, we proceed in much the same way as in 2.1.
Define the linear automorphism
$$t_\varphi :V_LV_L$$
by the formula
$$t_\varphi (xe^{_im_iB^i+_in_iA^i})=e^{\varphi (_in_iA^i)}xe^{_im_iB^i+_in_iA^i}.$$
$`(4.1)`$
Define $`V_{L,\varphi }`$ to be the vertex algebra equal to $`V_L`$ as a vector space with the following n-th product:
$$(xe^{_im_iB^i+_in_iA^i})_{(n),\varphi }(ye^{_im_i^{}B^i+_in_i^{}A^i})$$
$$=t_\varphi ^1(t_\varphi (xe^{_im_iB^i+_in_iA^i})_{(n)}t_\varphi (ye^{_im_i^{}B^i+_in_i^{}A^i})).$$
$`(4.2)`$
By definition,
$$t_\varphi :V_{L,\varphi }V_L$$
is a vertex algebra isomorphism. This provides us with a constant family of vertex algebras parametrized by $`\varphi `$ and we would like to study the behavior of this family as $`\varphi `$ tends to $`\mathrm{}`$. For this we have to impose certain restrictions on $`\varphi `$.
Following \[Bat\], call a piecewise linear function $`\varphi `$ convex if
$$\varphi (x)+\varphi (y)\varphi (x+y)\text{ all }x,yL_A.$$
$`(4.3)`$
The cone of all convex piecewise linear functions descends to the cone in $`H^2(X_\mathrm{\Sigma },)=PL(\mathrm{\Sigma })/L_B_{}`$, see Theorem 4.1. Denote this cone by $`K(\mathrm{\Sigma })`$ and its interior by $`K^0(\mathrm{\Sigma })`$. $`K^0(\mathrm{\Sigma })`$ consists of classes of all strictly convex piecewise linear functions, that is, of all those functions $`\varphi `$ for which equality in (4.1) is achieved if and only if $`x`$ and $`y`$ belong to the same cone in $`\mathrm{\Sigma }`$.
We see immediately that
(i) if $`\varphi `$ is convex piecewise linear, then the operations
$${}_{(n),\mathrm{}\varphi }{}^{}=\underset{\tau +\mathrm{}}{lim}{}_{(n),\tau \varphi }{}^{},n$$
are well defined and satisfy the Borcherds identities; denote the vertex algebra arising in this way by $`V_{L,\mathrm{}\varphi }`$;
(ii) if $`\varphi `$ is strictly convex piecewise linear, then $`V_{L,\mathrm{}\varphi }`$ is isomorphic to Borisov’s algebra $`V_L^\mathrm{\Sigma }`$.
These assertions mean that the family $`V_{L,\varphi }`$ produces a deformation of $`V_L^\mathrm{\Sigma }`$ with base equal to the cone of strictly convex piecewise linear functions. It is also immediate to see that if $`\varphi \varphi ^{}`$ is a linear function, then the two deformations $`V_{L,\tau \varphi }`$ and $`V_{L,\tau \varphi ^{}}`$, $`\tau 0`$, are equivalent. Therefore we have obtained the family of vertex algebras $`V_{L,\varphi }`$, $`\varphi K^0(\mathrm{\Sigma })`$, which is a deformation of $`V_L^\mathrm{\Sigma }`$ with base $`K^0(\mathrm{\Sigma })`$.
Denote by $`QH_\varphi ^{}(X_\mathrm{\Sigma },)`$ the quantum cohomology of $`X_\mathrm{\Sigma }`$ as defined in section 5 of \[Bat\]. Borisov’s differential is as follows
$$D=\underset{i=1}{\overset{N}{}}\mathrm{\Psi }^i(z)(\underset{j=1}{\overset{n}{}}(B^i,\xi _j)e^{\xi _j}(z)),$$
$`(4.4)`$
where $`\{\xi _1,\mathrm{},\xi _n\}`$ is the set of generators of all 1-dimensional cones in $`\mathrm{\Sigma }`$.
Theorem 4.2
$$H_{Q_0+D}(V_{L,\varphi })=QH_\varphi ^{}(X_\mathrm{\Sigma },).$$
Sketch of Proof. First of all,
$$(Q_0)^2=0,D^2=0,[Q_0,D]=0.$$
(The first two of these assertions are obvious, the last one is obtained in the same way as (2.4).) Hence there arises a spectral sequence completely analogous to the one used in 2.4. It converges and collapses:
$$H_{Q_0+D}(V_{L,\varphi })=H_D(H_{Q_0}(V_{L,\varphi }));$$
this is done in exactly the same way as in the proof of Theorem 2.4.
In part 1) of the proof of Theorem 2.3 the space $`H_{Q_0}(V_{L,\varphi })`$ was shown to be equal to the group algebra $`[L_A]`$ extended by grassman variables $`\mathrm{\Phi }_0^i`$ ($`i=1,\mathrm{},N`$). Thus $`H_{Q_0}(V_{L,\varphi })`$ is a Koszul complex, and the restriction of $`D`$ to this space equals the Koszul differential associated with the sequence
$$\underset{j=1}{\overset{n}{}}(B^i,\xi _j)e^{\xi _j},i=1,\mathrm{},N.$$
$`(4.5)`$
Therefore, $`H_D(H_{Q_0}(V_{L,\varphi }))`$ is the corresponding “Koszul cohomology”.
On the other hand, Batyrev defines $`QH_\varphi ^{}(X_\mathrm{\Sigma },)`$ to be the polynomial ring $`[z_1,\mathrm{},z_n]`$ modulo the sum of two ideals denoted $`P(\mathrm{\Sigma })`$ and $`Q_\varphi (\mathrm{\Sigma })`$. It follows from the proof of either Theorem 9.5 or Theorem 8.4 in \[Bat\] that
$$[L_A]=[z_1,\mathrm{},z_n]/Q_\varphi (\mathrm{\Sigma }).$$
Under this identification, the image of the ideal $`P(\mathrm{\Sigma })`$ in $`[L_A]`$ is generated by the elements (4.5) as follows from the comparison of (4.5) above and Definition 3.7 in \[Bat\]. The ring $`[L_A]`$ is Cohen-Macaulay; hence the sequence (4.5) is regular.
References
\[Bat\] V. Batyrev, Quantum cohomology rings of toric manifolds, Journées de Géométrie Algébrique d’Orsay (Orsay 1992), Astérisque, 218 (1993), 9-34; alg-geom/9310004.
\[B\] L. Borisov, Vertex algebras and Mirror symmetry, math.AG/9809094.
\[D\] V.I. Danilov, The geometry of toric varieties, Uspekhi Mat.Nauk, 33, No. 2 (1978), 85-134 (Russian); Russian Math. Surveys, 33, No. 2 (1978), 97-154.
\[FZ\] I.B. Frenkel, Y. Zhu, Vertex operator algebras associated to representations of affine and Virasoro algebra, Duke Math. J., vol. 66, No. 1, (1992), 123-168
\[G\] A. Givental, Equivariant Gromow-Witten invariants, Internat. Math. Res. Notices, 13 (1996), 613-663.
\[MSV\] F. Malikov, V. Schechtman, A. Vaintrob, Chiral de Rham complex, Comm. Math. Phys., 204 (1999), 439-473.
\[MS1\] F. Malikov, V. Schechtman Chiral de Rham complex. II Amer. Math. Soc. Transl. (2), vol. 194 (1999), 149-188
\[MS2\] F. Malikov, V. Schechtman Chiral Poincaré duality Math. Res. Lett. vol. 6 (1999), 533-546
F.M.: Department of Mathematics, University of Southern California, Los Angeles, CA 90089, USA; fmalikovmathj.usc.edu
V.S.: IHES, 35 Route de Chartres, 91440 Bures-sur-Yvette, France; vadikihes.fr |
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In this Letter we consider a simple toy model which illustrates a generic problem in string-inspired supergravity models with an anomalous $`U(1)_X`$. A large number of matter multiplets charged under $`U(1)_X`$ remain massless above the supersymmetry-breaking scale because of degeneracy of vacua solving the D-flatness conditions. We refer to these multiplets as “D-moduli.” For example, in the model described in Section 4.2 of , that we will refer to here as the FIQS model, there are 26 massless chiral multiplets associated with this degeneracy. In the toy model considered here, we find that the degeneracy is partially broken by the introduction of a generic supersymmetry-breaking term in such a way that the overall vacuum energy vanishes. However many scalar fields remain massless even after supersymmetry breaking. Specifically, in our toy model with the convention<sup>1</sup><sup>1</sup>1Our charge normalization is such that $`\mathrm{Tr}T_a^2=\frac{1}{2}`$ if $`T_a`$ is a generator of $`SU(N)`$. that $`\mathrm{Tr}Q_X>0`$, where $`Q_X`$ is the generator of $`U(1)_X`$, the only remaining flat directions are those for which a linear combination of fields with the lowest value $`q_0`$ of the $`U(1)_X`$ charge acquires a vev. The corresponding linear combination of chiral multiplets forms a massive vector multiplet with the $`U(1)_X`$ gauge fields, while the orthogonal combinations are massless. In the FIQS model, for example, the lowest $`U(1)_X`$ charge is $`8`$ with a 15-fold degeneracy, so 14 complex scalars (as well as 26 chiral fermions) would be massless if no other symmetries were broken. The scalar fields with $`q_X>q_0`$ would acquire soft masses of the same order as generically expected for squarks and sleptons. However in string-derived models such as this one, the D-moduli are charged under other $`U(1)`$’s, which partially lifts the degeneracy of the vacuum before supersymmetry breaking.
We first consider our toy model in the context of standard supergravity with scalars and their superpartners in chiral multiplets. Then we appeal more explicitly to the linear multiplet formulation of gaugino condensation as the mechanism of supersymmetry breaking, which we refer to as the BGW model. Neither of our models is realistic. For example, we neglect the dependence of the matter Kähler metric on T-moduli (i.e. breathing modes). While including this would considerably complicate the analysis, it seems unlikely to provide a mechanism (at least at tree level) for lifting the degeneracy of the vacuum. We also neglect additional, nonanomalous $`U(1)`$ couplings of the fields with large $`vev`$’s; as discussed below, we do not expect them to lift the vacuum degeneracy completely. We illustrate this point using the FIQS model, which itself should probably be considered a toy model for reasons discussed below. Finally we suggest a possible mechanism for lifting the remaining degeneracy after supersymmetry breaking, and comment on implications for cosmology.
For our toy chiral supergravity model, we assume a Kähler potential
$$K=\underset{i}{}|A_i|^2+\underset{i}{}|B_i|^2+\underset{i}{}|\mathrm{\Phi }_i|^2+K(M,\overline{M}),$$
(1)
a superpotential
$$W=\lambda _{ijk}A_iB_j\mathrm{\Phi }_k+W(M),$$
(2)
and a gauge group $`G_{\text{gauge}}=SU(3)_c\times U(1)_X`$. We denote $`U(1)_X`$ charges by:
$$Q_XA_i=n_iA_i,Q_XB_i=p_iB_i,Q_X\mathrm{\Phi }_i=q_i\mathrm{\Phi }_i,Q_XM_a=0.$$
(3)
We assume $`q_i0`$ and $`|q_i|𝒪(1)`$. We further assume that the $`q_i`$ are such that $`<D_X>=0`$ has a solution with $`<A_i>=<B_i>=0`$; in other words we assume that there are flat directions that allow supersymmetry to remain unbroken in the absence of the nonperturbatively induced superpotential $`W(M)`$, as was found in the string-derived models studied in . Scalar components of the gauge-charged superfields are given by $`a_i=A_i|`$, $`b_i=B_i|`$, $`\varphi _i=\mathrm{\Phi }_i|`$. We take $`A_i`$ and $`B_i`$ to be charged under $`SU(3)_c`$, while $`\mathrm{\Phi }_i`$ and $`M_a`$ are $`SU(3)`$ singlets. E.g., $`A_i`$ is a $`\mathrm{𝟑}`$ while $`B_i`$ is a $`\overline{\mathrm{𝟑}}`$. Couplings such as (2) occur in semi-realistic heterotic orbifold models . When $`\varphi _i0`$, generally required by D-flatness because of the Fayet-Illiopoulos (FI) term associated with the anomalous $`U(1)_X`$, some color triplets acquire large masses. We demand that $`SU(3)_c`$ remain unbroken at all scales. Therefore
$$a_i=b_i=0.$$
(4)
We interpret the gauge singlet superfields $`M_a`$ as moduli; the supersymmetry-breaking superpotential $`W(M)`$ in (2):
$$W(M)|\delta ,$$
(5)
is assumed to be generated by nonperturbative dynamics. Cancellation of the $`U(1)_X`$ anomaly by the Green-Schwarz (GS) mechanism induces an FI term $`\xi `$ in the $`U(1)_X`$ D-term :
$$D_X=\underset{i}{}q_i|\varphi _i|^2+\underset{i}{}n_i|a_i|^2+\underset{i}{}p_i|b_i|^2+\xi ,\xi =\frac{g_X^2}{24\pi ^2}\mathrm{Tr}Q_X^3.$$
(6)
Motivated by semi-realistic models of string-derived effective supergravity with dynamical supersymmetry-breaking, we assume (in units where $`m_P=1/\sqrt{8\pi G}=1`$):
$$|\delta |^2|\xi |1.$$
(7)
In addition to $`D_X`$, we also have the $`SU(3)_c`$ D-term $`D_c^{(r)}`$, $`r=1,\mathrm{},8`$, which does not play a role in the following analysis. Let capital indices $`I,J`$, etc., refer collectively to fields $`a_i,b_i,\varphi _i`$. Then the scalar potential of the toy model is given by:
$`V`$ $`=`$ $`{\displaystyle \frac{g_X^2}{2}}D_X^2+{\displaystyle \frac{g_c^2}{2}}D_c^{(r)}D_{(r)}^c+e^K[\delta ^{I\overline{J}}(W_I+WK_I)(\overline{W}_{\overline{J}}+\overline{W}K_{\overline{J}})`$ (8)
$`+K^{a\overline{b}}(W_a+WK_a)(\overline{W}_{\overline{b}}+\overline{W}K_{\overline{b}})3W\overline{W}].`$
We parameterize the vacuum value of the moduli sector F-term as
$$K^{a\overline{b}}(W_a+WK_a)(\overline{W}_{\overline{b}}+\overline{W}K_{\overline{b}})=\alpha |\delta |^2,\alpha 1.$$
(9)
The requirement of a vanishing cosmological constant gives
$`<V>`$ $`=`$ $`{\displaystyle \frac{g_X^2}{2}}D_X^2+e^{<K>}|\delta |^2\left(v^2+\alpha 3\right)=0,`$
$`D_X`$ $`=`$ $`{\displaystyle \underset{i}{}}q_i|v_i|^2+\xi ,`$ (10)
where $`v_i=\varphi _i`$ and $`v=\sqrt{_i|v_i|^2}`$. Since $`V`$ is gauge neutral, $`V/a_i`$ and $`V/b_i`$ are $`SU(3)`$-charged and vanish when (4) holds. The minimization condition for $`\varphi ^i`$ gives
$$<V_i>\frac{V}{\varphi _i}=0=\overline{v}_i\left[g_X^2q_iD_X+|\delta |^2e^{<K>}(v^2+\alpha 2)\right].$$
(11)
This implies that $`v_i0`$ for only one value of $`q_iq`$, which, as we shall see, must be negative under our assumption (7). Note that (11) require $`D_X|\delta |^2<|W|>`$, so that supersymmetry breaking is dominated by the moduli sector under our assumptions. (11) and (10) together imply
$$v^2+\alpha 3=O(|\delta |^2),v^2+\alpha 21,$$
(12)
so that $`D_X>0`$. Next we consider the spectrum of $`\mathrm{\Phi }^i`$. The fermion mass matrix takes the form
$$(M_f^2)_j^i=v^i\overline{v}^{\overline{ȷ}}\left[2g_X^2q^2+e^{<K>}|\delta |^2(v^2+\alpha )\right],$$
(13)
and the scalar mass matrix (in Landau gauge) takes the form
$`M_s^2`$ $`=`$ $`\left(\begin{array}{cc}M^2& N^2\\ (N^{})^2& (M^{})^2\end{array}\right),`$
$`(M^2)_j^i`$ $`=`$ $`(M_f^2)_j^iv^i\overline{v}^{\overline{ȷ}}g_X^2q^2+\delta _j^i\left[g_X^2q_iD_X+e^{<K>}|\delta |^2\left(v^2+\alpha 2\right)\right]`$
$`(N^2)_j^{\overline{ı}}`$ $`=`$ $`\overline{v}^{\overline{ı}}\overline{v}^{\overline{ȷ}}\left[g_X^2q^2+e^{<K>}|\delta |^2(v^2+\alpha 1)\right].`$ (14)
In the absence of supersymmetry breaking, $`D_X=\delta =0`$, the superfield
$$\mathrm{\Pi }=\frac{1}{v}\underset{i}{}\overline{v}^i\mathrm{\Phi }^i$$
(15)
is eaten by the $`U(1)_X`$ gauge supermultiplet to form a massive vector multiplet. The orthogonal combinations
$$D_\alpha =\underset{i}{}c_\alpha ^i\mathrm{\Phi }^i,\underset{i}{}v^ic_\alpha ^i=0,\underset{i}{}\overline{c}_\alpha ^ic_\beta ^i=\delta _{\alpha \beta },$$
(16)
are the massless D-moduli. When $`\delta 0`$, the moduli $`M_a`$ mix with $`\mathrm{\Pi }`$; in other words the “Goldstone” chiral multiplets associated with supersymmetry breaking and with $`U(1)_X`$ breaking mix at order $`|\delta |^2/g_X^2q^2v^2`$. There is no mixing between the D-moduli and the M-moduli, so the D-moduli masses can be read directly from (13) and (14) by setting $`v_i=0`$; one obtains $`M_{Df}^2=0,`$ and, using the vacuum conditions (10) and (11), for the scalars $`d_\alpha =D_\alpha |`$:
$$m_\alpha ^2=(q_\alpha +q)g_X^2<D_X>=\frac{(q_\alpha +q)}{q}\left[m_{\stackrel{~}{G}}^2+O(|\delta |^4)\right],$$
(17)
where
$$m_{\stackrel{~}{G}}=e^{<K>/2}\delta $$
(18)
is the gravitino mass. From (17) it is clear that $`m_\alpha ^20`$ if and only if $`q=\mathrm{min}(q_1,\mathrm{},q_{N_\mathrm{\Phi }})`$. The $`d_\alpha `$ that are linear combinations of $`\varphi ^i`$ with $`q_i=q`$ remain massless, while the others acquire masses of order of the gravitino mass.
We now turn to a more specific model for supersymmetry breaking via gaugino condensation, as realized in the linear multiplet formulation for the dilaton . Our model is an approximation to the BGW model, in that we neglect the moduli-dependence of the Kähler metric for the gauge-charged matter fields that we consider. With this approximation, the scalar potential takes the form<sup>2</sup><sup>2</sup>2The full scalar potential for the BGW model is given in .
$`V`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_X^2D_X^2+{\displaystyle \underset{I}{}}\left|(W_I+K_IW)e^{K/2}+\beta uK_I\right|^2`$ (19)
$`+f(\mathrm{})\left|u(1+b_a)4\mathrm{}We^{K/2}\right|^2{\displaystyle \frac{3}{16}}\left|ub_a4We^{K/2}\right|^2,`$
where $`u=e^{K/2}\stackrel{~}{u}(\mathrm{},t)`$ is the gaugino condensate which is determined by the equations of motion as a function of the dilaton $`\mathrm{}`$ and the T-moduli $`t`$ (treated here as constants<sup>3</sup><sup>3</sup>3When the dilaton is dynamical, $`g_Xg_X(\mathrm{})`$ and $`g_X^2\xi 2\mathrm{}\xi `$. (When nonperturbative string effects are neglected, $`g_X^2(\mathrm{})=2\mathrm{}`$ at the string scale.)), $`W=W(\varphi )`$, $`b_a`$ is the $`\beta `$-function coefficient for the condensing gauge group, and the function $`f(\mathrm{})`$ depends on the Kähler potential for the dilaton. The terms in (19) are, in a one-to-one correspondence, the counterparts in this model of the terms in (8). The vacuum conditions are, assuming $`<W(\varphi )>=<W_I>=0`$,
$`<V>`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_X^2D_X^2+f(\mathrm{})|u|^2(1+b_a)^2{\displaystyle \frac{3}{16}}|u|^2b_a^2+\beta ^2|u|^2v^2`$
$``$ $`{\displaystyle \frac{1}{2}}g_X^2D_X^2+\widehat{V}=0,`$
$`<V_i>`$ $`=`$ $`\overline{v}^{\overline{ı}}\left(q_ig_X^2D_X+\widehat{V}+\beta ^2<|u|^2>\right).`$ (20)
In the spirit of the previous section, we assume (7) with $`<|u|>|\delta |`$, so that $`D_X|\delta |^2,\widehat{V}|\delta |^4`$. As in the previous example, one chiral supermultiplet with $`q_i=q`$, the lowest $`U(1)_X`$ charge, forms a massive gauge supermultiplet with the $`U(1)_X`$ gauge superfield, while the remaining chiral superfields have massless fermions and scalar masses now given by
$$m_\alpha ^2=(q_\alpha +q)g_X^2<D_X>=\frac{(q_\alpha +q)}{q}\left[\left(\frac{4\beta }{b_a}\right)^2m_{\stackrel{~}{G}}^2+O(|\delta |^4)\right].$$
(21)
When the GS term that cancels the T-duality anomaly is included, the parameter $`\beta `$ is given by
$$\beta =(p_\alpha b_a)/4,$$
(22)
where $`p_\alpha `$ measures the coupling of the fields $`\varphi ^i`$ to the GS term. Here the situation is the same as for squarks and sleptons; if $`p_\alpha =0`$, $`m_\alpha m_{\stackrel{~}{G}}`$, while if the D-moduli couple to the GS term with the same strength $`b`$ as the T-moduli, $`m_\alpha m_t/2|b/b_a1|m_{\stackrel{~}{G}}`$. In the BGW model with $`b=b_{E_8}10b_a`$ we get $`m_\alpha 10m_{\stackrel{~}{G}}`$. However, in the presence of Wilson lines that break the gauge group to a phenomenologically viable one, in general $`b<b_{E_8}`$. In particular, in the FIQS model discussed below, with an $`SO(10)`$ condensing group, $`b=b_a`$, which implies that the moduli masses are much smaller than the gravitino mass, so the FIQS model is not viable in the context of the BGW supersymmetry breaking scenario.
As mentioned above, in more realistic models the D-moduli are charged under additional, nonanomalous gauge groups $`U(1)_a`$. Assuming affine level one, so that the gauge couplings are all equal at the string scale:<sup>4</sup><sup>4</sup>4Strictly speaking we should integrate out the heavy modes at the scale $`v^i.1`$ and run the couplings their values to the condensation scale $`<u^{\frac{1}{3}}>10^4`$; we neglect such renormalization effects here. $`g_a=g_Xg`$, the potential (8) or (19) takes the form
$$V=\frac{g^2}{2}\underset{a}{}D_a^2+\widehat{V}.$$
(23)
where we have set to zero D-terms corresponding to nonabelian gauge groups such as $`SU(3)`$ in our toy model. The minimization conditions take the form
$$0=<V_i>=g^2\underset{a}{}<D_a>q_i^a\overline{v}_i+<\widehat{V}_i>.$$
(24)
In the models we are considering, we may write
$$<\widehat{V}_i>=\overline{v}_if,$$
(25)
where $`f`$ is some function of the $`v^i`$ and the moduli $`vevs`$’s. Then for any $`i`$ such that $`v^i0`$, (24) implies:
$$0=g^2\underset{a}{}D_aq_i^a+f=g^2\underset{a}{}q_i^a\underset{j}{}|v^j|^2+g^2q_i^X\xi +f.$$
(26)
We are interested in models with F-flat and D-flat directions, i.e. in which the set of equations
$$D_a=0$$
(27)
has a solution along some F-flat direction. For example in the FIQS model, with $`a=1,\mathrm{}8,`$ (a continuous degeneracy of) solutions exist with nonvanishing $`vev`$’s for a set of 27 complex scalar fields<sup>5</sup><sup>5</sup>5Our notation differs slightly from that of .
$$\mathrm{\Phi }^iS_\alpha ^i,Y_A^i,\alpha =1,\mathrm{}5,A=1,\mathrm{}4,i=1,\mathrm{}3,$$
(28)
which are charged only under the $`U(1)_a`$, with the charges $`q_a`$ independent of the index $`i`$ and
$$q_X^S=8,q_X^Y=4,q_a^{Y_1}=q_a^{Y_2},q_a^{Y_3}=q_a^{Y_4}.$$
(29)
The set of equations (27) have solutions that break 6 of the 8 $`U(1)`$’s, including $`U(1)_X`$, leaving unbroken the weak hypercharge of the Standard Model and one additional $`U(1)`$; in the effective low energy theory at scales $`\mu \xi `$ there are no supermultiplets that carry both this latter $`U(1)`$ charge and Standard Model gauge charges . Since 6 $`U(1)`$’s are broken, 6 of the supermultiplets in (28) are eaten by massive vector multiplets, and we are left with 21 D-moduli supermultiplets, instead of the 26 we would have in the absence of additional $`U(1)_a`$ charges for these fields.
Now consider the effect of the supersymmetry breaking term $`f`$ in (24). The solution to (27) for $`aX`$ requires that at least one field $`Y`$ have a nonvanishing vacuum value. Therefore $`D_{aX}=0`$ is not a solution to (24) since the previous analysis without the additional $`U(1)_a`$ requires in this case that only $`<S>0`$ when $`f0`$. Hence we are led to solve<sup>6</sup><sup>6</sup>6Again we are oversimplifying; once supersymmetry is broken there is no reason to assume that the F-terms involving D-moduli couplings in the superpotential remain zero. the set (26) of coupled equations for the $`|v^i|^2`$. We have analyzed these equations using the math package Maple and the $`U(1)_a`$ charge assignments of the FIQS model, and find that the minimum corresponds to
$`Y_3^i=Y_4^i=0,{\displaystyle \underset{i}{}}|S_\alpha ^i|^2=f_\alpha (\xi ,g^2,f),`$
$`{\displaystyle \underset{i}{}}\left(|Y_1^i|^2+|Y_2^i|^2\right)=f_Y(\xi ,g^2,f),`$ (30)
with an additional constraint of the form $`f=f(g^2,\xi )`$ to assure vanishing of the cosmological constant. Now $`Y_1^i`$ and $`Y_2^i`$ correspond to 12 real fields constrained by one equation to give 11 moduli, and each of the 5 choices of $`\alpha `$ in $`S_\alpha ^i`$ correspond to 6 real fields subject to one constraint giving 5 moduli each. In this model there are 6 $`U(1)`$’s that get broken, so 6 moduli are eaten, leaving a total of
$$5\times 5+116=30$$
D-moduli. Note that while there are fewer “light” D-moduli ($`mm_{\stackrel{~}{G}}`$) than in the toy model with the same $`U(1)_X`$ charges but no additional $`U(1)`$’s, there are actually more massless scalars (30 instead of 28) after supersymmetry breaking.
We remark that the first condition in (30) (which corresponds to $`Y_3^{1i}=Y_3^{2i}=0`$ in the notation of ) has as a consequence that all of the down-type quarks are massless at tree level, and their masses must be generated by radiative corrections. Leaving aside the moduli problem alluded to above, this could be a phenomenological improvement of the model as compared with the solution of (27) in the absence of supersymmetry breaking. In that case there are both up- and down-type quark masses at tree level, but (unless unmotivated mixing is introduced in the Kähler potential) the CKM matrix is unrealistic: the heaviest up quark is not in the same $`SU(2)_L`$ gauge multiplet as the heaviest down quark. It has recently been shown that all down-type quark masses can be generated entirely from radiative corrections, subject to certain conditions on the high energy theory and supersymmetry breaking scenario . Whether or not viable quark masses can be gotten by this mechanism in the FIQS model is under investigation.
In the generic gaugino condensation model of , supersymmetry breaking arises from the Veneziano-Yankielowicz part of the superfield Lagrangian:
$$_{\text{VY}}=d^4\theta \frac{E}{8R}U\left[b^{}\mathrm{ln}(e^{K/2}U)+\underset{\alpha }{}b_\alpha \mathrm{ln}\mathrm{\Pi }_\alpha \right]+\text{h.c.}.$$
(31)
The values of $`b^{}`$ and $`b_\alpha `$ are determined by anomaly matching and are related to the $`\beta `$-function coefficient by
$$b_a=b^{}+\underset{\alpha }{}b_\alpha .$$
(32)
The (weight two) chiral field $`U`$ is the gaugino condensate superfield: $`U|=u`$, while the (weight zero) chiral fields $`\mathrm{\Pi }`$ are matter condensates. Condensation occurs provided there is also a superpotential for the matter condensates:
$$W(\mathrm{\Pi },T)=\underset{\alpha }{}c_\alpha (T)\mathrm{\Pi }_\alpha ,$$
(33)
where the moduli-dependence of the coefficient assures modular invariance. In it was assumed that $`\mathrm{\Pi }`$ is a composite operator containing fields charged only under the condensing gauge group. However in many models – such as the FIQS model with a hidden sector $`SO(10)`$ and matter in 16’s – there is no operator that can be constructed from these fields alone that is invariant under the $`U(1)_a`$, and the coefficient $`c_\alpha `$ must depend on the $`\mathrm{\Phi }^i`$. It is possible that these additional couplings of the D-moduli are sufficient to lift the remaining degeneracy of the vacuum – and they may also give $`O(m_{\stackrel{~}{G}})`$ masses to the D-moduli fermions. This is because (25) would no longer hold, so that we do not get the single condition (26), but rather several independent conditions from (24). An analysis of this case requires a careful treatment of renormalization effects. However, it appears likely that any masses generated by these additional couplings will be governed by the supersymmetry breaking scale.
A large number of light scalar fields is problematic for cosmology in realistic models. The D-moduli have no gauge couplings in the effective low energy theory since they are charged only under the $`U(1)`$’s that are broken near the string scale. Therefore, unless they have unsuppressed superpotential couplings to relatively light particles, they are subject to the same constraints as, e.g., the T-moduli . The problem is somewhat alleviated if they couple to the Green-Schwarz term with $`p_\alpha =b`$ in (22) and $`bb_a`$. In this case, like the T-moduli in the BGW model, their masses can exceed the gravitino mass by an order of magnitude. In $`Z_3`$ and $`Z_7`$ compactifications, with no T-moduli-dependent string threshold corrections , $`bb_a`$ where the inequality is saturated if there are no twisted sector fields that are charged under the gauge group $`𝒢_a`$. This is the case, in particular for the FIQS model that we used above to illustrate the case of many $`U(1)`$’s. As we noted previously, this is not a viable option in the BGW context, since it gives unacceptably small moduli masses. In addition, models with $`bb_a`$ alleviate problems associated with dilaton cosmology.
It is plausible that the fermions can be sufficiently diluted by inflation to be harmless. The decay and annihilation of the particles of the Minimal Supersymmetric Standard Model suppresses the energy density of a decoupled massless fermion relative to a neutrino only by a factor of about 20, so 20 such fermions would contribute the equivalent of one neutrino species to the energy density during Nucleosynthesis if there is no other suppression mechanism. Fermions with order $`TeV`$ masses would vastly overclose the universe unless they are inflated away or are sufficiently short lived. With only gravitational strength couplings, dimensional analysis suggests decay rates $`\mathrm{\Gamma }_\alpha m_\alpha ^3/m_P^2`$, which were shown to be marginally acceptable for the T-moduli and dilaton fermions. Any remnant of the broken $`U(1)`$’s in the tree-level couplings of the D-moduli would tend to suppress the decay rate. For example $`U(1)`$ invariant couplings would give rates $`\mathrm{\Gamma }_\alpha m_\alpha ^5/m_P^4`$ that are unacceptably small without sufficient dilution. An estimate of the D-moduli fermion lifetimes requires a detailed understanding of the effective theory below the scale where the $`U(1)_a`$ are broken.
The point that we wish to emphasize here is that there are, generically, many more light moduli than have been previously considered, which may imply much stronger constraints on their masses and/or couplings. A full analysis of the D-moduli spectrum in the context of the BGW model for supersymmetry breaking, including dynamical T-moduli and dilaton, will be given elsewhere. |
warning/0001/nlin0001027.html | ar5iv | text | # Mean- Field Approximation and a Small Parameter in Turbulence Theory
## 1 Introduction
The role of the mean-field theories and gaussian limits as starting points for understanding of such important physical phenomena as superconductivity, superfluidity, critical point, naming just a few, can hardly be overestimated. These theories, usually based on remarkable physical intuition and insight, provided mathematical and intellectual basis for investigation of much more difficult regimes in terms of deviations from the mean field solutions. The most recent example is a theory of anomalous scaling in a model of a passive scalar, advected by a random velocity field, which was developed as an expansion in powers of small parameters characterizing deviations from the two gaussian limits -. In a typical non-linear system, a gaussian limit corresponds to a weak coupling asymptotics and, as a consequence, to a “normal”, non-anomalous, scaling, which can often be obtained from a “bare’ or linearized problem. A good example of this behaviour is a fluid in thermodynamic equilibrium.
The large- Reynolds- number three-dimensional ($`3d`$) strong turbulence is characterized by an $`O(1)`$ energy flux $`=\nu \overline{(_iv_j)^2}`$, which in many flows is $`O(v_{rms}^3/L)`$, where $`L`$ is an integral scale of turbulence. In the inertial range, where $`k_d>>k\mathrm{}`$ or $`r/L<<1`$, the observed energy spectrum $`E(k)`$ is close to the one, proposed by Kolmogorov and the probability density $`P(\mathrm{\Delta }u)`$ with $`\mathrm{\Delta }u=u(x+r)u(x)`$, is far from the gaussian. Moreover, the experiments revealed that the moments of velocity difference $`S_{n,0}=\overline{(\mathrm{\Delta }u)^n}r^{\xi _n}`$ with the exponents $`\xi _n`$ given by very strange (“anomalous”) numbers, which cannot be obtained on dimensional grounds. This anomalous scaling and the very existence of the energy flux, resulting in the non-zero value of the third-order moment $`S_{3,0}=\overline{(\mathrm{\Delta }u)^3}O(r)`$, where $`𝐮𝐫=ur`$, implies a strongly non-gaussian process and an obvious lack of the mean-field limit.
The situation may not be so grim, however: all odd-order moments of transverse velocity differences in both $`2d`$ and $`3d`$ flows $`S_{0,2n+1}=\overline{(v(x+r)v(x))^{2n+1}}=0`$ with $`𝐯𝐫=\mathrm{𝟎}`$. This fact tells us that these components of velocity differences do not participate in the inter-scale energy transfer and there is no a’priori any reason for them not to obey gaussian statistics in some limiting cases. This is indeed true in two-dimensional turbulence in the inverse cascade range where $`l_f<<r<<L`$ and $`l_f`$ is a forcing scale.
Numerical and physical experiments on the external- force- driven two- dimensional turbulence showed that the moments of transverse velocity differences and even-order moments of the longitudinal ones are very close to the gaussian values and are characterized by the Kolmogorov scaling exponents $`S_{0,2n}S_{2n,0}r^{\xi _{2n}}`$ with $`\xi _{2n}=\frac{2n}{3}`$ -. The odd-order moments of longitudinal velocity differences are positive in a $`2d`$ flow, while they all are negative in $`3d`$. This observation tells that the most important distinction between two and three-dimensional turbulence in in dynamic role of the dissipation contributions: they are irrelevant in the two-dimensional inverse cascade range and are crucial (see below) for the small-scale dynamics of a three-dimensional flow where the forcing terms can be neglected (see below). Thus, we can assume that it is the dissipation fluctuations that are responsible for both strong deviations from the gaussian statistics and anomalous scaling in three dimensional flows. This assumption is consistent with the observation of the close-to-gaussian probability density of velocity differences in $`3d`$ turbulence at the scales $`rL`$ where the integral scale $`L`$ is defined as the one at which: $`S_{3,0}(L)=0`$ and the intermittent dissipation fluctuations disappear . It is clear that this range ($`rL`$) is not characterized by a well-defined scaling exponents.
As will become clear below, the forcing contribution to the equation for the probability density involves factor $`\mu 1cos(k_fr)`$. This means that at the scales $`r>>1/k_f`$ the parameter $`\mu 1`$ while $`\mu (k_fr)^2`$ when $`k_fr<<1`$. Thus, the forcing term must be important in the inertial range of two-dimensional turbulence ($`r>>l_f`$) and is irrelevant in the $`3d`$-inertial range with the positive energy flux. This change happens at some space dimensionality $`d=d_c`$ at which the energy flux changes sign. First calculation of $`d_c2.05`$ was conducted by Frisch et. al. within the framework of a simple closure model. The more physically transparent calculation can be performed for the Navier-Stokes equations driven by a random force having an algebraically decaying spectrum in the inertial range , where a one- loop small-scale- elimination procedure gives a correction to the bare viscosity:
$$\delta \nu =\nu \frac{d^2dϵ}{2d(d+2)}Re^2$$
where $`Re`$ is a properly defined Reynolds number corresponding to the eliminated small-scale velocity fluctuations and $`ϵ`$ is a parameter characterizing the forcing function. In case of Kolmogorov turbulence $`ϵ4`$. This relation shows that the role, the small-scales play in turbulence dynamics, depends on the space dimensionality $`d`$: the correction to viscosity is positive when $`d>d_c(ϵ)`$ where it changes sign. Physically, this means that the small-scale velocity fluctuations take energy from the large- scales motions (direct energy cascade from large-to-small structures) at $`d>d_c`$, while at $`d<d_c`$ they excite the large scale motions (spend their energy) giving rise to the inverse energy cascade. For $`ϵ=4`$ the critical dimensionality $`d_c2.56`$. The correct value of $`d_c`$ is not too important: what is crucial for the theory presented below is that the critical dimensionality, at which the flux changes its sign, exists. It will be shown below that $`dd_c0`$ is a small parameter of the theory enabling one to calculate an expression for the dissipation anomaly in the form, resembling the Kolmogorov refined similarity hypothesis.
Since in $`2d`$ the moments $`S_{0,2n}`$ show Kolmogorov scaling, the gaussian statistics of transverse velocity differences cannot correspond to the weak coupling limit. This problem was considered in the Ref. where, following Polyakov , the equation for the generating function for the problem of the Navier-Stokes turbulence was first introduced . An unusual symmetry of this equation enabled one to show that the solution was consistent with both Kolmogorov scaling and gaussian statistics. In this work a more detailed theory of two-dimensional turbulence is presented and the generalization to three-dimensional flows is considered. The main result of the paper is a model demonstrating how the deviations from the “normal scaling” and gaussian statistics appear in $`3d`$ when the strength of the dissipation term $`\beta >0`$ and the “scale” parameter $`ϵ=1r/L`$ deviate from zero.
This paper is organized as follows. In Section 2 the equation for the generating function, derived in is introduced. Some new exact relations between velocity structure functions, following from this equation, are derived in Section 3. The connection between scaling exponents of the moments of velocity differences and their amplitudes is established in Section 4. The mean-field derivation of the pressure term is given in Section 5 which is used to obtain a gaussian pdf in the two-dimensional flow in Section 6. In Sections 7 the small parameter of the theory is identified and used for derivation of the expression for the dissipation contributions to the eqaution for the pdf. Section 8 is devoted to solution of the equation and demonstration how anomalous scaling and deviations from gaussian statistics emergr from the theory. Conclusive remarks are presented in Section 9.
## 2 Equation for the Generating Function
The equations of motion are (density $`\rho 1`$):
$$_tv_i+v_j_jv_i=_ip+\nu ^2v_i+f_i;_iv_i=0$$
(1)
where $`𝐟`$ is a forcing function responsible for the large-scale kinetic energy production and in a statistically steady state the mean pumping rate $`P=\overline{𝐟𝐯}`$. In what follows we will be mainly interested in the probability density function of two-point velocity difference $`𝐔=𝐮(𝐱^{})𝐮(𝐱)\mathrm{\Delta }𝐮`$. The generating function is: $`Z=<exp(\lambda 𝐔)>`$. The equation for the generating function of velocity differences corresponding to (1) is:
$$\frac{Z}{t}+\frac{^2Z}{\lambda _\mu r_\mu }=I_f+I_p+D$$
(2)
with
$$I_f=<\lambda 𝚫𝐟e^{\lambda \mathrm{\Delta }𝐮}>$$
(3)
$$I_p=\lambda <e^{\lambda \mathrm{\Delta }𝐮}\mathrm{\Delta }(p)>\lambda <e^{\lambda 𝐯}(_\mathrm{𝟐}p(x_2)_\mathrm{𝟏}p(x_1))>$$
(4)
and
$$D=\nu \lambda <(_\mathrm{𝟐}^\mathrm{𝟐}𝐯(𝐱_\mathrm{𝟐})_\mathrm{𝟏}^\mathrm{𝟐}𝐯(𝐱_\mathrm{𝟏}))e^{\lambda 𝐔}>$$
(5)
The most interesting and surprising feature of (2) is the fact that, unlike in the problem of Burgers turbulence , the advective contributions are represented there in a closed form. This means that the theory, developed below, is free from the troubles related to Galileo invariance, haunting all schemes, based on renormalized perturbation expansions in powers of Reynolds number. To completely close the problem the expressions for $`I_p`$ and $`D`$ are needed. The equations (2)-(3) formulate the turbulence theory in terms of “only” two unknowns $`I_p`$ and $`D`$. The Kolmogorov refined similarity hypothesis stating that $`(\mathrm{\Delta }u)^3=\varphi _rr`$ where $`\varphi `$ is a scale-independent random process and $`_r`$ is a dissipation rate averaged over a ball of radius $`r`$ around point $`x`$, can be a promising starting point to a closure for the dissipation term $`D`$. This will be done below. The pressure term in (2)-(3) is also of a very specific and rather limited nature: all we have to know is the correlation functions $`<U_iU_j\mathrm{}U_m\mathrm{\Delta }p>`$. Thus, the definite targets needed for derivation of the closed equation for $`Z`$-functions are well-defined.
The generating function can depend only on three variables: $`\eta _1=r;\eta _2=\frac{\lambda 𝐫}{r}\lambda cos(\theta );\eta _3=\sqrt{\lambda ^2\eta _2^2};`$ and
$$Z_t+[_{\eta _1}_{\eta _2}+\frac{d1}{r}_{\eta _2}+\frac{\eta _3}{r}_{\eta _2}_{\eta _3}+\frac{(2d)\eta _2}{r\eta _3}_{\eta _3}\frac{\eta _2}{r}_{\eta _3}^2]Z=I_f+I_p+D$$
(6)
where, to simplify notation we set $`_{i,\alpha }\frac{}{x.\alpha }`$ and $`v(i)v(𝐱_𝐢)`$. Below we will often use $`_{\eta _i}_i`$. The functions $`I_p`$, $`I_f`$ and $`D`$ are easily extracted from the above definitions. Let us denote $`\mathrm{\Delta }uU`$ and $`\mathrm{\Delta }vV`$.
In the new variables the generating function can be represented as:
$$Z=<e^{\eta _2\mathrm{\Delta }u+\eta _3\mathrm{\Delta }v}><e^{\eta _2U+\eta _3V}>$$
with the mean dissipation rate $``$ defined by: $`\overline{\nu (_{x_i}u)^2}=\frac{1}{d}`$. Any correlation function is thus:
$$S_{n,m}<U^nV^m>=_2^n_3^mZ(\eta _2=\eta _3=0,r)$$
## 3 Relations between moments of velocity difference
Let us discuss some direct consequences of equations (1)-(6). The Navier-Stokes equations are invariant under transformation: $`vv`$ and $`yy`$. That is why: $`<(_yp(0)_yp(r))(\mathrm{\Delta }v)^m>0`$ if $`m=2n+1`$ with $`n>1`$ and is equal to zero if $`m=2n`$. It is also clear from the symmetry that $`<_xpU^{2n}>=0`$. It follows from the Navier-Stokes equations that $`<(^2U)U^{2n}>=<(^2V)V^{2n}>=0`$.
Multiplying (6) by $`\eta _3`$, applying $`_2^{2n1}`$ to the resulting equation gives as $`\eta _2=\eta _30`$
$$\frac{S_{2n,0}}{r}+\frac{d1}{r}S_{2n,0}\frac{(d1)(2n1)}{r}S_{2n2,2}=P(1cos(r/l_f))\frac{2(2n1)(2n2)}{d}S_{2n3,0}$$
(7)
The right side of (7) is $`O(r^2)`$ in a three-dimensional flow (3d) where $`r/l_f<<1`$ and and can be neglected. It is however is $`O(1)`$ in two dimensional turbulence in the inverse cascade range where $`r/l_f>>1`$. The dissipation terms do not contribute to this relation. For $`n=1`$ (7) gives a well-known incompressibility relation , :
$$\frac{S_{2,0}}{r}+\frac{d1}{r}S_{2,0}=\frac{d1}{r}S_{0,2}$$
(8)
Multiplying (6) by $`\eta _3`$ with subsequent $`_3_2^2`$ leads after setting $`\eta _2=\eta _3=0`$ as $`\nu 0`$ to:
$$\frac{S_{3,0}}{r}+\frac{d1}{r}S_{3,0}2\frac{d1}{r}S_{1,2}=(1)^d\frac{4}{d}P$$
(9)
where $`d=2;3`$ . Applying $`_3^3\eta _3`$ to (6) gives, as $`\nu 0`$:
$$\frac{S_{1,2}}{r}+\frac{d+1}{r}S_{1,2}=(1)^d\frac{4}{d}P$$
(10)
Substituting this into (9) yields a well-known Kolmogorov relation:
$$S_{3,0}\overline{(\mathrm{\Delta }u)^3}=(1)^d\frac{12}{d(d+2)}Pr$$
(11)
For $`2n=4`$ the relation (7) reads:
$$\frac{S_{4,0}}{r}+\frac{d1}{r}S_{4,0}=\frac{3(d1)}{r}S_{2,2}$$
(12)
In two-dimensional turbulence $`(d=2)`$ in the inverse cascade range one can neglect the dissipation contribution $`D`$ (see below) and derive:
$$\frac{S_{1,2n}}{r}+\frac{1+2n}{r}S_{1,2n}=n(2n1)PS_{0,2n2}2n<𝒫_{yv}(\mathrm{\Delta }v)^{2n1}>$$
(13)
where $`𝒫_{yv}_yv(x+r)_yv(x)`$. Another interesting relation, valid in $`2d`$ is obtained from (6) readily:
$$\frac{S_{2,2n}}{r}+\frac{1+2n}{r}S_{2,2n}=\frac{S_{0,2n+2}}{r}2n<𝒫_{yv}\mathrm{\Delta }u(\mathrm{\Delta }v)^{2n1}>+n(2n1)PS_{1,2n2}$$
(14)
In the direct cascade range, where the forcing contribution is $`O(r^2)0`$, the relation (14) for an arbitrary dimensionality $`d`$ reads:
$$\frac{S_{2,2n}}{r}+\frac{d1+2n}{r}S_{2,2n}=\frac{2n+d1}{2n+1}\frac{S_{0,2n+2}}{r}2n<𝒫_{yv}\mathrm{\Delta }u(\mathrm{\Delta }v)^{2n1}>$$
(15)
Now, let us multiply (6) by $`\eta _3`$, differentiate once over $`\eta _2`$ and three times over $`\eta _3`$. This gives:
$$\frac{S_{2,2}}{r}+\frac{d+1}{r}S_{2,2}=\frac{d+1}{3r}S_{0,4}2\overline{𝒫_{yv}\mathrm{\Delta }u\mathrm{\Delta }v}$$
(16)
This relation is correct since $`\nu \overline{^2v\mathrm{\Delta }u\mathrm{\Delta }v}=\nu \overline{^2u(\mathrm{\Delta }v)^2}=0`$.
2d simulations of Bofetta, Celani and Vergassola To achieve a true steady state these authors conducted a series of very accurate simulations of the problem (1) with the large-scale dissipation term $`D_L=\alpha 𝐯`$ in the right side (6). The moments of transverse velocity differences, reported in this paper ($`n2`$), were very close to their gaussian values. It is clear that this term introduces $`\alpha \lambda _\mu \frac{Z}{\lambda _\mu }`$ into the right side of (6) which is small in the inertial range where the non-linearity is large. One has to be careful though with the dangerous interval $`\mathrm{\Delta }v0`$ where the linear term is not small. We expect the negative- order structure functions with $`1<n<0`$ strongly depend on the functional shape of the otherwise irrelevant large-scale dissipation term. The same can be predicted for various conditional expectation values of dynamical variables, like pressure gradients and dissipation terms, for the fixed values of $`\mathrm{\Delta }v\mathrm{\Delta }u`$: near the origin where $`\mathrm{\Delta }v`$ and $`\mathrm{\Delta }u`$ are very small, the artificially introduced linear contributions to the Navier-Stokes equations dominates, producing large and non-universal deviations from the universal functions characterizing inertial range. For example, with addition of the linear dissipation, the relation (14) reads:
$$\frac{S_{2,2n}}{r}+\frac{1+2n}{r}S_{2,2n}=\frac{S_{0,2n+2}}{r}(2n+1)\alpha S_{1,2n}2n<𝒫_{yv}\mathrm{\Delta }u(\mathrm{\Delta }v)^{2n1}>+n(2n1)PS_{1,2n2}$$
(17)
modifying the balance (pressure contribution) in the range of small product $`\mathrm{\Delta }u\mathrm{\Delta }v`$ or $`r/L1`$. In the interval where $`\mathrm{\Delta }u\mathrm{\Delta }v`$ is not small ($`r0`$), the linear terms are small and can be neglected.
The results obtained in this section can also be derived with
$$Z=e^{\sqrt{d1}\eta _3V+\eta _2U}$$
(18)
with the properly defined moments of $`V`$ and $`U`$.
## 4 Asymptotic values of exponents in three-dimensional flows
In case of intermittent turbulence $`S_{m,n}=A_{m,n}r^{\xi _{m,n}}`$. We can see that in the inertial range of a three-dimensional flow ($`r/l_f<<1`$) the right side of (7) is negligible and, as a result, $`\xi _{2n,0}=\xi _{2n2,2}\xi _{2n}<\frac{2n}{3}`$. Substituting this into (7) gives immediately:
$$\xi _{2n}=(d1)[(2n1)\frac{A_{2n2,2}}{A_{2n,0}}1]$$
(19)
Let us introduce the probability density functions:
$$S_{2n,0}=\overline{U^{2n}}=P(U)U^{2n}𝑑U$$
(20)
and
$$S_{2n2,2}=\overline{U^{2n2}V^2}=P(U)U^{2n2}q(V|U)𝑑U𝑑V$$
(21)
where $`q(V|U)`$ is conditional pdf of $`V`$ for fixed value of $`U`$. It is clear that $`q(V|U)=q(V,U)`$, so that all odd-order moments of $`V`$ are equal to zero.. This expression can be rewritten as:
$$S_{2n2,2}=\overline{U^{2n2}V^2}=P(U)U^{2n2}Q_2(U)𝑑U$$
(22)
where $`Q_2`$ is a conditional expectation value of $`V^2`$ for a fixed value of $`U`$:
$$Q_2(U)=V^2q(V|U)𝑑V$$
(23)
Comparing the above expressions we observe that the amplitudes $`A_{2n,0}A_{2n2,2}`$ only if $`Q_2U^2`$. This seems rather improbable in the limit $`U\mathrm{}`$ $`(n\mathrm{})`$. As a result, the linear regime $`\xi _{2n}n`$ is equally improbable.
Saturation of exponents $`\xi _{2n}\xi _{\mathrm{}}=const`$ as $`n\mathrm{}`$ is possible on a rather wide class of probability densities. For example,
$$P(U)A\frac{1}{U}\frac{}{U}P(U)Q_2(U)$$
(24)
where $`A>0`$ is a constant. Then,
$$S_{2n,0}=(2n1)AP(U)Q_2(U)U^{2n2}𝑑U=(2n1)AS_{2n2,2}$$
(25)
Substituting this result into the expression for $`\xi _{2n}`$ gives $`\xi _{2n}(d1)(A1)=const`$. The relation (24) defines the large-$`U`$ asymptotics of the pdf $`P(U)`$ in terms of $`Q_2(U)`$:
$$P(U)\frac{1}{Q_2(U)}e^{A^U\frac{udu}{Q_2(u)}}$$
(26)
If $`Q_2(U)U^\beta `$, then, assuming the existence of all moments, it follows from (25) that $`\beta <2`$. The “lognormal “ pdf $`P(U)`$ corresponds to $`Q_2(U)U^2/2log(U)`$. The expression (24) also gives in the limit of large $`n`$:
$$S_{2n+1,0}=\overline{U^{2n+1}}=2nA\overline{V^2U^{2n1}}=2nAS_{2n1,2}$$
(27)
## 5 Pressure contributions
Due to the symmetries of the Navier-Stokes equations, neither pressure nor dissipation terms contributed to the expressions (7) -(12). To proceed further we have to evaluate $`I_p`$ and $`D`$.
First of all we see from (12) that $`\xi _{4,0}=\xi _{2,2}`$. Let us assume that in the inertial range $`S_{4,0}=A_{4,0}r^{\xi _{4,0}}`$, $`S_{0,4}=A_{0,4}r^{\xi _{0,4}}`$ and $`S_{2,2}=A_{2,2}r^{\xi _{2,2}}`$. Then, it is clear from (12) and (16) that neglecting the pressure contribution to (16) gives $`\xi _{4,0}=\xi _{0,4}`$.
$`𝐝=\mathrm{𝟐}`$. It will be shown below that in 2d the even-order moments of velocity differences are very close to the gaussian ones and all exponents are close to the K41 values $`\xi _n=n/3`$. Then, $`A_{4,0}=3A_{2,0}^2`$ and $`A_{0,4}=3A_{0,2}^2`$. It follows from (12) that $`A_{2,2}=\frac{7}{9}A_{4,0}`$. Taking into account that when $`d=2`$ the amplitudes $`\frac{5}{3}A_{2,0}=A_{0,2}`$, we conclude that without the pressure contribution the equations ( 12) and (16) are incompatible.
Following we introduce a conditional expectation value of the pressure gradient difference for a fixed value of $`\mathrm{\Delta }u`$, $`\mathrm{\Delta }v`$ and $`r`$:
$$<_yp(x+r)_yp(x)|\mathrm{\Delta }u,\mathrm{\Delta }v,r>\underset{m,n}{}\kappa _{m,n}(r)(\mathrm{\Delta }u)^m(\mathrm{\Delta }v)^n$$
(28)
where the functions $`\kappa _{m,n}(r)`$ ensure proper dimensionality of the corresponding correlation functions. The above expression explicitly assumes existence of an expansion of conditional expectation value (28). In general, this may not be true due to various singularities such as the ones arising in the dissipation contributions (see below). Since the pressure term involves only one spacial derivative, the ultra-violet singularity cannot appear. The infra-red singularity is not there at least in 2d where the integral scale is time-dependent (see below). Keeping only the first two terms of the expansion (28), produces a model for the pressure contributions:
$$<_yp(x+r)_yp(x)|\mathrm{\Delta }y\mathrm{\Delta }v>h\frac{\mathrm{\Delta }u\mathrm{\Delta }v}{r}b\frac{\mathrm{\Delta }v}{(Pr)^{\frac{2}{3}}}$$
(29)
Since in an incompressible and homogeneous flow $`\overline{\mathrm{\Delta }u𝒫_{xu}}=\overline{\mathrm{\Delta }v𝒫_{yv}}=0`$, the coefficients $`h`$ and $`b`$ are related as:
$$hS_{1,2}=bS_{0,2}(Pr)^{\frac{1}{3}}$$
(30)
Limiting the expansion of a conditional expectation value by the first terms resembles Landau’s theory of critical phenomena, well describing experimental data in a certain range of parameters variation. We will show below that in case of turbulence this approximation gives the results which are in agreement with the data. This may be a consequence of the fact that $`A_{0,4}A_{4,0}A_{2,2}=O(1)`$.
With $`\xi _n=n/3`$ it follows from (7), (13)-(15) and (29) that when $`d=2`$ and $`n\mathrm{}`$:
$$(\frac{2n(43h)}{3}+1)\frac{S_{1,2n}}{r}=n(2n1)PS_{0,2n2}+b\frac{2nS_{0,2n}}{(Pr)^{\frac{2}{3}}}$$
and
$$(\frac{2n(43h)}{3}+1)\frac{S_{2,2n}}{r}=\frac{S_{0,2n+2}}{r}+n(2n1)PS_{1,2n2}+b\frac{2nS_{1,2n}}{(Pr)^{\frac{2}{3}}}$$
In the limit $`n\mathrm{}`$ assuming that $`S_{1,2n}nS_{1,2n2}`$ one derives readily:
$$PrnS_{0,2n2}\frac{S_{2,2n+2}}{Prn}\frac{1}{Prn^2}S_{0,2n+4}$$
(31)
which is consistent with the gaussian pdf $`P(\mathrm{\Delta }v)`$ as $`\mathrm{\Delta }v\mathrm{}`$. Thus, the relation (29) implies the gaussian tails of the probability density. It is clear that due to the finite energy flux and relations (10)-(11), two-dimensional turbulence cannot be a gaussian process. All the relation (31) can tell us is that the even-order moments with $`n>>1`$, described by (31), can be close to the gaussian values. It will be shown below that transverse velocity differences, not directly involved in the inter-scale energy transfer, can obey gaussian statistics.
It is clear from (17) that the model (29) for the pressure contributions is wrong when the linear dissipation terms are added to the Navier-Stokes equations. In the limit of small $`\mathrm{\Delta }u\mathrm{\Delta }v`$ the balance is achieved when:
$$<𝒫_{yv}|\mathrm{\Delta }u,\mathrm{\Delta }v,r>h\frac{\mathrm{\Delta }u\mathrm{\Delta }v}{r}(\alpha +\frac{b}{(Pr)^{\frac{2}{3}}})\mathrm{\Delta }v$$
which differs from (29) in the range of small $`\mathrm{\Delta }u\mathrm{\Delta }v`$.
$`\mathrm{𝟑}𝐝`$. In the intermittent three-dimensional turbulence $`\xi _{2n}<\xi _{2n+1}`$. This produces strong restrictions on the structure of the pressure contributions to the equation (6). Let us assume that $`\xi _{2n,0}=\xi _{2,2n2}=\xi _{2n2,2}`$. Then, it is clear from (15) that the first term of expansion (28) has all right properties. The relation (15), involving the $`r`$-derivatives, is valid for an arbitrary value of $`n`$ and that is why any additional term of expansion (28) must not only depend on a proper power of $`n`$ but the functions $`\kappa _{m,n}(r)`$ must also reflect non-trivial dimensionalities caused by the anomalous scaling exponents $`\xi _n`$. It cannot be rigorously ruled out, though this possibility seems quite bizzarre. In what follows we will adopt the pressure model (29) in the three-dimensional case also.
## 6 Two-dimensional turbulence
Now we will be interested in the case of the two-dimensional turbulence in the inverse cascade range. If a two-dimensional (2d) fluid is stirred by a random (or non-random) forcing, acting on a scale $`l_f=1/k_f`$, the produced energy is spent on creation of the large-scale ($`l>l_f`$) flow which cannot be dissipated in the limit of large Reynolds number as $`\nu 0`$. This is a direct and most important consequence of an additional, enstrophy conservation, law, characteristic of two dimensional hydrodynamics . As a result, the dissipation terms are irrelevant in the inverse cascade range and we set $`D=0`$ in (6) and hope that in two dimensions the situation is greatly simplified. This hope is supported by recent numerical and physical experiments - showing that as long as the integral scale $`L_it^{\frac{3}{2}}`$ is much smaller than the size of the system, the velocity field at the scales $`L_i>>l>>l_f`$ is a stationary close-to-gaussian process characterized by the structure functions with the Kolmogorov exponents $`\xi _n=n/3`$. In a recent paper Bofetta, Celani and Vegrassola reported the results of very accurate numerical simulations of two-dimensional turbulence generated by a random force. No deviations from gaussian statistics of transverse velocity differences as well as from the Kolmogorov scaling $`\xi _n=n/3`$ were detected. and is equal to zero if $`m=2n`$.
The pressure gradient $`_yp=_y_i_j^2\mathrm{\Delta }v_i\mathrm{\Delta }v_j`$ and the difficulty in calculation of $`I_p`$ is in the integral over the entire space defined by the inverse Laplacian $`^2`$. The huge simplification, valid in 2d, comes from the fact that all contributions to the left side of equation (6) as well as $`I_f`$ are independent on time. This means that the integrals involved in the pressure terms cannot be infra-red divergent since in a two-dimensional flow $`L=L(t)t^{\frac{3}{2}}`$. We also have that $`i_p\alpha ^2i_p`$ when $`U,V\alpha U;\alpha V`$. Based on this and taking into account that $`<(\mathrm{\Delta }v)^{2n+1}(u_x^2+v_y^2+u_yv_x)>=0`$, we, in the limit $`\eta _20`$, adopt a low-order model (29) giving:
$$I_p=[h\frac{^2}{\eta _2\eta _3}+b\frac{\eta _3}{(Pr)^{\frac{2}{3}}}_3]Z(\eta _2=0,\eta _3,r)$$
(32)
In two dimensions the relation (32) combined with (6) in the limit $`\eta _20`$, solves the problem 2d turbulence.
Substituting (32), (29) into (6) and, based on (9)-(11), seeking a solution as $`\eta _20`$ as
$$Z(\eta _2,\eta _3,r)Z_3(\eta _3,r)\phi (\eta _2r^{\frac{1}{3}},\eta _2)Z_3(\eta _3,r)exp(\frac{1}{2}A_{2,0}(\eta _2Pr^{\frac{1}{3}})^2)(1+\frac{1}{2}A_{1,2}\eta _3^2\eta _2(Pr)+\frac{\eta _2^3}{4}+\mathrm{})$$
(33)
where ($`A_{1,2}=1/2`$), gives:
$$[_r+\frac{1}{r}+\frac{1h}{r}\eta _3_3]\frac{A_{1,2}Pr}{2}\eta _3^2Z_3=2P\eta _3^2Z_3+\frac{b\eta _3}{r^{\frac{2}{3}}}_3Z_3$$
(34)
Setting $`Z_3=Z_3(\eta _3r^{\frac{1}{3}})Z_3(X)`$ and $`h=4/3`$ one derives using, the relation (30) ($`b=\frac{2}{3A_{0,2}}`$):
$$2A_{0,2}XZ_3=_XZ_3$$
corresponding to a Gaussian solution with the correct width $`A_{0,2}`$. This fact serves as a consistency check that the gaussian is a solution for the PDF of transverse velocity differences.
The equation (34) defines a probability density function corresponding to the finite moments $`S_{2m,2n}(r)`$ only when $`h=4/3`$. This situation resembles Polyakov’s theory of Burgers turbulence reduced to an eigenvalue problem with a single eigenvalue corresponding to the pdf which is positive in the entire interval.
Having these exact results, and keeping in mind (33) one can integrate the equation over $`\eta _2`$ from $`i\mathrm{}`$ to $`0`$ to obtain :
$$\frac{Z_3}{r}+3(1hb)\frac{\eta _3}{r}\frac{Z_3}{\eta _3}=\frac{2P}{(Pr)^{\frac{1}{3}}}\eta _3^2Z_3$$
(35)
valid as long as
$$\frac{(\eta _3r^{\frac{1}{3}})^2}{8A_{2,0}^2}<<1$$
This constraint is an artifact of an approximate relation (33). As will be shown below (35) gives an exact gaussian solution and thus, is valid beyond above interval. This result is obtained choosing the integration function $`\mathrm{\Psi }(\eta _3,r)`$ to compensate the $`O(Z/r)`$ term violating the normalizability constraint $`Z(0,0,r)=1`$. Solution to (35) is:
$$Z_3=exp(\gamma \eta _3^2(Pr)^{\frac{2}{3}})$$
(36)
with the parameter $`\gamma =\frac{3}{(3(1hb)+1)}A_{0,2}`$ defining the width of the gaussian.
The first-order differential equation (35) for the generating function differences implies the underlying linear Langevin dynamics of transverse velocity differences. It is important that this equation in non-local in physical space but local in the Fourier one. The effective forcing, corresponding to the right-side of (35), is a non-local and solution- dependent.
To evaluate the single-point probability density, corresponding to velocity differences in the limit of large displacements $`r`$ , we can notice that the energy flux is not equal to zero only at $`r<<L`$. At the distances $`rL(t)`$ the zero value of the energy flux and symmetrization of the probability density ($`P(\mathrm{\Delta }u,L)=P(\mathrm{\Delta }u,L)`$ can be achieved only when the pressure contribution to the equation (6) compensates the advective terms. As a result, since $`D=0`$, we have:
$$Z_t=2P\eta _3^2Z$$
Seeking a solution as $`Z=Z(\eta _3\sqrt{t})Z(\eta _3v_{rms}(t))`$ gives a gaussian result:
$$Z=e^{\eta _3^2u_{rms}^2(t)}$$
Similar outcome is obtained for the case investigated in . When turbulence is stabilized at the large scales by an artificially introduced friction, the resulting equation is:
$$\eta _3_3Z=2p\eta _3^2Z$$
also leading to the gaussian pdf. To conclude this section we would like to discuss the physical meaning of the integral scale $`L`$. The integral scale of turbulence is a scale at which the flux decreases to zero and at which $`S_{3,0}(L)=0`$.
## 7 Small Parameter in Turbulence Theory in Three Dimensions
Three-dimensional turbulence is a notoriously difficult problem due to absence of the small parameter. This can be illustrated on an example of a coarse-graining procedure, which was extremely successful in engineering turbulence simulations. Consider the wave-number $`k`$ in the inertial range. Let us denote $`𝐯^<(𝐤)`$ the Fourier components of the velocity field with all modes $`𝐯^>(𝐪)=0`$ when $`q>k`$. The coarse-grained field in the physical space is defined then:
$$𝐯_r(𝐱,t)=_{k<\frac{1}{r}}𝐯^<(𝐤)e^{i𝐤𝐱}d^3k$$
The equation of motion for $`v_i^<(k)`$ in the Fourier space resembles the Navier-Stokes equation with effective viscosity
$$\nu (k)(\frac{(dd_c)(d+\frac{1}{2})}{3.2d(d+2)})^{\frac{1}{3}}^{\frac{1}{3}}k^{\frac{4}{3}}$$
with $`d_c2.56`$ plus high order non-linearities. The parameter $`3.2`$ in the above relation, evaluated at $`d=3`$, is in fact a weak function of space dimensionality and $`=P=O(1)`$. This expression is derived assuming a close-to-gaussian statistics of the small-scale turbulence with
$$\overline{v_i(k,\omega )v_j(k^{},\omega ^{})}\frac{k^d}{i\omega +\nu (k)k^2}\delta (k+k^{})\delta (\omega +\omega ^{})$$
which is accurate at $`dd_c0`$ (see below).
In the physical space the effective viscosity of the coarse-grained field:
$$\nu _r(dd_c)^{\frac{1}{3}}N(r^4)^{\frac{1}{3}}v_r^2\tau _r$$
defines the relaxation time $`\tau _r`$ which is a characteristic time of interaction of the field $`v_r`$ with the eliminated modes acting on the scales $`l<r`$.
The difficulty of the theory is in the higher non-linearities
$$(v_r)^n\tau _r^{n1}v_r$$
and the dimensinless expansion parameter
$$\frac{\tau _r}{\theta _r}_rv_r\tau _r\frac{\tau _rv_r}{r}$$
is nothing but the ratio of translational time $`\theta _r`$, characterizing the tendency of the “large-scale” longitudinal velocity fluctuation at the scale $`r`$ to form a small-scale structure (“shock”) in the absence of pressure, to the relaxation time strongly influenced by the pressure gradients contributions. In both 2d and 3d these times are of the same order and that is why trancation of the expansion is a very difficult problem.
This is not so in the vicinity of $`d=d_c`$. Let us assume that the theory can be analytically continued to the non-integer dimensions . Then, since the energy flux is $`O(1)`$ we have:
$$\frac{}{r}<\mathrm{\Delta }u(\mathrm{\Delta }v)^2>\frac{}{r}<(\mathrm{\Delta }u)^3>\nu <(\frac{v_{ri}}{r_j})^2>=O(1)$$
This means that
$$v_{r,rms}(dd_c)^{\frac{1}{6}}(r)^{\frac{1}{3}}$$
the dissipation wave-number $`k_d(dd_c)^{\frac{1}{4}}(\frac{}{\nu _0^3})^{\frac{1}{4}}`$ and, as a consequence,
$$\frac{\tau _r}{\theta _r}(dd_c)^{\frac{1}{2}}$$
which will serve as a small parameter of the theory when $`dd_c0+`$.
These relations tell us that the turbulent intensity grows to infinity with $`dd_c0`$ where the energy flux changes its sign. The time needed to reach the steady stae is estimated easily:
$$T(dd_c)^{\frac{1}{3}}^{\frac{1}{3}}r^{\frac{2}{3}}$$
after which a close-to-Kolmogorov spectrum is expected both above and below $`d_c`$. Thus, at $`d=d_c`$ the the flow is unsteady.
The above results enable one to derive a plausable estimate for the dissipation term $`D`$. It is clear from the Navier-Stokes equations that:
$$=\frac{1}{2}_iv_iv_j^2\frac{1}{2}v_t^2_iv_ip$$
The coarsed-grained expression in the law frequency limit
$$\frac{1}{2}\frac{}{r_r}v_{ri}v_{rj}^2(1+O(dd_c))$$
(37)
To arrive in the expression for $`D`$ we assume $`v_r\mathrm{\Delta }v`$ leading to an expression very similar to Kolmogorov’s refined similarity hypothesis.
## 8 Three-Dimensional Flow
The most important feature of two-dimensional turbulence, considered in a previous section, is irrelevance of the dissipation processes in the inverse cascade range when $`d<d_c`$. It is this irrelevance which was responsible for the gaussian probability density of transverse velocity differences.
The model for the dissipation contribution $`D`$ in the limit $`\eta _20`$ is evaluated from (37) readily. We would like to keep at least some information about $`\mathrm{\Delta }v`$ and the expression must be invariant under transformation $`𝐯𝐯`$ and $`𝐱𝐱`$. In addition, the expression must be local in physical space. Based on this considerations:
$$_vc(d)\mathrm{\Delta }u\mathrm{\Delta }v\frac{\mathrm{\Delta }v}{r}$$
where $`_v`$ is a dissipation rate of the “v-component to kitenic energy” $`K_v=\frac{1}{2}v^2`$. Locality of this model is clear since $`_r\mathrm{\Delta }v=_1v(x_1)+_2v(x_2)`$. The problem is in evaluation of the coefficient $`c(d)`$ since, in principle, it can be singular at $`d=d_c`$. Indeed, it is clear that the dissipation term is zero at $`dd_c`$. However, the point $`d=d_c`$ separating inverse and direct cascade ranges, is a singularity due to infinitly large amplitudes of velocity fluctuations. All we can say at this point is that the pdf can be represented as a sum of even and odd functions of $`\mathrm{\Delta }v`$. The symmetric part has a growing with $`dd_c0+`$ width, while the width of odd one is $`O(1)`$. The behaviour of $`c(d)`$ in the vicinity of $`d_c`$ is not clear. We feel that it is $`O(1)`$ at $`d>d_c`$ and zero at $`dd_c`$.
Thus, we have:
$$Dc(d)\eta _3_{\eta _2}_{\eta _3}_rZc(d)\eta _3^2<\mathrm{\Delta }u\mathrm{\Delta }v\frac{\mathrm{\Delta }v}{r}e^{\eta _2\mathrm{\Delta }u+\eta _3\mathrm{\Delta }v}>+O(_r_{\eta _2}Z)$$
This expression obeys the basic symmetries of the Navier-Stokes equation. The expression, similar to (38), was being used in engineering turbulent modelling based on the low-order in the Reynolds number coarse-grained equations.
The last term in the right side of (38), simply modifies the coefficient in front of the first term in the left side the equation (6) and does not generate anything new. The expression (38) resembles Kolmogorov’s refined similarity hypothesis, connecting the dissipation rate, averaged over a region of radius $`r`$, with $`(\mathrm{\Delta }u)^3`$. Thus, in the limit $`\eta _20`$:
$$[_{\eta _1}_{\eta _2}+\frac{2}{r}_{\eta _2}+(1+h)\frac{\eta _3}{r}\frac{^2}{_{\eta _2}\eta _3}+c(d)\eta _3_{\eta _2}_{\eta _3}_r]Z(\eta _2=0,\eta _3,r)=\frac{2P}{3}(1cos(k_fr))Zb\frac{\eta _3}{r^{\frac{2}{3}}}_3Z$$
The $`dd_c>0`$ counterpart of equation (35) is:
$$[_{\eta _1}+(1+h+b)\frac{\eta _3}{r}\frac{}{\eta _3}+c(d)\eta _3_{\eta _3}_r]Z(\eta _2=0,\eta _3,r)=a\frac{2P}{3}\frac{1cos(k_fr)}{r^{\frac{1}{3}}}\eta _3^2Z_3$$
(38)
with $`\mathrm{\Psi }(\eta _3)`$ chosen in such a way that the generating function $`Z(0,0,r)=1`$. We consider two limiting cases.
$`\mathrm{𝐒𝐦𝐚𝐥𝐥}\mathrm{𝐒𝐜𝐚𝐥𝐞𝐃𝐲𝐧𝐚𝐦𝐢𝐜𝐬}`$.
Inverse Laplace transform of (40) without right side gives an equation for the pdf $`P(\mathrm{\Delta }v,r)`$:
$$\frac{P}{r}+\frac{1+3\beta }{3r}\frac{}{V}VP\beta \frac{}{V}V\frac{P}{r}=0$$
(39)
where $`\beta c(d)`$. Since $`S_{0,3}=0`$, the coefficients in (41) are chosen to give $`s_{0,3}=\overline{|\mathrm{\Delta }v|^3}=a_3Pr`$ with an undetermined amplitude $`a_3`$. This is an assumption of the present theory, not based on rigorous theoretical considerations. Seeking a solution in a form $`S_{0,n}=<(\mathrm{\Delta }v)^n>r^{\xi _n}`$ gives:
$$\xi _n=\frac{1+3\beta }{3(1+\beta n)}n\frac{1.15}{3(1+0.05n)}n$$
(40)
which was derived in together with $`\beta 0.05`$. It follows from (42) that: $`P(0,r)r^\kappa `$ where $`\kappa =\frac{1+3\beta }{3(1\beta )}0.4`$ for $`\beta =0.05`$. Very often the experimental data are presented as $`P(X,r)`$ where $`X=V/r^\mu `$ with $`2\mu =\xi _20.696`$ for $`\beta =0.05`$. This gives $`P(X=0,r)r^{\kappa +\mu }r^{0.052}`$ compared with the experimental data by Sreenivassan : $`\kappa +\mu 0.06`$.
Let us write $`P(V,r)=r^\kappa F(\frac{V}{r^\kappa },r)=r^\kappa F(Y,r)`$, so that $`F`$ obeys the following equation:
$$(1\beta )r\frac{F}{r}+\beta \kappa \frac{}{Y}Y^2\frac{F}{Y}\beta Yr\frac{^2F}{Yr}=0$$
(41)
Next, changing the variables again $`\mathrm{}<y=Ln(Y)<\mathrm{}`$, substituting this into (43) and evaluating the Fourier transform of the resulting equation gives:
$$(1\beta )r\frac{F}{r}+\beta \kappa (ikk^2)Fik\beta r\frac{F}{r}=0$$
(42)
with the result: $`Fr^{\gamma (k)}`$, where $`\gamma (k)=\beta \kappa \frac{ik+k^2}{1\beta i\beta k}Ln(r/L)`$ with $`r/L<<1`$. We have to evaluate the inverse Fourier transform:
$$F=_{\mathrm{}}^{\mathrm{}}𝑑ke^{iky}e^{\gamma (k)}$$
(43)
in the limit $`y=O(1)`$ and $`r0`$ so that $`Ln(r/L)\mathrm{}`$. The integral can be calculated exactly. However, the resulting expression is very involved. Expanding the denominator of $`\gamma (k)`$ gives :
$$F=_{\mathrm{}}^{\mathrm{}}𝑑ke^{ik(y+\frac{\beta \kappa (Ln(r))}{1\beta })}e^{\frac{\beta \kappa (1+\beta )|Ln(r)|}{(1\beta )}k^2}$$
(44)
and
$$F\frac{1}{\sqrt{\mathrm{\Omega }(r)}}exp(\frac{(Ln(\xi ))^2}{4\mathrm{\Omega }})$$
(45)
with $`\xi =V/r^{\frac{\kappa }{1\beta }}`$ and $`\mathrm{\Omega }(r)=4\beta \kappa \frac{1+\beta }{1\beta }|Ln(r/L)|`$.
To understand the range of validity of this expression, let us evaluate $`<V^n>`$ using the expression (47) for the pdf. Simple integration, neglecting $`O(\beta ^2)`$ contributions, gives: $`<V^n>r^{\alpha _n}`$ with $`\alpha _n=(1+3\beta )(n\beta (n^2+2))/3`$. Comparing this relation with the exact result (43) we conclude that the expression for the pdf, calculated above, is valid in the range $`n>>1`$ and $`\beta n<<1`$. The properties of the pdf in the range $`3\xi 15`$ are demonstrated on Fig. 1 for $`r/L=0.1;0.01;0.001`$. The log-normal distribution (48), is valid in a certain (wide but limited) range of the $`V`$\- variation. It is clear from (42) that neglecting the dissipation terms ($`c(dd_c)\beta =0`$) leads to $`\xi _n=n/3`$, i. e. disappearance of anomalous scaling of moments of velocity differences. This result agrees with the well-developed phenomenology, attributing intermittency to the dissipation rate fluctuations: the stronger the fluctuations, the smaller the fraction of the total space they occupy , . To the best of our knowledge, this is the first work leading to multifractal distribution of velocity differences as a result of approximations made directly on the Navier-Stokes equations.
To investigate the probability density $`P(Y,r)`$ in the limit $`Y0`$ we introduce an expansion:
$$F(Y,r)=\underset{n}{}C_nY^{2n}f_{2n}(r)$$
(46)
Substituting this into (44) gives:
$$f_{2n}(\frac{r}{L})^{\frac{\beta \kappa 2n(2n1)}{1\beta (1+n)}}$$
(47)
It is seen from (48)-(49) that the pdf starts bending from the log-normal slope (47) toward $`_YF(Y,r)=0`$ at $`Y=0`$ at:
$$Y<(\frac{r}{L})^{\frac{\beta \kappa }{12\beta }}$$
(48)
This inequality shows that as $`r0`$ the pdf develops a narrow cusp at the origin $`Y=0`$. If the probability density is plotted in the dimensionless variable $`X`$, the bending starts at $`Xr^{0.07}`$. This value was calculated, as above, with $`\beta =0.05`$.
Large-scale limit: $`r/L1`$. Now, let us investigate the large-scale limit $`\frac{r}{L}1`$. Realizing that $`\beta `$ can be an $`O(1)`$ constant we, for the illustration purposes, will investigate the large scale limit pretending that $`\beta (d)0`$ as $`dd_c`$. This is also useful since the estimated value of $`\beta 0.05`$ at $`d=3`$ which is numerically small. In this limit the right side of (41) is $`O(\eta _3^2Z)`$ and cannot be neglected. Repeating the procedure leads to an equation:
$$\frac{P}{r}+\frac{1+3\beta }{3r}\frac{}{V}VP\beta \frac{}{V}V\frac{P}{r}=a\frac{P}{(\mathrm{Pr})^{\frac{1}{3}}}\frac{^2Z}{V^2}$$
(49)
where $`a`$ is a proportionality coefficient and $`rL`$. As one can see from this equation in the limit of small $`\beta `$ the solution to this equation approaches gaussian. It is also clear that for any finite $`\beta `$, the tails of the pdf are strongly non-gaussian when
$$\beta Y^2>>1$$
(50)
This estimate means that, according to the theory presented above, the perturbative treatment of deviations from the mean-field gaussian theory is possible but it involves two parameters: the ratio $`ϵ=1\frac{r}{L}<<1`$ and $`\beta <<1`$. The fact that the “real-life” $`\beta 1/20`$ may explain why the experimentally observed pdf of the large-scale ($`ϵ<<1`$ ) velocity fluctuations was so close to the gaussian \[see and references therein\].
It is also seen from (52) that at $`(\frac{r}{L})^2\beta 0.05`$ the pdf is dominated by a gaussian central part.
## 9 Conclusions
The equation (6) formulates theory of turbulence in terms of “only” two unknowns: pressure and dissipation terms $`I_p`$ and $`D`$, respectively. It provides a mathematical testing ground for various analytic expressions and models obtained from numerical simulations.
Armed with the experimental and numerical data, supporting gaussian statistics of transverse velocity differenced in two-dimensional flows, we showed that the mean field approximation (the lowest-order term of the expansion (28)) for the pressure contributions (29) leads to both Kolmogorov scaling and gaussian statistics of transverse velocity differences. In addition, the equation (6) shows that the single-point pdf’s in $`2d`$ turbulence are gaussian. It is to be stressed that $`2d`$ turbulence cannot be a gaussian process and probability density $`P(\mathrm{\Delta }u,\mathrm{\Delta }v,r)=`$ is not a gaussian. It is only pdf $`P(\mathrm{\Delta }v,r)=P(\mathrm{\Delta }u,\mathrm{\Delta }v,r)𝑑\mathrm{\Delta }u`$ is a gaussian. This statement violates no dynamic constraints.
One of the most interesting outcomes of the present theory is a discovery of existence of the two time-scales in the system which are very different in the vicinity of $`d=d_c`$. This difference enables one to coarse-grain the Navier-Stokes equations and neglect all high-order non-linearities, generated by the procedure. Using this result the model for for the dissipation term $`D`$ was derived.
The as yet unresolved ambiguity of this model is its behaviour as $`dd_c`$. If transition from 3d to the non-intermittent state at $`d<d_c`$ is smooth, then $`\beta 0`$ and the resulting equation shows onset of both anomalous scaling and non-gaussian statistics. The transition can be singular, however: right at $`d>d_c`$ the coefficient $`\beta `$ can become $`O(1)`$ and a weakly intermittent state and weak coupling limit do not exist. In this case, due to existence of the small parameter, enabling evaluation of the dissipation expression $`D`$, the theory non-perturbatively predicts both the shape of the pdf and scaling exponents provided the small parameter $`S_{3,0}/(S_{2,0})^{\frac{3}{2}}0`$ as $`dd_c`$. This result is possible since $`D=D_0+O(dd_c)`$ avd even in the limit $`dd_c`$, the model $`D=O(1)`$. At $`dd_c<0`$ , $`D=0`$ leading to the gaussian pdf of transverse velocity differences. Experimental investigation of hydrodynamics in a non-integer space dimension is impossible. However, it was demonstarted by Jensen that a force-driven shell model is yields the changing sign of the energy flux upon variation of a leading parameter. Numerical solution at a critical point (zero flux) demonstrated an unsteady state with the growing total energy and the energy spectrum concentrated in the vicinity of $`k_f`$. The calculat ion also gave kolmogorov energy spectrum at $`{}_{}{}^{\prime \prime }d>d_c^{\prime \prime }`$ with growing Kolmogorov constant as $`dd_c`$. It is not yet clear how the eddy viscosity approximation works for the shell model, but, since the $`O(1)`$ energy flux is fixed by the forcing function, the growth of kinetic energy must be related to a relaxation time $`\tau _r0`$ at $`d=d_c`$ and a corresponding small parameter. This result also shows that the phenomenon is very robust: all one needs is a point at which energy cascade changes its direction. The results of the shell model investigation will be published elsewhere .
The expression (47) is similar to the one obtained in a ground-breaking paper by Polyakov on the scale-invariance of strong interactions, where the multifractal scaling and the pdf were analytically derived for the first time -. In the review paper Polyakov noticed that the exact result can be simply reproduced considering a cascade process with a heavy stream (particle) transformed into lighter streams at each step of the cascade (fission). Due to the relativistic effects the higher the energy of the particle, the smaller the angle of a cone, accessible to the fragments formed as a result of fission. Thus, the larger the number of a cascade step, the smaller is the fraction of space occupied by the particles .
The theory, presented in this paper describes many experimental observations. Still, understanding of the limits of validity of expression (29) is crucial for the final assessment of the theory. The relation (32) shows that (29) is consistent with the gaussian tails of the pdf. However, at the present stage we are unable to prove that (29) is the only expression leading to this result. The problem is that without experimental detection of at least some deviations from the the gaussian statistics of transverse velocity differences, one will not be able to understand the limits of validity of (29). Given the state-of-the-art of numerical simulations, this goal may not be that simple.
## 10 Acknowledgment
I am grateful to A. Polyakov for most interesting and stimulating discussions. Help from K.R. Sreenivasan and M. Vergassola, who shared with me their sometimes unpublished experimental data , was most usefull. Recent numerical results of M. Jensen produced not only support of conclusions of this paper but also demonstrated a surprising power and importance of a shell model. I benefited a lot from conversations with B. Shraiman, P. Constantin, C. Doering, T. Gotoh and T. Kambe.
## 11 References
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9. U. Frisch and J.D. Fournier, Phys. Rev. A, 17, 747 (1978)
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13. L.D.Landau and E.M. Lifshitz, Fluid Mechanics, Pergamon Press, Oxford, 198 14. A.S.Monin and A.M.Yaglom, “Statistical Fluid Mechanics” vol. 1, MIT Press, Cambridge, MA (1971)
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V. Yakhot, Phys. Rev.E, Phys. Rev. E, 57, 1737 (1998)
19. K.R. Sreenivasan, private communication
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warning/0001/nucl-th0001023.html | ar5iv | text | # Shell Corrections for Finite-Depth Deformed Potentials: Green’s Function Oscillator Expansion Method
## I Introduction
The positive-energy spectrum of the average single-particle potential plays a role in the description of weakly bound nuclei for which the Fermi level approaches zero (see Ref. ). For these nuclei, important for both nuclear structure and nuclear astrophysics studies, special care should be taken when dealing with the particle continuum which seriously impacts many nuclear properties, including bulk nuclear properties (e.g., masses, radii, shapes) as well as nuclear dynamics (i.e., excitation modes).
In two earlier papers , a macroscopic-microscopic method was applied to nuclei far from the beta stability line. It has been demonstrated that the positive-energy single-particle spectrum does impact the results significantly, and the systematic error in binding energies, due to the neglect or the improper treatment of the particle continuum, can be as large as several MeV at the neutron drip line. In the first paper , both spherical and deformed nuclei were considered, and the continuum was approximated by a limited number of quasistationary states which resulted from a diagonalization of the Woods-Saxon average potential in a harmonic oscillator basis. In the vicinity of the neutron drip line, the result of the Strutinsky smoothing (standard averaging method) becomes unreliable and it deviates from the result of the semiclassical Wigner-Kirkwood expansion. The semiclassical method, which does not use the positive-energy spectrum explicitly, gives a more reliable estimate of the shell correction than the standard method (see Refs. and references quoted therein).
In the following study , carried out for spherically symmetric nuclei, a more detailed comparison was carried out between the Strutinsky smoothing method and the Wigner-Kirkwood expansion. In the Strutinsky method, the continuum effect was taken into account exactly, i.e., by calculating the continuum part of the level density from the derivative of the scattering phase shift with respect to single-particle energy. The smooth part of the continuum level density has been calculated by means of the contour integration along a path in the complex energy plane . Although the continuum level density was treated properly, it has been concluded that in most nuclei the plateau condition of the Strutinsky method could not be met. Therefore, this condition was replaced with the requirement of the linear energy dependence of the mean level density. This modification (which widens the range of the applicability of the Strutinsky procedure considerably) is referred to as the generalized Strutinsky method. The new procedure has proved to be very useful; in most cases it gives results reasonably close to the estimate of the semiclassical method. The exceptions are the neutron drip-line nuclei in which the neutron Fermi level approaches zero and the semiclassical procedure diverges .
In the present work, the effect of the particle continuum on shell correction is further studied for both spherical and deformed nuclei by using the recently introduced Green’s function approach . The advantage of this method is that, by diagonalizing the finite single-particle potential in a square-integrable basis, one can get rid of the spurious contribution of the particle gas to the level density. We shall refer to this novel method as the new method in order to distinguish it from the commonly used standard smoothing procedure (or: old method) in which the spurious contribution from the particle gas is not subtracted (but diminished by using a reduced number of basis states).
The paper is organized as follows. Section II contains a brief review of the shell-correction method and describes several versions of the smoothing procedure used. The numerical results of shell-correction calculations are presented in Sec. III for both neutron-rich and proton-rich nuclei. Finally, conclusions are drawn in Sec. IV.
## II Strutinsky Smoothing Procedure
### A Basic Definitions
In the macroscopic-microscopic approach , the shell correction,
$$\delta E_{\mathrm{shell}}=E_{\mathrm{s}.\mathrm{p}.}\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.},$$
(1)
is defined as the difference between the total single-particle energy $`E_{\mathrm{s}.\mathrm{p}.}`$,
$$E_{\mathrm{s}.\mathrm{p}.}=\underset{i\mathrm{occ}}{}ϵ_i,$$
(2)
and the smooth single-particle energy $`\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.}`$. The shell correction represents the fluctuating part of the binding energy resulting from the quantal single-particle shell structure.
For the sake of simplicity, we shall assume that the single-nucleon energy spectrum is that of a one-body Hamiltonian,
$$\widehat{H}=\widehat{T}+\widehat{V},$$
(3)
with an average single-particle potential $`\widehat{V}`$. In practice, the potential contains a deformed Woods-Saxon potential, the spin-orbit term and the Coulomb potential. Since the central potential is finite, the spectrum of $`\widehat{H}`$ is composed of bound states with discrete negative eigenvalues ($`ϵ_i<0`$) and the continuum of scattering states with positive energies ($`ϵ>0`$). Consequently, the single-particle level density is
$$g(ϵ)=g_\mathrm{d}(ϵ)+g_\mathrm{c}(ϵ),$$
(4)
where
$$g_\mathrm{d}(ϵ)=\underset{i}{}2\delta (ϵϵ_i)$$
(5)
is the level density of the discrete (bound) states and $`g_\mathrm{c}(ϵ)`$ is the continuum level density (it will be specified later). The factor $`2`$ in Eq. (5) appears due to the two-fold Kramers degeneracy of the deformed single-particle energy levels.
In the shell-correction method , $`\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.}`$ is calculated by employing the smoothed level density $`\stackrel{~}{g}(ϵ)`$ obtained from $`g(ϵ)`$ by folding it with a smoothing function $`f(x)`$:
$$\stackrel{~}{g}(ϵ)=\frac{1}{\gamma }_{\mathrm{}}^+\mathrm{}𝑑ϵ^{}g(ϵ^{})f\left(\frac{ϵ^{}ϵ}{\gamma }\right).$$
(6)
In practical applications, for the folding function $`f(x)`$ one is usually taking a product of a Gaussian weighting function, $`\frac{1}{\sqrt{\pi }}\mathrm{exp}(x^2)`$, and a corresponding curvature correction polynomial of the order $`p`$ which is an associated Laguerre polynomial $`L_{p/2}^{1/2}(x)`$ ($`p`$-even). The smoothed level density (6) defines both the smooth single-particle energy
$$\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.}=_{\mathrm{}}^{\stackrel{~}{\lambda }}ϵ\stackrel{~}{g}(ϵ)𝑑ϵ,$$
(7)
and the smoothed Fermi level $`\stackrel{~}{\lambda }`$. The latter is obtained from the particle number equation:
$$N=_{\mathrm{}}^{\stackrel{~}{\lambda }}\stackrel{~}{g}(ϵ)𝑑ϵ.$$
(8)
The smooth single-particle energy and the resulting shell correction have to be defined unambigously. Therefore, they must neither depend on the smoothing range $`\gamma `$ nor on the order $`p`$ of the curvature correction. This requirement, referred to as the plateau condition, can be written as
$$\frac{d\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.}}{d\gamma }=0,\frac{d\stackrel{~}{E}_{\mathrm{s}.\mathrm{p}.}}{dp}=0.$$
(9)
Of course, since one wants to eliminate the oscillations due to the shell structure, the smoothing range $`\gamma `$ should be greater than the average energy distance between neighboring major shells, $`\mathrm{}\omega _0`$$``$41/$`A^{1/3}`$ MeV .
For infinite potentials, such as an infinite square well, harmonic oscillator, and a deformed Nilsson potential, one can always find a range of the smoothing parameters $`\gamma `$ and $`p`$ in which the smooth single-particle energy is independent of the values $`\gamma `$ and $`p`$ . For finite-depth potentials, however, additional complications arise due to (i) the presence of positive-energy continuum and (ii) the difficulties with meeting the plateau condition. We shall discuss these points in the following.
### B Effect of the Unbound Spectrum
The need for calculating the continuum level density, $`g_\mathrm{c}(ϵ)`$, appears whenever one deals with finite-depth potentials. For spherically symmetric potentials, the continuum level density is defined by means of the scattering phase shifts $`\delta _{lj}(ϵ)`$:
$$g_\mathrm{c}(ϵ)=\frac{1}{\pi }\underset{l,j}{}(2j+1)\frac{d\delta _{lj}(ϵ)}{dϵ}.$$
(10)
For realistic nuclear potentials, phase shifts have to be calculated by numerically solving the radial Schroedinger equations and by matching the wave function to the asymptotic solution at a distance where the nuclear potential can be neglected. This procedure has to be carried out for every partial wave below a certain angular momentum cut-off on a fine mesh in the positive-energy region. In order to prevent sudden jumps in $`g_\mathrm{c}(ϵ)`$ around narrow resonances, a new calculational method employing the Cauchy theorem was introduced in Ref. . Here, the complex energies $`w_i`$ of the Gamow resonances (poles of the S-matrix) are localized first, then a contour of the complex energy plane is chosen. The contour, denoted by $`L`$, should go far away from the poles. The mean level density $`\stackrel{~}{g}(ϵ)`$ is then calculated as a sum over bound and those resonant states which lie between $`L`$ and the real energy axis and an integral term along a contour:
$$\stackrel{~}{g}(ϵ)=\underset{i}{}f\left(\frac{ϵw_i}{\gamma }\right)+_L𝑑wg_\mathrm{c}(w)f\left(\frac{ϵw}{\gamma }\right).$$
(11)
Apart from the numerical errors, this procedure gives the continuum level density exactly. Therefore, we call this approach as numerically exact or, simply, exact.
For deformed single-particle potentials, the continuum level density has a more complicated form and can be expressed by the on-shell S-matrix $`S(ϵ,\widehat{k},\widehat{k}^{})`$ as
$$g_\mathrm{c}(ϵ)=\frac{1}{2i\pi }\mathrm{Tr}[S(ϵ,\widehat{k},\widehat{k}^{})^{}\frac{d}{dϵ}S(ϵ,\widehat{k},\widehat{k}^{})].$$
(12)
If one wants to use expression (12) for calculating the continuum level density, one has to determine the S-matrix by solving the coupled system of differential equations for each value of $`ϵ`$ . In practice this is a difficult task.
In the old method, the single-particle Hamiltonian is diagonalized in a square-integrable basis formed from the eigenstates of an infinite potential. This potential can be either a finite-range potential contained in an impenetrable box (i.e., an infinite wall at a certain distance) or the harmonic oscillator potential. Since the number of basis states is always assumed to be finite, the diagonalization results in a discrete set of eigenstates. The eigenstates with negative energy approximate the bound states of the original Hamiltonian, while the positive-energy quasi-bound states mock-up the effect of the particle continuum in a very crude way. It has early been realized that in the application of this method one should not use too large a basis; otherwise, the level density around the zero energy would increase dramatically. (In fact it diverges as the basis size goes to infinity.) In order to avoid this catastrophe, the use of a harmonic oscillator basis with 12-14 harmonic oscillator shells was recommended . One of the objectives of this paper is to perform the critical evaluation of the standard smoothing method by calculating the continuum level density in a more reliable way using the Green’s function approach.
### C Green’s Function Method
Although the Green’s function approach to the single-particle level density was developed long ago (see, e.g., Refs. ), it is somehow surprising that so far it has not been widely applied. Below, we briefly summarize the main features of this method. More details can be found in Refs. .
A Hamilton operator with an infinite potential, $`\widehat{H}_{\mathrm{}}`$, has only discrete energy eigenvalues and its eigenfunctions are all square integrable. Therefore, in this case, the single-particle level density is given by Eq. (5). By introducing the Green’s operator, $`\widehat{G}_{\mathrm{}}(z)=(z\widehat{H}_{\mathrm{}})^1`$, $`g_\mathrm{d}(ϵ)`$ can be written as
$$g_\mathrm{d}(ϵ)=\frac{1}{\pi }\mathrm{Im}\left\{\mathrm{Tr}\left[\widehat{G}_{\mathrm{}}(ϵ)\right]\right\}.$$
(13)
As discussed in Ref. , for a Hamiltonian $`\widehat{H}`$ containing a finite potential the full level density (4) becomes
$$g(ϵ)=\frac{1}{\pi }\mathrm{Im}\left\{\mathrm{Tr}\left[\widehat{G}^+(ϵ)\widehat{G}_0^+(ϵ)\right]\right\},$$
(14)
where $`\widehat{G}^+(z)=(z\widehat{H}+i0)^1`$ and $`\widehat{G}_0^+(z)=(z\widehat{H}_0+i0)^1`$ is the free outgoing Green’s operator associated with $`\widehat{H}_0=\widehat{T}`$. The interpretation of Eq. (14) is straightforward: the second term contains the contribution to the single-particle level density originating from the gas of free particles.
Let us now introduce an approximation to the exact expression (14). To this end, we diagonalize $`\widehat{H}`$ and $`\widehat{H}_0`$ in an orthonormal basis formed from the $`M`$ square-integrable basis functions. The resulting approximate eigenenergies of $`\widehat{H}`$ and $`\widehat{H}_0`$ are denoted by $`e_i`$ and $`e_i^0`$, respectively ($`i`$=1,…,$`M`$). This procedure amounts to a projection of both Hamiltonians into the $`M`$-dimensional Hilbert space of square-integrable basis functions. The level density (4) can then be approximated by the difference of the discrete level densities of the two projected Hamiltonians:
$$g_M(ϵ)=\underset{i=1}{\overset{M}{}}2\delta (ϵe_i)\underset{i=1}{\overset{M}{}}2\delta (ϵe_i^0).$$
(15)
By increasing the dimension $`M`$, the bound eigenvalues of $`\widehat{H}`$ converge to the exact single-particle energies while the positive-energy eigenvalues will tend to approach zero energy. The eigenvalues $`e_i^0`$, which are obviously different from the positive energies $`e_i`$ for any finite $`M`$, can, in fact, compensate for the spurious increase of the level density around the zero energy if the smoothing procedure (6) is applied:
$$\stackrel{~}{g}_M(ϵ)=\frac{1}{\gamma }_{\mathrm{}}^+\mathrm{}𝑑ϵ^{}g_M(ϵ^{})f\left(\frac{ϵ^{}ϵ}{\gamma }\right).$$
(16)
It has been shown in Ref. that the exact smoothed level density $`g(ϵ)`$ can be reproduced by $`\stackrel{~}{g}_M(ϵ)`$ in the limit of large $`M`$.
In this work we employ the (stretched) harmonic oscillator basis. As discussed below, thirty oscillator shells are sufficient to guarantee the convergence of results. In the standard (old) method, one takes much fewer states (12-14 oscillator shells) and the second term in Eq.(15) is not subtracted. Consequently, the results are not stable as $`M`$ is varied.
### D Generalized Shell-Correction Method
In the standard Strutinsky smoothing method, when applied to finite potentials, it is difficult to meet the plateau condition (9) . The more detailed study of Ref. demonstrated that the plateau condition can seldom be satisfied even if the particle continuum is properly accounted for. Fortunately, it is possible to replace the standard plateau condition with a new requirement which yields unambigous shell-correction values . By comparing the smoothed Strutinsky level densities with those obtained in the semiclassical Wigner-Kirkwood method, it was found that they are in good agreement, apart from the low and the high ends of the spectra. In the intermediate energy region, the average level density shows linear dependence on $`ϵ`$. (The linearity of the semiclassical level density for heavy nuclei was noticed already in Ref. .) Guided by this observation, the shell-correction method was generalized by replacing the plateau condition with the requirement that in an energy interval $`[ϵ_l,ϵ_u]`$ which is wider than the average distance between neighboring major shells,
$$ϵ_uϵ_l=1.5\mathrm{}\omega _0,$$
(17)
the deviation of $`g(ϵ,\gamma ,p)`$ from linearity should be minimal . In practice, one has to minimize the deviation
$$\chi ^2(\gamma ,p)=_{ϵ_l}^{ϵ_u}\left[g(ϵ,\gamma ,p)abϵ\right]^2𝑑ϵ,$$
(18)
where the parameters $`a`$ and $`b`$ are uniquely determined for each value of $`\gamma `$ and $`p`$ by the method of least squares.
Figure 1 displays $`\chi ^2`$ as a function of $`\gamma `$ and $`p`$ for the neutrons in the spherical superheavy nucleus $`Z`$=114, $`N`$=184. It is seen that $`\chi ^2`$ has two minima for each $`p`$ value, and there is a clear correlation between $`p`$ and $`\gamma `$ (minima are shifted to larger values of $`\gamma `$ with increasing $`p`$). The position of the first minimum is between 1.12 and 1.48 $`\mathrm{}\omega _0`$ (i.e., between 6.6 MeV and 8.6 MeV). The second minimum appears at larger $`\gamma `$ values, namely between 7.7 MeV and 12 MeV. As demonstrated in Ref. , the average level densities corresponding to the first minimum of $`\chi ^2`$ practically do not depend on $`p`$ in the negative energy region. This is also valid for the second minimum in $`\chi ^2`$. However, a difference can be seen if one compares average level densities calculated at different minima. The inset of Fig. 1 shows $`\stackrel{~}{g}(ϵ)`$ for both minima at $`p`$=10. The level density corresponding to the lower-$`\gamma `$ minimum preserves its linearity for a wider range of energies, and the linearity is best fulfilled in the energy region which is midway at the bottom and the top of the potential. Therefore, in our calculations \[Eq. (18)\] we fixed the energy interval so that it is centered around the half of the energy of the lowest single-particle level. Far from this central region $`g(ϵ)`$ varies rapidly and this forces the smoothed level density to oscillate. These oscillations can be viewed as end effects, and they are largely independent of the shell structure. For example, for the harmonic oscillator potential whose spectrum has no natural upper bound, these oscillations occur around the bottom of the potential well. For the Woods-Saxon potential, additional oscillations occur around the threshold energy.
It is well known that the realistic value of the smoothing parameter $`\gamma `$ has to lie in a certain energy interval . The value of $`\gamma `$ should be large enough to wipe out shell effects in the energy range of a typical distance between shells, but it should not be much larger to avoid bringing the threshold oscillations down to lower energies. To prevent this, in our calculations we always use the $`\gamma `$ values corresponding to the first minimum of $`\chi ^2`$. For the case shown in Fig. 1, shell correction changes by 0.3 MeV if one uses the value of $`\gamma `$ at the second minimum instead. While in this case the change in shell correction is well within the uncertainty of the model, the difference is more pronounced for lighter nuclei. For the very light nuclei, the end effects dominate the energy dependence of $`\stackrel{~}{g}(ϵ)`$ in the whole energy range. Consequently, the Strutinsky smoothing method cannot be meaningfully applied to these systems.
## III Results
### A Model Parameters
In the calculations presented in this work, we have used the average, axially deformed Woods-Saxon (WS) potential of Ref. , which contains the central part, the spin-orbit term, and the Coulomb potential for protons. The potential depends on a set of deformation parameters, $`\beta _\lambda `$, defining the nuclear surface:
$$R(\theta ;𝜷)=C(𝜷)r_{}A^{1/3}\left[1+\underset{\lambda }{}\beta _\lambda Y_{\lambda 0}(\theta )\right],$$
(19)
where the coefficient $`C`$ assures that the total volume enclosed by the surface (19) is conserved. The Coulomb potential has been assumed to be that of the charge $`(Z1)e`$ uniformly distributed within the deformed nuclear surface. We employed the set of WS parameters introduced in Ref. . For details pertaining to the WS model, see Refs. .
The deformed WS Hamiltonian was diagonalized in the deformed harmonic oscillator basis using the computer code of Ref. . For the diagonalization we took all the (stretched) oscillator states having the principal quantum number less or equal than $`N_{\mathrm{max}}`$. (In short, we took $`N_{\mathrm{max}}`$ deformed shells.) The diagonalization of $`\widehat{H}_0`$ was carried out in precisely the same basis.
When adopting the new scheme, the important question is how many harmonic oscillator shells are needed in order to reproduce the exact results of Ref. . Naturally, the value of $`N_{\mathrm{max}}`$ depends on the size and the shape of the potential to be diagonalized, and also on the oscillator frequency $`\mathrm{}\omega =\eta \mathrm{}\omega _0`$ (the default value of $`\eta `$ is 1.2).
The convergence of $`\delta E_{\mathrm{shell}}`$ for neutrons as a function of $`N_{\mathrm{max}}`$ is illustrated in Fig. 2 for <sup>132</sup>Sn and <sup>154</sup>Sn (which is weakly bound). One can see that for both nuclei the shell correction obtained in the new procedure quickly converges to the exact value, and at $`N_{\mathrm{max}}`$=30 the agreement is very satisfactory. Therefore, in the following, we shall use 30 oscillator shells when applying the new method.
As expected, the results of calculations using the standard method do not stabilize with $`N_{\mathrm{max}}`$. However, for the recommended values of $`N_{\mathrm{max}}`$=12 and 14, the shell correction produced with the old method differs from the exact value by less than 1 MeV. However this apparent agreement seems to be accidental. We display in Fig. 2c and 2d, respectively, the single-particle energy and the smoothed single-particle energy for <sup>154</sup>Sn as a function of $`N_{\mathrm{max}}`$. One can notice that at $`N_{\mathrm{max}}`$=12 the total single-particle energy differs from the exact value by about 8 MeV. Since the corresponding smoothed single-particle energy is also shifted by about 8.5 MeV, the resulting shell correction differs only by 0.8 MeV from the exact value. Consequently, an acceptable agreement for $`\delta E_{\mathrm{shell}}`$ comes as a result of cancellation between two large numbers, each subject to large errors. As one approaches the neutron drip line, the accurate calculation of single-particle energies requires a rather high number of shells and/or a basis optimization with respect to the parameter $`\eta `$ which determines the oscillator length. This is illustrated in Fig. 3 which shows the convergence of the total neutron single-particle energy for <sup>120</sup>Zr at a large quadrupole deformation $`\beta _2`$=0.6. Clearly, for a weakly bound and deformed nucleus one needs at least $`N_{\mathrm{max}}`$=30 oscillator shells to reach the convergence with the standard value of $`\eta `$. Of course, by increasing the oscillator length, i.e., by choosing a smaller value of $`\eta `$, one can improve the convergence significantly for a system with a spatially extended density. In this example, one can arrive at a reasonably accurate value of $`E_{\mathrm{s}.\mathrm{p}.}`$ by using $`N_{\mathrm{max}}`$=20 and $`\eta `$=0.8. Another practical way of improving the convergence of single-particle energies is to use the modified oscillator basis obtained by means of the local scaling transformation .
Figure 4 shows the neutron smoothed level density for <sup>132</sup>Sn, a relatively well bound nucleus, calculated with different methods. The new method with $`N_{\mathrm{max}}`$=30 describes very well $`\stackrel{~}{g}(ϵ)`$ in the whole region of negative energies. This proves that the Green’s function approach can be used with confidence, even for weakly bound systems. On the other hand, the average level density obtained with the standard method never stops increasing, and its deviation from the exact result shows up already at $`ϵ`$=–18 MeV. The result displayed in Fig. 4 demonstrates that even for well-bound nuclei, the shell corrections calculated using the old method are prone to significant errors.
### B Deformation Effects
In order to investigate the deformation dependence of shell corrections, we performed calculations for <sup>100,110,120</sup>Zr as a function of $`\beta _2`$ (other deformation parameters were assumed to be zero). In the Green’s function variant, the generalized plateau condition was used; the resulting values of $`\gamma `$ were also employed in the standard Strutinsky calculations. The results are shown in Fig. 5. As expected, the most pronounced difference between the results of the two methods is for the weakly bound nucleus, <sup>120</sup>Zr. This difference does depend on $`\beta _2`$; part of it can be attributed to the deformation dependence of the smoothing width. (It should be noted that in the deformed calculations of Ref. $`\gamma `$ was assumed to be constant.)
Since the effect of the particle continuum on $`\delta E_{\mathrm{shell}}`$ should be less pronounced for systems that are bound better, one would expect the two methods to yield similar results for lighter Zr isotopes. However, as seen in Fig. 4, the difference between both methods is negligible only at very low energies, $`ϵ`$$`<`$18 MeV. For higher values of $`ϵ`$ (or Fermi level $`\stackrel{~}{\lambda }`$), the difference between the smoothed level densities is not negligible and it is not even a monotonous function of $`\stackrel{~}{\lambda }`$. As a result, one can notice in Fig. 5 that for <sup>110</sup>Zr the results of both methods are very close while in a better bound nucleus of <sup>100</sup>Zr they differ more. It is interesting to note that for <sup>100</sup>Zr and <sup>110</sup>Zr the deformation dependence of $`\delta E_{\mathrm{shell}}`$ is very similar in both variants of calculations. This is in accord with the observation made in Ref. that the difference between the shell corrections obtained in the standard method and the semiclassical approach depends rather weakly on deformation.
As discussed earlier, two main sources of the error of the standard smoothing procedure are (i) the error in determining the total single-particle energy, and (ii) the uncertainty of smoothing that influences the value of the smooth single-particle energy. Since (i) and (ii) are not independent, a large cancellation takes place which might reduce the total error. However, it is a priori difficult to predict how large this difference is and how strongly it would affect the predicted position of the drip line. In order to shed some light on this question, we carried out a comparison between shell corrections calculated with the two methods for the spherical Sn and $`Z`$=114 isotopes, and for the deformed Zr and Er isotopes. (In the latter case, we fixed deformation at $`\beta _2`$=0.2.) These nuclei represented medium-mass and heavy nuclei where the generalized shell-correction method can be applied. The results are presented in Fig. 6. Except for well-bound Zr and $`Z`$=114 isotopes, the difference between the shell corrections calculated using the old and the new methods are on the order of MeV, and, except for the Sn isotopes, it is rather large when approaching $`\stackrel{~}{\lambda }=0`$ (drip line). Although our calculations do not aim at determining the actual position of the drip line, they give a reasonably good estimate for the uncertainty of the old procedure. If one identifies the drip line with the neutron number where $`\lambda `$ becomes positive (this assumption is usually violated in actual calculations because of the lack of self-consistency between the microscopic and macroscopic parts of the energy formula ), our limited calculations suggest that the drip-line predictions by the standard method are prone to severe uncontrolled uncertainties.
### C Modification of the Green’s Function Method to the Proton Case
In the presence of the long-range Coulomb potential, the free Hamiltonian appearing in Eq. (14) has to be modified. Indeed, for the protons, the asymptotic behavior of the scattering states is that of the Coulomb functions, not plane waves. Therefore, in this case, for the free Hamiltonian we take
$$\widehat{H}_0=\widehat{T}+V_{\mathrm{Coul}},$$
(20)
with $`V_{\mathrm{Coul}}`$ being the Coulomb potential. The role of the Coulomb term is to effectively push the continuum up in energy to the top of the Coulomb barrier. The results are insensitive to the radius of the Coulomb potential in the free Hamiltonian. As a matter of fact, even a point Coulomb potential can be used in Eq. (20) .
Figure 7 shows the smoothed proton level density in the proton-rich nucleus <sup>180</sup>Pb. One can see that already with a rather low value of $`N_{\mathrm{max}}`$=19 the new method reproduces the exact smoothed level density in the whole range of negative energies, and $`N_{\mathrm{max}}`$=30 gives an excellent agreement with the exact result. The reason that relatively low values of $`N_{\mathrm{max}}`$ are sufficient in the proton case is that even slightly unbound states (narrow proton resonances) are well localized due to the confining effect of the Coulomb barrier.
## IV Conclusions
In this work, we employed the Green’s function oscillator expansion method to calculations of shell corrections. For spherical nuclei, the new method has proved to be a fast and very accurate approximation to the exact procedure. It also allows for a straightforward generalization to deformed shapes. In essence, the method is based on two simultaneous diagonalizations in a large oscillator basis. The first diagonalization involves the actual one-body Hamiltonian while the other one is carried out for the free Hamiltonian representing the particle gas whose contribution to the level density should be subtracted. For the neutrons, the free Hamiltonian is given by the kinetic energy operator, while for the protons it also includes the Coulomb potential. In practice, the space of 30 (stretched) oscillator shells is sufficient to guarantee the stability of results. This relatively large (but still tractable) space is necessary not only for the proper treatment of the free gas but also for the accurate calculations the total single-particle energy.
As demonstrated in our study, the use of the standard smoothing procedure can lead to serious deviations when extrapolating off beta stability. In particular, the particle drip lines predicted in the traditional approach can be very uncertain. (The systematic error in $`\delta E_{\mathrm{shell}}`$, due to the particle continuum, can be as large as several MeV at the neutron drip line.) According to our calculations, the error on $`\delta E_{\mathrm{shell}}`$ depends weakly on deformation in most cases. It is only for the weakly bound nuclei that the difference between the old and new methods exhibits a sizable deformation dependence.
There is no simple “fix” that would cure the deficiencies of the standard Strutinsky procedure when applied to finite-depth potentials. One does need the large basis in order to guarantee the stability of $`E_{\mathrm{s}.\mathrm{p}.}`$. On the other hand, at these large values of $`N_{\mathrm{max}}`$, the smooth single-particle energy becomes unreliable due to the unphysical increase of the quasi-bound levels around the threshold. We believe that the new Green’s function method, together with the generalized plateau condition, is a very useful tool that should be employed in future global calculations of nuclear masses in the framework of the one-body (macroscopic-microscopic) description and in level-density calculations for spherical and deformed nuclei. Of course, the new procedure does not remove the generic problem of the lack of the self-consistency condition between the microscopic and macroscopic Fermi energies . Recently, the Green’s function method, based on self-consistent potentials obtained in Hartree-Fock and relativistic mean-field calculations, was used to extract shell corrections in the spherical superheavy nuclei . Although these calculations were not done using the oscillator basis expansion method but directly in the coordinate space, their main principle is the same as that discussed in this paper.
###### Acknowledgements.
This research was supported in part by the Hungarian National Research Fund (OTKA T026244 and T029003) and the U.S. Department of Energy under Contract Nos. DE-FG02-96ER40963 (University of Tennessee), DE-FG05-87ER40361 (Joint Institute for Heavy Ion Research), and DE-AC05-96OR22464 with Lockheed Martin Energy Research Corp. (Oak Ridge National Laboratory). |
warning/0001/gr-qc0001070.html | ar5iv | text | # Coordinate singularities in harmonically-sliced cosmologies
## I Introduction
In the 3+1 formulation of general relativity, the freedom in choosing a coordinate system through which to describe a spacetime is replaced by freedom in choosing two gauge quantities: the lapse function $`N`$, which controls the slicing of the spacetime into a foliation of spatial hypersurfaces, and the shift vector $`N^i`$, which defines a set of reference world lines threading those spatial slices. Numerical simulations based on the 3+1 formulation are required to prescribe a method for evaluating the lapse and the shift in terms of the geometrical quantities (the intrinsic metric $`h_{ij}`$ and the extrinsic curvature $`K_{ij}`$) and the matter fields (the density $`\rho `$, momentum $`J^i`$, and stress $`S^{ij}`$) that are known on each spatial slice. The use of a harmonic slicing condition to determine the lapse function is appealing for both practical and theoretical reasons, and the incorporation of such a slicing condition in numerical codes is becoming increasingly common. However, recent work (Alcubierre and Massó , Geyer and Herold ) has found evidence that the use of harmonic slicing can produce pathological behaviour in numerically evolved spacetimes.
This paper addresses the question as to how suitable harmonic slicing is for work in numerical relativity. In section II the main features of the harmonic slicing condition are reviewed. Following an approach similar to that of Geyer and Herold the behaviour of ‘planar’ harmonic slicings of the Minkowski spacetime is analysed in section III; it is demonstrated there that such slicings always cover the whole of the spacetime. A similar statement does not however hold true for the Kasner spacetime. Results derived in section IV show that coordinate singularities of the type found by Alcubierre appear in harmonic slicings of the Kasner spacetime under reasonably general conditions. Section V shows that, furthermore, the results for the Kasner spacetime carry over directly to harmonic slicings of the more general class of Gowdy $`T^3`$ spacetimes. In section VI results are presented from numerical simulations of the Kasner spacetime which use harmonic slicing. The coordinate singularities predicted by the analysis of earlier sections are indeed encountered, and the behaviour of the numerical solutions at times just prior to the formation of coordinate singularities is examined. Section VII concludes the discussion by considering the implications that these results have for the use of harmonic slicing in numerical relativity.
## II Harmonic Slicing in Numerical Relativity
When considering Einstein’s equations as a 3+1 evolution system, a spacetime is described in terms of a foliation $`\{\mathrm{\Sigma }_t:t\}`$ of spatial slices, which defines, in effect, a time coordinate $`t`$ on the spacetime. (The details of this idea are discussed by York , and familiarity with that material is assumed.) The manner in which a spacetime is sliced is a very important issue when the 3+1 formulation is used as a basis for performing numerical simulations. This section recalls some basic ideas on slicing conditions in numerical relativity and, in particular, describes how harmonic slicings of spacetimes are constructed.
Traditionally, work in numerical relativity has been based on maximal or constant mean curvature (CMC) slicing conditions, initially developed for this purpose by Eardley, Smarr and York , among others. A spatial slice has constant mean curvature if $`K`$, the trace of the extrinsic curvature tensor, takes a constant value on that slice. If values for $`K`$ are specified across a range of slices as a function of the time coordinate $`t`$, then the 3+1 evolution equation for the extrinsic curvature yields an elliptic equation determining the value of the lapse function $`N`$ on each slice:
$$\mathrm{\Delta }NN[K_{ij}K^{ij}+4\pi (\rho +S)]=K^{}(t),$$
(1)
where $`\mathrm{\Delta }=h^{ij}_i_j`$. If the function $`K(t)`$ is identically zero then the spacetime is maximally sliced. Maximal slicing has been found to work well in numerical simulations of asymptotically flat spacetimes, while for closed cosmologies (in which at most one maximal slice can exist) the more general CMC slicing condition is appropriate. Based on their behaviour in simple examples and their success in numerical simulations, it is believed that the maximal and CMC slicing conditions will in general produce foliations which cover most, if not all, of a spacetime being investigated, and which at the same time avoid getting too close to any singularities that may form in that spacetime.
Interest in the use of harmonic slicing has arisen in recent years because of its connection with work that has been done in reformulating Einstein’s equations as an explicitly hyperbolic system (see the review by Reula). To date, all of the known 3+1 formulations of general relativity as a strongly hyperbolic evolution system with only physically relevant characteristic speeds determine the lapse through some form of the harmonic slicing condition.
Harmonic slicing contrasts in several basic ways with maximal and CMC slicing. The lapse function $`N`$ in a harmonically-sliced spacetime is determined either through an algebraic condition or a simple evolution equation, and not through an elliptic condition as in equation (1). This ‘local’ specification of the lapse is a great advantage from the point of view of implementing a numerical code to evolve solutions to Einstein’s equations since elliptic partial differential equations are computationally very expensive to solve. The harmonic slicing condition also differs from the maximal and CMC slicing conditions in that it specifies a relationship between two adjacent slices in a foliation rather than determining properties of each individual slice: while an isolated spacetime slice can be characterized as maximal or of constant mean curvature, there is no such thing as an individual harmonic slice. One consequence of this is that the choice of lapse function on the initial slice of a foliation is arbitrary if the harmonic slicing condition is used.
To simplify the following discussion, all spacetime foliations considered in this paper are assumed to have shift vectors $`N^i`$ which are identically zero. Since the shift vector controls only the positioning of spatial coordinates on slices and not how the slices themselves are arranged, this assumption does not limit the generality of the results derived here on the appearance of coordinate singularities.
The harmonic slicing condition for determining the lapse $`N`$ can be expressed in either algebraic form,
$$N=Q(t,x^k)\sqrt{deth_{ij}},$$
(2)
or as an evolution equation,
$$\frac{N}{t}=N(\mathrm{ln}Q)_{,t}N^2K,$$
(3)
where the slicing density $`Q(t,x^i)`$ is an arbitrary (positive) function of the foliation coordinates which does not depend on any evolved variables. If the slicing density is independent of the foliation time coordinate $`t`$, then the harmonic slicing is described as simple, and in this case a choice of the value of the lapse $`N`$ on an initial slice completely determines the value of the slicing density $`Q`$. The alternative situation is described as generalized harmonic slicing, and it is clear from equation (2) that any spacetime foliation can be constructed using generalized harmonic slicing and a particular form for the slicing density $`Q`$. If the slicing density has only a ‘separable’ dependence on the time coordinate of the form $`Q(t,x^i)=f(t)\overline{Q}(x^i)`$, then the resultant foliation has the same slices as the foliation produced by using $`\overline{Q}(x^i)`$ as the slicing density, but with the slices labelled by a different time coordinate. In general it is not obvious how a useful slicing density $`Q`$ which has a non-trivial dependence on the time coordinate $`t`$ may be chosen, and so the simple form of harmonic slicing is most often used in practice.
The main question addressed in this paper is that of how suitable the harmonic slicing condition is for numerical work. Several authors have already considered this question from various different viewpoints. Bona and Massó have shown that the (simple) harmonic slicing condition can be used to foliate several standard spacetimes, and also that ‘focusing singularities’ are avoided by the slicing condition in much the same way as they are in maximally-sliced spacetimes, with the slices of the foliation not reaching the singularity in a finite coordinate time. (It should be noted that for a harmonically-sliced spacetime a focusing singularity is essentially a point at which the lapse $`N`$ becomes zero—this is in contrast to the behaviour found at the ‘gauge pathologies’ described below.) Cook and Scheel have investigated the construction of well-behaved harmonic foliations for Kerr-Newmann black hole spacetimes.
The work of Alcubierre and Massó is of particular relevance to the present discussion. They have shown that ‘gauge pathologies’ (described as ‘coordinate shocks’ in the earlier paper) can occur in numerical simulations based on hyperbolic formulations of Einstein’s equations which use harmonic slicing. (In fact, a range of gauge conditions are considered by the authors, with simple harmonic slicing—the only gauge choice of interest here—corresponding to the special case $`f=1`$ in their formulation. None of the alternative gauge choices they use correspond to the generalized harmonic slicing condition.) These gauge pathologies manifest as a loss of continuity at points in the evolved solution with, in particular, large spikes appearing in the lapse. After the time at which the pathologies appear the numerical solution no longer converges at the expected order. Alcubierre and Massó explain the appearance of gauge pathologies in terms of nonlinear behaviour in the hyperbolic evolution equations. However this fails to adequately answer questions about how common the gauge pathologies are, what happens to the foliation at the points where pathologies appear, and what approaches can be used to prevent the pathologies from occurring.
Simple harmonic slicings for the Schwarzschild and Oppenheimer-Snyder spacetimes have been investigated by Geyer and Herold . The approach they use is based on an alternative representation of equations (2) and (3) in the simple harmonic case: if $`T`$ is a scalar function on a spacetime, the level surfaces of which represent the spatial slices of a foliation $`\{\mathrm{\Sigma }_T\}`$, then the spacetime will be harmonically sliced (in the simple sense) if
$$\mathrm{}Tg^{\mu \nu }_\mu _\nu T=0.$$
(4)
By numerically integrating equation (4) with respect to known background metrics, Geyer and Herold construct simple harmonic slicings which they compare to maximal slicings of the spacetimes. Some properties of a slicing can be determined straightforwardly from its time function $`T`$, and in particular the lapse $`N`$ can be evaluated through the equation
$$g^{\mu \nu }(_\mu T)(_\nu T)=1/N^2,$$
(5)
where it is clear that the vector field normal to the foliation must remain timelike if the lapse is to have a positive real value. In fact, Geyer and Herold find that for simple harmonic slicings of the Oppenheimer-Snyder spacetime, foliations which are initially timelike can at later times become null or spacelike, with the development of the foliation thus terminating at what it seems appropriate to call a coordinate singularity. The lapse becomes infinite as these singular points are reached, and this is consistent with the behaviour found at the gauge pathologies of Alcubierre and Massó.
In what follows, Geyer and Herold’s approach is applied to Minkowski, Kasner and Gowdy spacetimes, with the intention being to gain a better understanding of the circumstances under which coordinate singularities appear, and in particular to determine whether they are a rare or a common feature of harmonic slicings.
## III Harmonic Slicings of the Minkowski Spacetime
In the present work, equations (4) and (5) are used to investigate the formation of coordinate singularities in (simple) harmonic slicings of cosmological models. The following approach is employed. The metric $`g_{\mu \nu }`$ of the spacetime being investigated is assumed to be known with respect to a coordinate system $`(t,x,y,z)`$. An alternative foliation of the spacetime is constructed by taking as an initial slice one of the constant time hypersurfaces of the background coordinate system: the foliation time coordinate $`T(t,x,y,z)`$ is given an initial value
$$T(t_0,x,y,z)=t_0,$$
(6)
for some value $`t_0`$ of the coordinate $`t`$. (It is assumed here that the hypersurface $`t=t_0`$ is spacelike.) The lapse function of the foliation can be specified arbitrarily on the initial slice and this determines via equation (5) the value of the first derivative of $`T`$ away from that slice:
$$T_{,t}(t_0,x,y,z)=\frac{1}{N_0(x,y,z)\sqrt{g^{00}(t_0,x,y,z)}}.$$
(7)
Equations (6) and (7) provide initial data for the wave equation (4), and the solution $`T`$ describes a new simple harmonic slicing of the spacetime.
As an example of how equation (4) can be used to find coordinate singularities, consider harmonic foliations of the Minkowski spacetime, written in standard coordinates as
$`g_{\mu \nu }dx^\mu dx^\nu =dt^2+dx^2+dy^2+dz^2.`$
Suppose that a ‘planar’ slicing of the spacetime is constructed such that the time function $`T`$ depends only on the Minkowski coordinates $`t`$ and $`x`$. Then equation (4) takes the form of the one-dimensional wave equation
$`T_{,tt}=T_{,xx},`$
which has the general solution
$`T(t,x)=f(t+x)+h(tx),`$
for arbitrary functions $`f(u)`$ and $`h(u)`$. The lapse associated with this time function can be found from equation (5):
$`1/N^2`$ $`=`$ $`T_{,x}{}_{}{}^{2}T_{,t}^2`$
$`=`$ $`4f^{}(t+x)h^{}(tx),`$
and the lapse will be well behaved as long as the function $`f^{}(t+x)h^{}(tx)`$ is positive.
For most problems the next step in the analysis would be to use equations (6) and (7) to specify a value for the lapse on an initial slice of the foliation. However the present case is sufficiently simple that the appearance of coordinate singularities can be studied without needing to specify an initial slice. If the lapse becomes infinite at a point in the foliation then the function $`f^{}(t+x)h^{}(tx)`$ must be zero there, and it is clear that the function must then be zero at all points along a line $`t+x=\text{constant}`$ or $`tx=\text{constant}`$. Consequently, any spacelike slice of that foliation must include a point in it at which the lapse is infinite. It follows that for simple harmonic slicings of the ‘planar’ Minkowski spacetime no coordinate singularities will be present, and the foliation will cover the whole of the spacetime, provided that the initial slice of the foliation is everywhere spacelike. (This result is consistent with the numerical simulations of Minkowski spacetime reported on by Alcubierre . No gauge pathologies are discovered for ‘planar’ harmonic slicings—the $`f=1`$ case—of flat spacetime, although they are found in the spherically symmetric case.)
In the following sections a similar analysis is performed for more complicated spacetimes, and coordinate singularities are found to be much more common than the above result for the Minkowski spacetime would suggest.
## IV Harmonic Slicings of the Kasner Spacetime
In this section, equation (4) is used to construct simple harmonic foliations of the Kasner spacetime, with the main point of interest being the question of whether or not coordinate singularities (gauge pathologies in the terminology of Alcubierre and Massó ) form in the slicings. Later, in section VI, these results are compared to numerical simulations of the spacetime.
The present work uses the axisymmetric Kasner model described by the metric
$$ds^2=t^{1/2}(dt^2+dx^2)+t(dy^2+dz^2),$$
(8)
for which a cosmological singularity is present at time $`t=0`$. For convenience in the following analysis (and to strengthen the connection with the Gowdy cosmologies examined in section V) a three-torus topology is imposed on the Kasner spacetime: periodicity is assumed in each of the coordinates $`x`$, $`y`$ and $`z`$ over the range $`[0,2\pi )`$.
As in the simple example for the Minkowski spacetime given in the previous section, here ‘planar’ slicings of the Kasner spacetime are considered for which the time function $`T`$ depends only on the coordinates $`t`$ and $`x`$. A straightforward calculation then shows that for the metric (8), equation (4) for the time function of a harmonically-sliced foliation is equivalent to
$$T_{,tt}T_{,xx}+t^1T_{,t}=0.$$
(9)
This equation can be solved by decomposing $`T`$ spatially as a sum of Fourier modes (recalling that, since the model is spatially closed, the function $`T`$ must be periodic in the coordinate $`x`$) and the general solution is found to be
$`T`$ $`=`$ $`a_0\mathrm{ln}t+b_0`$ (11)
$`+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[Z_{+n}(t)\mathrm{cos}(nx)+Z_n(t)\mathrm{sin}(nx)\right],`$
with
$`Z_{\pm n}(t)=a_{\pm n}J_0(|n|t)+b_{\pm n}Y_0(|n|t)\text{ for }n0,`$
where $`a_{\pm n}`$ and $`b_{\pm n}`$ are constants, and $`J_\nu `$ and $`Y_\nu `$ are Bessel functions of the first and second kinds .
Equation (5) for the lapse function $`N`$ of the foliation can be written in the present case as
$$\sqrt{t}(T_{,t}{}_{}{}^{2}T_{,x}{}_{}{}^{2})=1/N^2,$$
(12)
and the requirement that the foliation be well behaved can then be expressed as
$$T_{,t}>|T_{,x}|,$$
(13)
where it is assumed that the orientation of the time function $`T`$ of the foliation is the same as that of the Kasner time $`t`$, and thus that $`T_{,t}`$ is positive. If the time function $`T`$ of equation (11) fails to satisfy the condition (13) at a point, then the harmonic slicing must have a coordinate singularity there.
The constants $`a_{\pm n}`$ and $`b_{\pm n}`$ in equation (11) can be determined by specifying a value $`N_0(x)`$ for the lapse on an initial hypersurface $`t=t_0`$ as in equations (6) and (7). In fact, rather than choosing the initial value for the lapse directly, it is more convenient to choose an initial value $`I(x)`$ for the time derivative of $`T`$, with this being specified as a Fourier sum,
$`T_{,t}(t_0,x)`$ $`=`$ $`I(x)`$ (14)
$`=`$ $`k_0+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left[k_{+n}\mathrm{cos}(nx)+k_n\mathrm{sin}(nx)\right],`$ (15)
where the constants $`k_{\pm n}`$ must be such that $`I(x)`$ is everywhere positive. The value of the lapse on the initial slice can then be determined through equation (7):
$`N_0(x)={\displaystyle \frac{1}{t_0{}_{}{}^{1/4}I(x)}}.`$
The particular time function $`T`$ of equation (11) which fits the initial data of equations (6) and (14) is found to be
$$T=t_0+k_0t_0\mathrm{ln}(t/t_0)\frac{\pi t_0}{2}\underset{n=1}{\overset{\mathrm{}}{}}\left[J_0(nt)Y_0(nt_0)Y_0(nt)J_0(nt_0)\right]\left[k_{+n}\mathrm{cos}(nx)+k_n\mathrm{sin}(nx)\right].$$
(16)
The simplest harmonic foliations of the Kasner metric (8) are the homogeneous ones for which the time function $`T`$ is independent of the coordinate $`x`$. Then, to within a linear rescaling, the function $`T`$ equals $`\mathrm{ln}t`$, and so the cosmological singularity in the spacetime occurs in the harmonic foliation only in the limit as $`T`$ tends to negative infinity. Clearly the foliation covers the entire Kasner spacetime and contains no coordinate singularities.
The behaviour of homogeneous harmonic slicings is, however, not representative of the general case. As an example, figure 1 shows the foliation which results from setting the lapse on the initial slice to be
$$N_0(x)=\frac{1}{1+\frac{1}{2}\mathrm{sin}x}\text{for }t_0=1,$$
(17)
which corresponds to the choice of non-zero coefficients $`k_0=1`$ and $`k_1=1/2`$ in equation (14). When equation (12) is evaluated to determine the lapse on this foliation, it is found that in some spacetime regions to the future of the initial slice the quantity $`1/N^2`$ is non-positive; figure 2, which plots $`1/N^2`$, is filled in where this happens. Coordinate singularities must appear in any slices of the foliation which intersect these regions, and the development of the foliation cannot proceed beyond a time $`T=T_{\text{sing}}2.45`$ with the lapse $`N`$ tending to infinity at points as the limiting slice is approached. In contrast, if the foliation is extended backwards in time from the initial slice towards the cosmological singularity, then the lapse appears to be well behaved on all of the slices.
It can in fact be shown that, in general, harmonic foliations of the Kasner spacetime always develop coordinate singularities at sufficiently late (future) times. If the time derivative of the function in equation (16) is evaluated at a fixed spatial coordinate $`x=x_0`$, then
$`T_{,t}(t,x_0)={\displaystyle \frac{k_0t_0}{t}}+{\displaystyle \frac{\pi t_0}{2}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}nc_n(x_0)\left[J_1(nt)Y_0(nt_0)Y_1(nt)J_0(nt_0)\right],`$
where it is assumed that the value $`x_0`$ has been chosen such that at least one of the values $`c_n`$ is non-zero. For sufficiently large values of the coordinate $`t`$, the approximations (9.2.5) and (9.2.6) of reference can be applied to the Bessel functions in this equation. The result is that
$`T_{,t}(t,x_0){\displaystyle \frac{k_0t_0}{t}}+\sqrt{{\displaystyle \frac{\pi }{2t}}}t_0{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\sqrt{n}c_n\left[Y_0(nt_0)\mathrm{cos}(nt\frac{3}{4}\pi )J_0(nt_0)\mathrm{sin}(nt\frac{3}{4}\pi )\right]={\displaystyle \frac{k_0t_0}{t}}+{\displaystyle \frac{\mathrm{\Theta }(t)}{\sqrt{t}}},`$
where $`\mathrm{\Theta }(t)`$ is a periodic function which takes both positive and negative values. It follows that, for some sufficiently large value of the coordinate $`t`$, the value of $`T_{,t}`$ must be negative at a point. Since this violates the condition (13), it must be the case that a coordinate singularity is present in the foliation.
The above argument shows that coordinate singularities must in general appear in harmonic slicings of an expanding Kasner cosmology but gives no indication of how much of the spacetime a slicing will cover before it becomes pathological. To investigate this, consider the simple case of a time function $`T`$ from equation (16) which, like the foliation produced by the initial data (17), contains only one non-trivial mode of amplitude $`k`$:
$`T`$ $`=`$ $`t_0+k_0t_0\mathrm{ln}(t/t_0)`$ (19)
$`k{\displaystyle \frac{\pi t_0}{2}}\left[J_0(nt)Y_0(nt_0)Y_0(nt)J_0(nt_0)\right]\mathrm{sin}(nx),`$
where $`0<k<k_0`$. For a fixed value of the Kasner time $`t`$ the foliation will be well behaved (in that condition (13) is satisfied) if and only if
$`{\displaystyle \frac{k_0t_0}{t}}>{\displaystyle \frac{kn\pi t_0}{2}}\underset{x}{\mathrm{max}}\left\{|A(t)\mathrm{cos}(nx)|B(t)\mathrm{sin}(nx)\right\},`$
where
$`A(t)`$ $`=`$ $`J_0(nt)Y_0(nt_0)Y_0(nt)J_0(nt_0),`$
$`B(t)`$ $`=`$ $`J_1(nt)Y_0(nt_0)Y_1(nt)J_0(nt_0),`$
and this condition is equivalent to
$$F(t)t^2\left[A(t)^2+B(t)^2\right]<\left(\frac{2k_0}{kn\pi }\right)^2.$$
(20)
It is straightforward to show that $`F^{}(t)0`$. Furthermore, if the time $`t`$ is assumed to be large enough for approximations to be applied to the Bessel functions in $`A(t)`$ and $`B(t)`$, then $`F(t)t\mathrm{\Pi }(t)`$ where $`\mathrm{\Pi }(t)`$ is a positive, periodic function. An estimate of the time $`t=t_{\text{sing}}`$ at which a coordinate singularity develops can then be seen to obey
$$t_{\text{sing}}(k_0/k)^2,$$
(21)
and so, for a harmonic slicing which is perturbed from homogeneity by a single mode of amplitude $`k`$, the amount of the spacetime covered by the slicing becomes infinite as $`k`$ tends to zero.
Considering now the behaviour of the slicing in the region of spacetime between the initial slice and the cosmological singularity, it follows from equation (20) that, for the single mode case (19), coordinate singularities never appear. To see this, note that the function $`F(t)`$ must satisfy the inequality at time $`t=t_0`$ since the lapse is required to be well behaved on the initial slice, and since it is an increasing function of time it must then also satisfy the inequality at all earlier times. In the limit as $`t`$ tends to zero (while $`x`$ is constant), the Bessel function behaviour is known from equations (9.1.10–13) of reference , and the time function $`T`$ must have the form
$`T`$ $``$ $`t_0\left[k_0+kJ_0(nt_0)\mathrm{cos}(nx)\right]\mathrm{ln}t`$ (23)
$`+\text{(terms that tend to a constant)}.`$
Since the factor multiplying $`\mathrm{ln}t`$ is always positive, the value of $`T`$ must tend to negative infinity as the cosmological singularity is approached. Thus a single mode harmonic slicing will always cover the whole of the Kasner spacetime to the past of the initial slice.
Returning to the general case in which an arbitrary number of modes are present in the time function $`T`$, the above argument can be used to derive a simple condition on the initial data (14) which ensures that a slicing is well behaved in the region of spacetime between the initial slice and the cosmological singularity. Suppose that equation (16) is split up as
$`T=t_0+{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(T^{(+n)}+T^{(n)}\right),`$
where
$`T^{(+n)}`$ $`=`$ $`k_0^{(+n)}t_0\mathrm{ln}(t/t_0)k_{+n}{\displaystyle \frac{\pi t_0}{2}}\left[J_0(nt)Y_0(nt_0)Y_0(nt)J_0(nt_0)\right]\mathrm{cos}(nx),`$
$`T^{(n)}`$ $`=`$ $`k_0^{(n)}t_0\mathrm{ln}(t/t_0)k_n{\displaystyle \frac{\pi t_0}{2}}\left[J_0(nt)Y_0(nt_0)Y_0(nt)J_0(nt_0)\right]\mathrm{sin}(nx),`$
and the values $`k_0^{(\pm n)}`$ are arbitrary subject to the condition
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}\left(k_0^{(+n)}+k_0^{(n)}\right)=k_0.`$
The condition (13) that the foliation be free of coordinate singularities will then certainly be satisfied if
$`T^{(+n)}{}_{,t}{}^{}>|T^{(+n)}{}_{,x}{}^{}|\text{and}T^{(n)}{}_{,t}{}^{}>|T^{(n)}{}_{,x}{}^{}|n,`$
and, as the analysis of the single mode case shows, this will be true for $`0<t<t_0`$ if it is true for $`t=t_0`$. The time derivative of each function $`T^{(\pm n)}`$ must therefore be positive on the initial hypersurface $`t=t_0`$, and for this to be the case the values $`k_0^{(\pm n)}`$ must be chosen such that
$`k_0^{(+n)}>|k_{+n}|\text{and}k_0^{(n)}>|k_n|n.`$
This can always be done if
$$k_0>\underset{n=1}{\overset{\mathrm{}}{}}\left(|k_{+n}|+|k_n|\right),$$
(24)
and so any initial value for the lapse which satisfies this condition must produce a foliation which is free of coordinate singularities to the past of the initial slice. In addition to this, if the behaviour of the time function $`T`$ is considered as $`t`$ tends to zero, then it is straightforward to show (in analogy with equation (23) for the single mode case) that condition (24) ensures that $`T`$ tends to negative infinity for all values of $`x`$ as the cosmological singularity is approached, and hence that the foliation covers all of the spacetime up to the singularity. The condition (24) on the initial lapse is sufficient but not necessary for the harmonic foliation to be well behaved to the past of the initial slice. In fact, experimentation with different choices of initial data (14) suggests that, even when condition (24) is not satisfied, harmonic foliations may never develop coordinate singularities as the cosmological singularity is approached, although no proof (or disproof) of this conjecture has yet been constructed.
## V Harmonic Slicings of Inhomogeneous Cosmologies
To summarize the results of the previous section, foliations of the Kasner spacetime (8) based on the simple harmonic slicing condition (4) have been investigated under the assumptions that one slice of the foliation coincides with a $`t=\text{constant}`$ hypersurface of the original metric, and that the foliation is independent of two of the standard spatial coordinates. It is found that when the Kasner cosmology is expanding, the slices of the foliation must eventually develop coordinate singularities (except when the lapse is initially constant), with this happening at late times for foliations which are initially close to homogeneous. In contrast, when the Kasner cosmology is collapsing, the results suggest that coordinate singularities never develop in the foliation and that the slices extend all the way to the cosmological singularity. (This is certainly true if the lapse on the initial slice is chosen such that the coefficients of equation (14) satisfy the condition (24), and it may in fact be the case for all reasonable choices of initial lapse.)
The behavioural differences found in harmonic foliations according to whether the Kasner spacetime is expanding or collapsing can be explained by considering the relationship between the lapse function $`N`$ and the mean curvature $`K`$ of a slice. In any foliation, the mean curvature provides a measure of the local convergence of world lines running normal to the spatial slices:
$$K=_\mu n^\mu =\mathrm{\pounds }_n\left(\mathrm{ln}\sqrt{deth_{ij}}\right),$$
(25)
where $`n^\mu `$ is the vector normal to the slices, $`\mathrm{\pounds }_n`$ is the Lie derivative along that vector, and $`\sqrt{deth_{ij}}`$ is the spatial volume element. The value of $`K`$ is positive when the foliation world lines are locally converging, and negative when they are expanding. The harmonic slicing condition can be written in terms of the mean curvature of the foliation using equation (3):
$`{\displaystyle \frac{(1/N)}{t}}=K,`$
where the harmonic slicing is assumed to be simple with zero shift. It follows that when the mean curvature $`K`$ is positive (as it may be expected to be for slices of a collapsing cosmology) the value of $`1/N`$ will increase, and so the lapse will approach (but usually not reach) zero. (The ability of the harmonic slicing condition to avoid ‘focusing singularities’ at which the lapse vanishes is discussed by Bona and Massó .) Conversely, when the mean curvature is negative (as it typically will be in an expanding spacetime) the value of $`1/N`$ will decrease, and if it reaches zero then a coordinate singularity of the type investigated in this paper will develop in the foliation.
The obvious extension of the above analysis is to the consideration of harmonic slicings of spacetimes more general than the Kasner model. In fact it turns out to be straightforward to extend the results of section IV to a class of inhomogeneous cosmological models : the unpolarized Gowdy spacetimes on $`T^3`$ (the three-torus). These spacetimes are vacuum models of spatially closed cosmologies with planar symmetry, and can be represented by the metric
$`ds^2`$ $`=`$ $`e^{\lambda /2}t^{1/2}(dt^2+d\theta ^2)`$ (27)
$`+e^Pt\left[d\sigma ^2+2Qd\sigma d\delta +(Q^2+e^{2P})d\delta ^2\right].`$
The spatial coordinates $`\theta `$, $`\sigma `$ and $`\delta `$ range from $`0`$ to $`2\pi `$. A cosmological singularity occurs at time $`t=0`$ and the model expands forever through positive values of $`t`$. The functions $`\lambda `$, $`P`$ and $`Q`$ depend on the coordinates $`\theta `$ and $`\tau `$ only and are required by the topology of the model to be periodic in $`\theta `$. The vacuum Einstein equations for the metric (27) reduce to evolution and constraint equations for the variables $`P`$, $`Q`$ and $`\lambda `$ which cannot in general be solved exactly. Moncrief has shown that no singularities appear in the Gowdy $`T^3`$ spacetimes for $`0<t<+\mathrm{}`$, and that this coordinate range describes the maximal Cauchy development of the models. Thus, no coordinate singularities of the type investigated in this paper can be present in the standard Gowdy slicing.
If the analysis of the previous section for harmonic slicings of the axisymmetric Kasner spacetime (8) is applied to the Gowdy metric (27) (taking the planar symmetry of the slicing to coincide with the symmetry of the metric) then it is found that the time function $`T`$ must satisfy
$`T_{,tt}T_{,\theta \theta }+t^1T_{,t}=0,`$
while the lapse $`N`$ of the associated foliation is defined through
$`e^{\lambda (t,\theta )/2}\sqrt{t}(T_{,t}{}_{}{}^{2}T_{,\theta }{}_{}{}^{2})=1/N^2.`$
Comparing these expressions to equations (9) and (12), it is clear that the Gowdy $`T^3`$ spacetimes admit the same harmonic foliations as the axisymmetric Kasner spacetime, and furthermore that the condition that foliations must satisfy to be free of coordinate singularities is the same. It thus turns out that all of the results summarized above on the behaviour of harmonic slicings of a Kasner cosmology carry over directly to a fairly general class of planar cosmological models. (In addition, it may be noted that the class of Gowdy spacetimes on $`T^3`$ includes as special cases the full set of Kasner spacetimes with the same topology. Thus the restriction of the discussion in section IV to only the axisymmetric Kasner model (8) can be relaxed; the same results hold for all of the Kasner spacetimes.)
## VI Numerical Simulations of the Kasner Spacetime
The results of section IV show that coordinate singularities are a generic feature of harmonic foliations of the (expanding, axisymmetric) Kasner cosmology. It follows then that numerical simulations of that spacetime which use the harmonic slicing condition will be forced to terminate after a finite number of time steps, regardless of their accuracy. The practical details of this are examined in the present section.
Numerical simulations are performed based on initial data for the Kasner spacetime given by the metric (8) on an initial slice at time $`t=t_0=1`$ using the harmonic slicing condition (2) and the initial value (17) for the lapse. The foliation of the Kasner spacetime which results is the one pictured in figure 1, and the numerical simulation thus allows the prediction made in section IV—that coordinate singularities prevent the foliation from developing beyond a time $`T=T_{\text{sing}}2.45`$—to be verified. The slicing density of the foliation is given the inhomogeneous form
$$Q(T,X,Y,Z)=\frac{t_0^{3/4}}{1+\frac{1}{2}\mathrm{sin}X},$$
(28)
where $`(T,X,Y,Z)`$ is the coordinate system of the numerical simulation, and is distinct from the coordinate system $`(t,x,y,z)`$ of the original Kasner metric except on the initial slice which is labelled $`T=t_0`$ and on which $`(X,Y,Z)=(x,y,z)`$. The shift vector $`N^i`$ is taken to be identically zero, and the lack of dependence of the slicing density $`Q`$ on the foliation time $`T`$ means that the harmonic slicing is of the simple type. All of the quantities involved in the numerical simulation are independent of the spatial coordinates $`Y`$ and $`Z`$ so that the evolved spacetime in effect has planar symmetry. The numerical simulation is one dimensional and uses a high-resolution numerical scheme within an adaptive mesh refinement code to evolve a first-order hyperbolic formulation of the vacuum Einstein equations; a complete description of the code is given in reference .
Figure 3 displays the $`(T,X)`$ coordinate system of the numerical simulation as it appears in the $`(t,x)`$ coordinates of the homogeneous background metric (8). This comparison between coordinate systems is made using an algorithm, described in reference , for tracking the positions of the observers at rest in the slices of the foliation, based on the values $`N(T,X)`$ taken by the lapse during the course of the simulation. The observer world lines are plotted up until a simulation time $`T=2.3`$. Although the inaccuracies in the numerical solution which are present at that time are not sufficient to stop the simulation, they do make it increasingly difficult to accurately track the positions of the observers. The $`T=\text{constant}`$ surfaces reconstructed in figure 3 can be seen to be in good agreement with the exact results shown in figure 1.
The analysis of section IV (as presented in figures 1 and 2) shows that at a time $`T=T_{\text{sing}}2.45`$ the lapse $`N`$ must become infinite at points of the evolved foliation. Figure 4 shows the actual value taken by the lapse at the earlier time $`T=2.3`$ in the simulation, and it can be seen that two sharp spikes are already present in the solution. As the evolution progresses these spikes grow in size, and if the evolved solution were exact they would reach infinite heights within a short time. As they grow, the two spikes also narrow and become closer together. The reason for this can be seen from the paths of the coordinate observers plotted in figure 3: in the region $`0<x<\pi `$ where the coordinate singularities eventually form, the world lines normal to the foliation are diverging (with respect to the background coordinates), and so the region is resolved by a diminishing number of grid points. It should be noted that this divergence of observers is not itself the cause of the coordinate singularities; it could in principle be counteracted by an appropriate choice of shift vector for the foliation, but this would not prevent the lapse from becoming infinite. The narrowing and effective coalescence of the spikes is problematic for the simulation since features which are smaller than the grid spacing $`\mathrm{\Delta }X`$ cannot be accurately resolved, and, regardless of the number of grid points used, eventually the numerical solution must fail to accurately model the behaviour of the exact solution. In fact, with the spikes being inadequately resolved, the evolved value of the lapse fails to become infinite (or even the computer representation of this) as the coordinate singularities are reached, and it is possible for the simulation to continue beyond the time at which the slices of the foliation cease to be spacelike in the exact solution. However, the numerical solution has no physical meaning past the points at which coordinate singularities form, and in particular, as noted by Alcubierre and Massó , it will no longer converge as the grid is refined.
The spikes that develop in the lapse as the coordinate singularities are approached are at face value very similar to some features that are seen in numerical solutions for collapsing one- and two-dimensional inhomogeneous cosmologies, as studied by Hern and Stewart . In both cases spiky features appear in the evolved variables which, as the simulation progresses, increase in height and decrease in width until they can no longer be resolved by the numerical grid used in the simulation. The question then arises as to whether the features seen in the inhomogeneous cosmologies have any physical relevance or whether they are, like the features in the Kasner simulation described above, entirely a coordinate effect.
In reference it is concluded that the spiky features found in the inhomogeneous cosmological models are not simply artifacts of the coordinate systems (or the metric variables) on which the simulations are based. Part of the evidence for this comes from an examination of the profiles of the evolved variables with respect to proper distances rather than coordinate distances: the narrowing of the spiky features becomes even more pronounced when proper distances are used. It is interesting to note that the opposite result is observed when a similar analysis is applied to the spikes seen in the lapse in figure 4. If the metric component $`h_{11}`$, which measures the proper distance in the $`X`$-direction of the simulation, is examined for the data plotted in figure 4, it is found that sharp spikes (towards positive infinity) are coincident there with the spikes in the lapse, and hence that the apparently narrow spikes are in fact spread over large physical distances. As figure 5 demonstrates, if instead of being plotted against the coordinate $`X`$, the lapse is shown as a function of the proper distance along the $`X`$-axis, then the spikes no longer appear as notable features of the solution. This approach could prove to be a useful way of distinguishing genuine physical features in an evolved solution from effects caused by coordinate singularities in harmonic slicings.
## VII Discussion
It has been demonstrated in sections IV and V of this paper that ‘planar’ foliations of Kasner and Gowdy $`T^3`$ cosmologies generated using the simple harmonic slicing condition will, except in special cases, always terminate at coordinate singularities without covering the whole of the spacetimes. When these analytic results are considered together with the descriptions of coordinate singularities in specific examples of ‘spherically symmetric’ harmonic foliations given by Alcubierre , and Geyer and Herold , a pessimistic picture of the nature of the harmonic slicing condition emerges: it appears likely that, except in special cases, foliations generated by harmonic slicing will eventually terminate at coordinate singularities. (Seemingly this pathological behaviour can occur whenever the foliation locally undergoes a protracted period of expansion).
The coordinate singularities examined here arise as part of exact solutions for harmonically-sliced spacetimes, and they cannot be avoided simply by employing alternative numerical methods in simulations. This is clearly a drawback to the use of harmonic slicing in numerical work, and by extension to the use of hyperbolic formulations of Einstein’s equations which typically rely on this type of slicing. It should be appreciated however that no known slicing condition is guaranteed to completely cover an arbitrary spacetime, and potential problems with coordinate singularities do not necessarily outweigh the advantages of using hyperbolic formulations in numerical relativity. In section VI the behaviour of a numerically evolved solution is examined as it approaches a coordinate singularity, and it is found that by considering measures of proper distance it may be possible to distinguish the effects of coordinate singularities from genuine physical features of a spacetime. In general though the most reliable approach to identifying the presence of a coordinate singularity in a numerical solution is the one recognized by Alcubierre and Massó : convergence of the solution is lost when a coordinate singularity develops.
Of course, the formation of coordinate singularities has only been discussed here for the simple harmonic slicing condition, and it is possible that the freedom in choosing the time dependence of the slicing density in generalized harmonic slicing could be put to use in controlling the development of the spacetime foliation such that coordinate singularities are avoided. Although when used in the context of a hyperbolic formulation the slicing density is formally required to be independent of the evolved variables, in practice there seems to be no problem in allowing occasional ‘corrections’ to be made to its value in response to the behaviour of the foliation; in effect this amounts to intermittently halting the simulation and choosing a new value for the lapse on the current slice. (An important point here is that changes to the slicing density should be made only at fixed time intervals, rather than after a fixed number of time steps, since otherwise the exact solution being sought will depend on the resolution of the simulation.) This ‘piecewise harmonic’ form of slicing has been used in some preliminary tests employing simple heuristics for making alterations to the value of the slicing density, with the resulting foliations being examined using the world line integration algorithm which was used to produce figure 3. (The algorithm was in fact developed for this purpose.) As yet however no definite conclusions regarding the effectiveness of this approach have been reached.
###### Acknowledgements.
Thanks must go to John Stewart for overseeing this research. The author was supported by a studentship from the Engineering and Physical Sciences Research Council. |
warning/0001/astro-ph0001223.html | ar5iv | text | # 1 INTRODUCTION
## 1 INTRODUCTION
The main goal of this paper is to present a detailed analysis of the following problem: Given a beam structure, a level of noise, a certain partial coverage, and a pixelization, how well can we assign a temperature to each pixel?. In other words, how well can we deconvolve the beam to get appropriate temperatures in the pixels?. Here, temperatures are considered to be appropriate when the resulting maps lead to good physical spectra. Hereafter, any beam reversion leading to right spectra is referred to as a ”S–deconvolution”. In the absence of noise, the possibility of performing a good S–deconvolution essentially depends on the ratio between the beam area and that of the chosen pixel. In practice, S–deconvolution is not feasible for too large values of this ratio; in other words, if we fix the beam, S–deconvolution is not feasible for too small pixel sizes. For a given beam and a certain mathematical method, there is a minimum pixel size allowing S–deconvolution. For values smaller than this minimum, too many pixels can be placed inside the beam and S–deconvolution is not possible. The minimum size corresponding to two S–deconvolution methods has been estimated in various cases. Both methods lead to similar minimum values of the pixel size around $`\theta _{_{FWHM}}/2`$. These values depend on the level of the uncorrelated noise produced by the instruments.
It is worthwhile to emphasize that we are not interested in assigning temperatures to hundreds of pixels located inside the beam. This assignation can be useful in other contexts; however, in our case, the important point is that the spectra contained in the S–deconvolved maps must be similar to the true physical spectra (up to the scales corresponding to the pixel size). Unfortunately, this type of deconvolution requires a moderate number of pixels inside the beam. This number will be estimated below in various cases.
From § 2 to § 4, pixelization is considered in the framework of pure beam reversion, without analyzing particular experiments; however, in § 5, we focus our attention on PLANCK multifrequency observations and, then, pixelization is discussed taking into account both previous conclusions about beam reversion and some physical constraints due to difraction, observational strategy, et cetera.
In Sáez, Holtmann & Smoot (1996) and Sáez & Arnau (1997), the modified power spectrum
$$E_{\mathrm{}}(\sigma )=\frac{32\pi ^3}{(2\mathrm{}+1)^2}_{\alpha _{min}}^{\alpha _{max}}C_\sigma (\alpha )P_{\mathrm{}}(\mathrm{cos}\alpha )\mathrm{sin}\alpha d\alpha $$
(1)
was described. Functions $`P_{\mathrm{}}`$ are the Legendre polinomial normalized as follows: $`P_{\mathrm{}}P_{\mathrm{}^{}}d(\mathrm{cos}\theta )=[(2\mathrm{}+1)/8\pi ^2]\delta _{\mathrm{}\mathrm{}^{}}`$. As explained in those papers, this type of spectrum can be easily found from both theory and maps. Comparisons of the modified spectra obtained from theory with those extracted from simulated or observational maps are appropriate to take into account pixelization, partial coverage and beam features, simultaneously. Sometimes, the estimation and use of the well known $`C_{\mathrm{}}`$ coefficients –although possible– is not the best procedure. In Eq. (1), the effect of pixelization is simulated by the angle $`\alpha _{min}`$, which is the angle separating two neighbouring nodes, while the angle $`\alpha _{max}`$ depends on the area of the covered region; this angle is to be experimentally obtained (Sáez & Arnau, 1997), it is smaller than the size of the map and large enough to include as much scales as possible. The autocorrelation function $`C_\sigma (\alpha )`$ is
$$C_\sigma (\alpha )=\left(\frac{\delta T}{T}\right)_\sigma (\stackrel{}{n}_1)\left(\frac{\delta T}{T}\right)_\sigma (\stackrel{}{n}_2),$$
(2)
where $`\alpha `$ is the angle formed by the unit vectors $`\stackrel{}{n}_1`$ and $`\stackrel{}{n}_2`$, the angular brackets stand for an average over many full realizations of the CMB sky and, quantity $`(\delta T/T)_\sigma (\stackrel{}{n})`$ is the temperature contrast in the direction $`\stackrel{}{n}`$ after smoothing with a Gaussian beam having a certain $`\sigma =0.425\theta _{_{FWHM}}`$. The modified spectra are used below to analyze some simulated maps.
## 2 SIMULATIONS
The angular power spectrum $`C_{\mathrm{}}=_{m=\mathrm{}}^{m=\mathrm{}}|a_\mathrm{}m|^2/(2\mathrm{}+1)`$ is only an auxiliary element in our estimations. We are not particularly interested in any choice and, consequently, we have used the same spectrum as in Sáez, Holtmann & Smoot (1996). It corresponds to the minimum cold dark matter model with a baryonic density parameter $`\mathrm{\Omega }__B=0.03`$ and a reduced Hubble constant $`h=0.5`$. The $`C_{\mathrm{}}`$ coefficients have been taken from Sugiyama (1995) and renormalized according to the four-year COBE data ($`Q_{rms_PS}18\mu K`$ , Gorski et al. 1996). Our simulations are performed by using the Fast Fourier Transform (see Sáez, Holtmann & Smoot 1996 and Bond & Efstathiou 1987) and, then, a certain beam is used to average temperatures; thus, we obtain maps which must be deconvolved with the same beam. After S–deconvolution, the resulting map must be compared with the initial one.
## 3 IDEAL BEAM S–DECONVOLUTION
Two methods are proposed to perform a S–deconvolution of the beam in the absence of noise (ideal case). The efficiency of these methods is verified and the limits for their application are discussed. For appropriate coverages and beam structures, the size of the smallest pixels compatible with beam reversion is estimated in each case. The conclusions obtained in this section are important to understand realistic S–deconvolution in noisy maps (§ 4).
### 3.1 BEAM
We begin with a Gaussian spherically symmetric beam. If the direction of the beam center is $`\stackrel{}{n}`$, the measured temperature $`T(\stackrel{}{n})`$ is given by the following average:
$$T(\stackrel{}{n})=\frac{1}{2\pi \sigma ^2}T^{}(\theta ,\varphi )e^{[\theta ^{}(\stackrel{}{n})]^2/2\sigma ^2}\mathrm{sin}\theta d\theta d\varphi ,$$
(3)
where $`\sigma `$ defines the beam size, the angles $`\theta `$ and $`\varphi `$ are the spherical coordinates of a certain pixel, the element of solid angle is $`d\mathrm{\Omega }=\mathrm{sin}\theta d\theta d\varphi `$ and, quantity $`\theta ^{}`$ is the angle formed by the direction ($`\theta `$, $`\varphi `$) and the observation direction $`\stackrel{}{n}`$.
Small pixels can be considered as surface elements and, consequently, Eq. (3) can be discretized as follows:
$$T(\stackrel{}{n})=\frac{1}{2\pi \sigma ^2}\underset{i}{}T_i^{}e^{(\theta _i^{})^2/2\sigma ^2}\frac{dS_i^2}{R^2},$$
(4)
where the subscript $`i`$ stands for the $`i`$-th pixel; here, $`T_i^{}`$ and $`dS_i`$ are the temperature and the area of the $`i`$-th pixel, respectively, and $`\theta _i^{}`$ is the angle formed by this pixel and the beam center. The exponential tends rapidly to zero as $`\theta _i^{}`$ increases beyond $`\sigma `$; hence, only a small number of pixels are significant in order to estimate $`T(\stackrel{}{n})`$. Furthermore, due to technical reasons, the beam could receive energy from a reduced number of pixels (not from all the significant pixels in an ideal infinite Gaussian beam). By these reasons, we assume that only $`q\times q`$ pixels are relevant and we give various values to number $`q`$; these pixels cover a square patch centered at the same point as the beam.
An asymmetric beam of the form
$$W=\frac{1}{2\pi \sigma ^2}e^{[a^2(\theta \theta ^{})^2+a^2(\varphi \varphi ^{})^2]/2\sigma ^2}$$
(5)
has been also considered. The parameter $`a`$ defines the degree of asymmetry.
Figure 1 illustrates, for $`q=7`$, three situations corresponding to beams and pixelizations considered below.
### 3.2 COVERAGE
Our choice of an appropriate partial coverage is based on some results obtained in previous papers. In Sáez, Holtmann & Smoot (1996), it was shown that CMB maps close to $`20^{}\times 20^{}`$ can be simulated –with good accuracy– neglecting curvature and using the Fourier transform. The effects of partial coverage were analyzed in detail in Sáez & Arnau (1997), these authors proved that, although a $`20^{}\times 20^{}`$ map does not suffice to get the angular power spectrum, a few tens of $`20^{}\times 20^{}`$ maps allow us to find the power spectrum with a high accuracy for $`\mathrm{}>200`$. On account of this previous result, forty $`20^{}\times 20^{}`$ maps are used in this paper to get each angular power spectrum. These maps would cover about $`40\%`$ of the sky. In the above papers, it was also argued that uncorrelated noise is not expected to be important for coverages greater than or equal to $`20^{}\times 20^{}`$ and for the noise level expected in modern CMB experiments. This is true in the sense that the power spectrum can be calculated in the presence of this noise, but the noise can be problematic for S–deconvolution as we will discuss below. These comments point out the interest of considering $`20^{}\times 20^{}`$ maps and motivate our choice of these regions to begin with our analysis of the S–deconvolution procedure.
### 3.3 PIXELIZATION AND EQUATIONS TO BE SOLVED
In order to compute the integral in (3) using Eq. (4), the area $`dS_i`$ is not required to be independent on $`i`$; namely, no equal area distributions of pixels are necessary. Furthermore, if we take a large enough number of small pixels covering all the region contributing significantly to the integral (3), moderated variations in the pixel shapes are also admissible. In spite of these comments, equal area and equal shape pixelizations are advantageous –at least from the mathematical point of view– as it is shown along the paper.
Since a $`20^{}\times 20^{}`$ region is approximately flat, small squares with edges of angular lenght $`\mathrm{\Delta }\theta =\mathrm{\Delta }\varphi =\mathrm{\Delta }`$ define an approximately regular pixelization. The number of pixels per edge is $`N=20/\mathrm{\Delta }`$, where angle $`\mathrm{\Delta }`$ must be given in degrees. Our measurements would cover the pixelized region with a certain strategy. The beam center should point towards each pixel $`\alpha `$ various times and, then, the mean of the resulting measures could be considered as the smoothed temperature $`T_\alpha `$ corresponding to pixel $`\alpha `$. From these $`T_\alpha `$ values, a new temperature which does not involve the beam effect, $`T_i^{}`$, must be assigned to each pixel $`i`$; namely, the beam must be deconvolved. All along § 3, it is assumed that sistematic errors have been corrected and that the uncorrelated instrument noise is negligible. If $`T_\alpha `$ is the temperature measured by the instrument when the beam center is pointing towards the center of the pixel $`\alpha `$, according to Eq. (4), we can write
$$T_\alpha =\underset{i}{}A_{\alpha i}T_i^{}$$
(6)
with
$$A_{\alpha i}=\frac{1}{2\pi \sigma ^2}e^{(\theta _i(\stackrel{}{n}_\alpha )^{})^2/2\sigma ^2}\frac{dS_i^2}{R^2},$$
(7)
where $`T_i^{}`$ is the true temperature in the i-th pixel; namely, the S–deconvolved temperature we are looking for. Similar equations hold for the asymmteric beam (5).
It is evident that the temperatures $`T_i^{}`$ only define a S–deconvolution if they are an approximate solution of the linear Eqs. (6). Only in this case, the temperatures $`T_i^{}`$ are similar to the true temperatures averaged by the beam and, consequently, the resulting maps contain the right spectra. This fact strongly restrict the methods for S–deconvolution. Two of them are described in next section. These equations can be written in the matrix form $`T=AT^{}`$, where $`T`$ and $`T^{}`$ are arrays of $`N\times N`$ numbers (one number for each pixel) and A is a $`N^2\times N^2`$ matrix. The element $`A_{\alpha i}`$ weights the contribution of the pixel $`i`$ to the smoothed temperature at pixel $`\alpha `$. Quantity $`A_{\alpha i}`$ is assumed to be significant only in the $`q\times q`$ pixels mentioned in § 3.1.
### 3.4 METHODS AND RESULTS
Two methods are used to estimate the S–deconvolved temperature $`T_i^{}`$ corresponding to pixel $`i`$: In the first one, Eqs. (6) are solved as a linear system of algebraic equations (hereafter LS–deconvolution) where the independent terms are the observed temperatures $`T_\alpha `$. In the second method, equation (3) is considered as a convolution and, then, Fourier transform (FT) and the deconvolution theorem are used to get $`T_i^{}`$ on the nodes of the 2D Fourier grid (hereafter FTS–deconvolution).
#### 3.4.1 LS–DECONVOLUTION
Iterative methods (Golub & van Loan, 1989; Young, 1971) can be used to solve the system (6). In any of these methods, the matrix $`A`$ is split as follows A=M-P, where M is any matrix which can be easily inverted. In matrix form, the iteration scheme reads as follows:
$$T^{(n+1)}=M^1PT^{(n)}+M^1T,$$
(8)
where the superscript $`n`$ stands for the n-th iteration. The necessary and sufficient condition for convergence is that the spectral radius of the matrix $`Q=M^1P`$ is smaller than unity. This radius is the maximum of $`|\lambda _i|`$, where $`\lambda _i`$ are the eigenvalues of $`Q`$. Hereafter, the so-called Jacobi method is used. This method corresponds to a particular choice of the matrix $`M`$. This matrix is assumed to be the matrix formed by the diagonal of $`A`$, which is denoted $`A__D`$. A sufficient condition for the convergence of the Jacobi method is that matrix $`A`$ is diagonal dominant ($`|A_{ii}|>_{ji}|A_{ij}|`$ for any $`i`$). The dimension of the matrices $`A`$, $`A__D`$ and $`P`$ are $`N^2\times N^2`$, where N is the number of pixels per edge in the map. Since this number is greater than $`10^2`$ in all the practical cases, the dimension of the above matrices is very great and they cannot be stored. Fortunately, this storage is not necessary; in fact, if Eq. (8) is rewritten using indices
$$A_{ii}T_i^{(n+1)}=T_i\underset{j<i}{}A_{ij}T_j^{(n)}\underset{j>i}{}A_{ij}T_j^{(n)},$$
(9)
we see that all the elements of the matrix $`A`$ necessary to get $`T_i^{(n+1)}`$ can be obtained when they are necessary, without storage. Of course, a given element can be calculated various times, but no storage is necessary at all.
Equation (8) can be considered as a matrix equation in which each element is a block. Matrix $`A`$ is split in $`N\times N`$ blocks and vectors $`T`$ and $`T^{}`$ in N arrays (blocks) of dimension $`N`$. Then the resulting $`N\times N`$ blocks appear to have a q-diagonal structure and diagonal domination gives limit values for quantities $`\lambda =e^{}(\mathrm{\Delta }^2/2\sigma ^2)`$ and $`\theta _{_{FWHM}}/\mathrm{\Delta }`$. The maximum value of $`\lambda `$ and the corresponding maximum value of the ratio $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ are given in the first and second columns of Table 1, respectively, for various $`q`$ values. The third column gives the ratio $`S_q/S__B`$, where $`S_q`$ is the area covered by the $`q\times q`$ significant pixels and $`S__B`$ is the area of the beam (a circle with radius $`\theta _{_{FWHM}}`$). The top panel of Fig. 1 illustrates the case $`q=7`$. Eleven pixels are located inside the beam circle (radius $`\theta _{_{FWHM}}`$).
Taking into account that, in general, the condition $`\lambda <\lambda _{max}`$ is only a sufficient condition for the convergence of the Jacobi method (not necessary), we have studied numerically many cases in which $`\lambda \lambda _{max}`$ with the hope of getting new convergent cases. The result is that the Jacobi method has never converged for $`\lambda >\lambda _{max}`$. This result points out that, in practice, the condition $`\lambda \lambda _{max}`$ is also necessary for LS–deconvolution. This information suffices for us. A rigorous mathematical study about the necessary character of this condition is not appropriate here. From the intuitive point of view, the existence –in practice– of a $`\lambda _{max}`$ value is an expected result, in fact, in the absence of a $`\lambda _{max}`$, it would be possible to assign deconvolved temperatures to billions (an arbitrary number) of small pixels placed inside the beam and, furthermore, the right spectrum could be recovered up to the spatial scales of these small pixels; this would be a nonsense.
When LS–deconvolution applies, it is an accurate method. In fact, for $`\theta _{_{FWHM}}=8.8^{}`$ and $`\mathrm{\Delta }=4.6875^{}`$, about 24 iterations suffice to get a very good deconvolved map. After these iterations, the relation $`[_{i=1}^{N^2}(T_i^{(n+1)}T_i^{(n)})^2]^{1/2}/N^2<10^7`$ is satisfied and, consequently, the method is converging towards a certain map. The question is: are the numerical iterations converging to a good S–deconvolved map with the right spectrum? In order to answer this question, we proceed as follows: (1) the Fast Fourier Transform (FFT) is used to do a simulation ($`S_1`$) which does not involve either beam or noise, (2) a second map ($`S_2`$) is obtained –from $`S_1`$– by using a certain beam for smoothing, (3) the map $`S_2`$ is deconvolved using the Jacobi method to get a new map ($`S_3`$), (4) the three above steps are repeated forty times (see § 3.2) and, (5) the method described in the introduction (see also Sáez, Holtmann & Smoot 1996 and Sáez & Arnau 1997) is used to obtain the modified spectrum in three cases: before smoothing (from $`S_1`$ maps), after smoothing (from $`S_2`$ maps) and, after deconvolution (from $`S_3`$ maps), these spectra are hereafter referred to as $`E_1\mathrm{}`$, $`E_2\mathrm{}`$, and $`E_3\mathrm{}`$, respectively. If deconvolution is a good enough S–deconvolution, spectrum $`E_3\mathrm{}`$ should be comparable with $`E_1\mathrm{}`$. Whatever the deconvolution method may be, these five steps allow us to analize the resulting deconvolved maps. In all the Figures of this paper which show the three above spectra, pointed, dashed, and solid lines correspond to $`E_1\mathrm{}`$, $`E_2\mathrm{}`$, and $`E_3\mathrm{}`$, respectively. The top right panel of Fig. 2 shows the resulting spectra for the LS–deconvolution under consideration ($`\theta _{_{FWHM}}=8.8^{}`$ and $`\mathrm{\Delta }=4.6875^{}`$). We see that the dotted line ($`E_1\mathrm{}`$) is almost indistinguishable from the solid one ($`E_3\mathrm{}`$) for $`\mathrm{}<2000`$. This result qualitatively proves the goodness of the iterative LS–deconvolution. In order to compare the spectra $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ quantitatively, the following quantities are calculated and presented in Table 2 (entries 7 and 8): The mean, $`M1`$, of the quantities $`\mathrm{}(\mathrm{}+1)E_{\mathrm{}}\times 10^{10}`$ corresponding to the spectrum $`E_1\mathrm{}`$ (col. ), the mean, $`M2`$, of the differences $`E_1\mathrm{}E_3\mathrm{}`$ (col. ), the mean MA of $`|E_1\mathrm{}E_3\mathrm{}|`$ (col. ), and the typical deviation, $`\mathrm{\Sigma }`$, of the differences of column (4) (col. ). The above quantities are computed in appropriate $`\mathrm{}`$–intervals (col. ). Entries 7 and 8 show values of $`|M2|`$, $`MA`$ and $`\mathrm{\Sigma }`$ much smaller than $`|M1|`$, which means that the spectra $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are very similar in both $`\mathrm{}`$–intervals (40,1000) and (1000,2000). It is also remarkable that $`|M2|`$ is much smaller than MA, which means that spectrum $`E_3\mathrm{}`$ oscillates around spectrum $`E_1\mathrm{}`$ giving positive and negative values of $`E_1\mathrm{}E_3\mathrm{}`$ which cancel among them.
Table 2 compares other pairs of spectra displayed in Figs. 2 and 3. Column (1) gives the Figure and panel where each pair of spectra are displayed. Tables 3 and 4 have the same structure as Table 2, but they compare pairs of spectra contained in Figs 4 and 5, respectively. The interpretation of the data exhibited in these Tables is straightforward. For a given entry, the compared spectra are similar if the quantities $`|M2|`$, $`MA`$ and $`\mathrm{\Sigma }`$ are much smaller than $`|M1|`$. Given two entries with similar $`|M1|`$ values, the smaller the values of $`|M2|`$, $`MA`$ and $`\mathrm{\Sigma }`$, the greater the similarity between the spectra (the better the deconvolution if we are comparing $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ spectra).
#### 3.4.2 FTS–DECONVOLUTION
The method based on the FT only can be used in the case of small enough coverages (almost flat regions) allowing a uniform pixelization. The region covered by the observations should be a square with no much more than $`20^{}`$ per edge; thus, curvature can be neglected and the covered area can be considered as a square where the angles $`\theta `$ and $`\varphi `$ play the role of cartesian coordinates. Equation (3) can be then seen as a convolution of the function $`U(\theta ,\varphi )=T^{}(\theta ,\varphi )\mathrm{sin}(\theta )`$ with the Gaussian beam function $`W=\frac{1}{2\pi \sigma ^2}e^{[(\theta \theta ^{})^2+(\varphi \varphi ^{})^2]/2\sigma ^2}`$, where coordinates $`\theta ^{}`$ and $`\varphi ^{}`$ define the observation direction $`\stackrel{}{n}`$. Then, the deconvolution theorem ensures that the Fourier transform of function $`U`$ is
$$U(\stackrel{}{k})=T(\stackrel{}{k})/W(\stackrel{}{k}),$$
(10)
where $`\stackrel{}{k}`$ is a vector in the 2-dimensional Fourier space. Given a smoothed map and a window function, we can find their Fourier transforms $`T(\stackrel{}{k})`$ and $`W(\stackrel{}{k})`$ and, then, Eq. (10) plus an inverse FT allows us to find the deconvolved map. Unfortunately, the use of the FT is not compatible with spherically assymmetric rotating beams. If one of these beams measures in such a way that its orientation changes from measure to measure, Eq. (3) is not a convolution anymore and, consequently, the FTS–deconvolution does not apply; hence, the FT can be used either in the case of a spherically symmetric beam or in the case of a nonspherical nonrotating beam which measures preserving its orientation.
The maximum value of the ratio $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ compatible with FTS–deconvolution has been derived using simulations. In all the $`20^{}\times 20^{}`$ simulations, we have taken $`N=256`$ ($`\mathrm{\Delta }=4.6875^{}`$), while quantity $`\theta _{_{FWHM}}`$ has been varied appropriately. The code for FTS–deconvolution has been run in each case. This code follows the five steps of the process described above for analyzing deconvolved maps. The left panels of Fig. 2 show the spectra $`E_1\mathrm{}`$, $`E_2\mathrm{}`$, and $`E_3\mathrm{}`$, for different values of $`\theta _{_{FWHM}}`$. The top left panel, which corresponds to $`\theta _{_{FWHM}}=10^{}`$, shows that the spectra $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are very similar in all the $`\mathrm{}`$–interval (2,2000). For $`\theta _{_{FWHM}}11^{}`$ (middle left panel), these spectra are very similar for $`2\mathrm{}1000`$, while they become a little diferent in the interval (1000,2000). For $`\theta _{_{FWHM}}>11^{}`$ , the differences between $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ grow rapidly and, for $`\theta _{_{FWHM}}=11.5^{}`$ (bottom left panel), these spectra are very different. This is quantitatively confirmed by the numbers presented in entries 1 to 6 of Table 2, where we see that, in the interval (1000,2000), quantities $`|M1|`$, $`MA`$ and $`\mathrm{\Sigma }`$ increase as $`\theta _{_{FWHM}}`$ does. The maximum value of $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ appears to be $`2.3`$. This means that, in the case $`1.87<\theta _{_{FWHM}}/\mathrm{\Delta }2.3`$, the FTS–deconvolution applies and the LS–deconvolution does not. The most dense grid compatible with FTS–deconvolution is shown in the middle panel of Fig. 1. In this case, around 16 pixels can be placed inside the beam circle.
An asymmetric nonrotating beam of the form (5) has been also deconvolved for various values of the parameters $`a`$ and $`\theta _{_{FWHM}}`$. The middle right panel of Fig. 2 shows the results of the FTS–deconvolution for $`a=1.29`$ and $`\theta _{_{FWHM}}=7.75^{}`$. These results are good in the full $`\mathrm{}`$–interval (2,2000) (see entries 9 and 10 of Table 2). This choice of the parameters $`a`$ and $`\sigma `$ simulates an asymmetric beam whose effective $`\theta _{_{FWHM}}`$ along the $`\theta `$-axis ($`\varphi `$-axis) is $`\theta _1^{eff}=\theta _{_{FWHM}}/a6^{}`$ ($`\theta _2^{eff}=a\theta _{_{FWHM}}10^{}`$). This beam and the most dense pixelization allowing its FTS–deconvolution are shown in the bottom panel of Fig. 1. Around nine pixels are located inside the beam ellipse. Other $`a`$ and $`\sigma `$ values have been also considered to conclude that FTS–deconvolution is possible when both $`\sigma _1^{eff}/\mathrm{\Delta }`$ and $`\sigma _2^{eff}/\mathrm{\Delta }`$ are smaller than $`2.3`$ (this constraint is equivalent to that obtained in the spherically symmetric case). Finally, we have compared the effect of a symmetric beam with $`\theta _{_{FWHM}}=10^{}`$ and that of the asymmetric nonrotating beam described above. In order to do this comparison we have obtained forty $`S2`$ maps with each beam. The spectra obtained from these S2 maps are displayed in the bottom right panel of Fig. 2 and quantitatively compared in entries 11 and 12 of Table 2. The solid (pointed) line corresponds to the asymmetric (symmetric) beam. Results show that the deformations of the original spectrum produced by these beams are different; namely, that the $`E_2\mathrm{}`$ spectra are distinct (significant asymmetry).
Vanishing instrumental noise has been assumed so far; nevertheles, partial coverage introduces a kind of sky noise (an uncertainty). In order to estimate this noise for a coverage of forty $`20^{}\times 20^{}`$ maps, we compare the $`S1`$ spectrum extracted from these maps with the theoretical spectrum (which would correspond to many realizations of the full sky). Both spectra are presented in Fig. 3, where the solid (pointed) line corresponds to the S1 (theoretical) spectrum. The quantitative comparison of these spectra is given in entries 13 and 14 of Table 2. We see that the deviations with respect to the true spectrum –produced by the partial coverage under consideration– are greater than those produced by good S-deconvolutions (compare entries 13 and 14 of Table 2 with the pairs of entries 1–2, 7–8, and 9–10. Compare also the corresponding panels in the Figures). The deviations decrease as the coverage increases.
## 4 BEAM DECONVOLUTION IN NOISY MAPS
In § 3, negligible uncorrelated noise has been assumed, thus, for a given beam, a minimum pixel size for LS–deconvolution and another one for FTS–deconvolution have been found. These minima define theoretical restrictions for admissible pixelization; nevertheless, in the presence of uncorrelated noise, stronger restrictions on pixel sizes could appear and, consequently, S–deconvolution could be impossible for some sizes close to the minimum size obtained in the absence of noise. Which is the minimum size allowing FTS–deconvolution in the presence of a certain level of uncorrelated noise? We are going to study this question.
As it is well known, the amount of noise in a map depends on the observing time per pixel, $`t_{pix}`$, which is inversely proportional to the pixel area. At pixel i, the noise contributes to the temperature an amount $`\delta T_i^^N`$. It is assumed that the noise is uncorrelated and has uniform variance $`\sigma __N^2`$; i.e., $`\delta T_i^^N\delta T_j^^N=\sigma __N^2\delta _{ij}`$. The relation $`\sigma __N=s/(t_{pix})^{1/2}`$ can be used to estimate the level of uncorrelated noise in the pixelized map, where $`s`$ is the detector sensitivity (see Knox, 1995). For $`s=200\mu K\sqrt{(}sec)`$ , $`\mathrm{\Delta }=4.6875^{}`$, and a year of uniform full-sky coverage, the level of noise is $`\sigma __N93.9\mu K`$. Furthermore, using the $`C_{\mathrm{}}`$ numbers of § 2 with $`\mathrm{\Delta }=4.6875^{}`$ and $`\theta _{_{FWHM}}=8.8^{}`$, the expected signal $`S=[(1/4\pi )(_{\mathrm{}}(2\mathrm{}+1)C_{\mathrm{}}e^{\sigma ^2\mathrm{}^2})]^{1/2}`$ takes on the value $`S108\mu K`$. Therefore, for the above choice of $`s`$, $`\theta _{_{FWHM}}`$, and $`\mathrm{\Delta }`$, the signal to noise ratio $`S/\sigma __N`$ is close to $`1`$ and, consequently, noise cannot be neglected a priori in order to do beam S–deconvolution. In other realistic cases, the situation is similar. Since $`\sigma __N`$ is proportional to $`(t_{pix})^{1/2}`$, too long observation times would be necessary to rise $`S/\sigma __N`$ significantly. Fortunately, technological progress leads to smaller and smaller $`s`$ values and, accordingly, the ratio $`S/\sigma __N`$ increases. In the PLANCK project of the European Spatial Agency there are two instruments: the Low Frequency Instrument (using radiometers) and the High Frequency Instrument (using bolometers). One of the radiometers ($`\nu =100GHz`$ and $`\theta _{FWHM}=10^{}`$), working at $`20^{}K`$, would produce a noise $`\sigma __N54\mu K`$ (for pixels with $`\mathrm{\Delta }=4.6875^{}`$) during a year of uniform coverage; moreover, for one of the bolometers ($`\nu =143GHz`$ and $`\theta _{FWHM}=10.3^{}`$), which would work at $`0.1^{}K`$, the noise would be $`\sigma __N16.4\mu K`$ for the same pixels and time coverage. These data are taken into account below in order to analyze the perspectives of beam S–deconvolution in the framework of the most accurate project for anisotropy detection in small angular scales (the PLANCK mission of the European Spatial Agency).
In spite of the fact that LS–deconvolution has been very useful in order to analyze and understand the existence of a minimum size for pixelization, in practice, only the FTS–deconvolution has been used –so far– in the noisy case. The maximum value of $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ compatible with FTS–deconvolution depends on the level of uncorrelated noise. For this type of noise, if the average in Eq. (2) is performed on many sky realizations, the resulting $`C_\sigma (\alpha )`$ values must be very small; nevertheless, if the average is done in a $`20^{}\times 20^{}`$ patch, the $`C_\sigma (\alpha )`$ values can be relevant, which means that, on the patch, the noise is not properly uncorrelated (its spectrum is unknown). This fact is important in order to understand the effect of this noise on FTS–deconvolution of $`20^{}\times 20^{}`$ maps. In the presence of a certain noise which is independent on the signal, FTS–deconvolution could be performed using the so-called optimal Wiener filter (Press et al. 1988). In such a case, the spectrum of the function $`U`$ –defined above– is estimated as follows:
$$U(\stackrel{}{k})=T(\stackrel{}{k})\mathrm{\Phi }(\stackrel{}{k})/W(\stackrel{}{k}),$$
(11)
where
$$\mathrm{\Phi }(\stackrel{}{k})=1\frac{|N(\stackrel{}{k})|^2}{|U(\stackrel{}{k})|^2+|N(\stackrel{}{k})|^2}.$$
(12)
In the absence of noise, function $`\mathrm{\Phi }`$ takes on the form $`\mathrm{\Phi }(\stackrel{}{k})=1`$ and Eq. (11) reduces to Eq. (10). Equations (11) and (12) cannot be used in practice to deconvolve the beam (unknown spectrum of a given $`20^{}\times 20^{}`$ noise realization); nevertheless, these equations are useful to understand why the noise can be neglected in some cases. The maximum of the $`|U(\stackrel{}{k})|^2`$ values corresponding to forty $`20^{}\times 20^{}`$ simulations –based on the model of § 2– has been estimated to be $`1.7`$, while the maximum of $`|N(\stackrel{}{k})|^2`$ obtained from the same number of simulations of pure noise (uncorrelated in great regions) has appeared to be proportional to the level of noise $`\sigma __N`$. For $`\sigma __N=16.4\mu K`$, the resulting maximum is $`2.7\times 10^3`$. In this case –and also for any current or planned experiment– the amplitude corresponding to $`|N(\stackrel{}{k})|^2`$ is much smaller than that of $`|U(\stackrel{}{k})|^2`$. This smallness –relative to that of the signal– indicates that, in realistic noisy cases, the filter function is close to $`\mathrm{\Phi }(\stackrel{}{k})=1`$ and the following question arises: Is it possible to take $`\mathrm{\Phi }(\stackrel{}{k})=1`$ (noise neglection) to reverse beam smoothing? No theoretical arguments have been found to answer this question. Numerical simulations have been necessary. Results obtained from simulations are displayed in Fig. 4. In all the cases studied, the noise has been neglected and Eq. (10) has been used to perform FTS–deconvolution. Good results indicate that noise neglection is appropriate.
In the left panels of Fig. 4, the level of noise is $`16.4\mu K`$. For $`\theta _{_{FWHM}}=8.8^{}`$ (top left panel) spectra $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are quasi indistinguishable in the $`\mathrm{}`$–interval (2,2000). For $`\theta _{_{FWHM}}9.5^{}`$, a small difference between these spectra appears for $`1000\mathrm{}2000`$ (middle left panel). Finally, $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are clearly different for $`\theta _{_{FWHM}}=10^{}`$ (bottom left panel). Hence, the maximum value of $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ is close to 2 (around 12 pixels inside the beam circle). The right panels of Fig. 4 show the same analysis as the left panels for a noise level of $`54.6\mu K`$. In this case, spectra $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are similar in all the interval (2,2000) for $`\theta _{_{FWHM}}<8^{}`$ (top right panel), small discrepancies in the interval (1000,2000) have already appeared for $`\theta _{_{FWHM}}8.5^{}`$ (middle right panel) and, finally, for $`\theta _{_{FWHM}}9^{}`$ these discrepancies are important (bottom right panel). For this level of noise, the maximum value of $`\theta _{_{FWHM}}/\mathrm{\Delta }`$ is close to 1.8 (around 10 pixels inside the beam circle). This qualitative analysis of Fig. 4 is confirmed by Table 3 where quantities $`M1`$, $`M2`$, $`MA`$, and $`\mathrm{\Sigma }`$ are presented in all the cases.
## 5 PIXELIZATION AND DECONVOLUTION IN REALISTIC EXPERIMENTS
Limitations of the S–deconvolution process have been discussed along the paper. If the pixel size $`\mathrm{\Delta }`$ is taken to be similar to the beam radius $`\theta _{_{FWHM}}`$, the angular power spectrum can be only estimated for $`\mathrm{}\mathrm{}_{max}`$ with $`\mathrm{}_{max}180/\mathrm{\Delta }180/\theta _{_{FWHM}}`$; however, for $`\mathrm{\Delta }\theta _{_{FWHM}}/2`$, the angular power spectrum can be evaluated up to $`\mathrm{}_{max}360/\theta _{_{FWHM}}`$. In these formulae, $`\mathrm{\Delta }`$ and $`\theta _{_{FWHM}}`$ must be written in degrees. We see that, for $`\theta _{_{FWHM}}10^{}`$, the pixelization $`\mathrm{\Delta }\theta _{_{FWHM}}`$ ($`\mathrm{\Delta }\theta _{_{FWHM}}/2`$) allows us to get the spectrum up to $`\mathrm{}_{max}1080`$ ($`\mathrm{}_{max}2160`$); therefore, for a given beam, the choice of the best feasible pixelization is crucial in order to get maximum information from observations. The minimum pixel size compatible with S–deconvolution (for the methods used in the paper) is hereafter denoted $`\mathrm{\Delta }_{_{DE}}`$.
In realistic experiments, various effects –apart from deconvolution– conditionate the choice of the most appropriate pixelization. In order to discuss these effects, let us focus our attention on PLANCK project (see Tauber, 1999). With a telecope having a diameter $`D`$, the minimum pixel size allowed by difraction is roughly $`\mathrm{\Delta }_{_{DI}}1.22c/D\nu `$, where $`c`$ is the speed of light and $`\nu `$ its frequency; hence, this minimum size depends on $`\nu `$. In a multifrequency experiment, there is a minimum pixel size $`\mathrm{\Delta }_{_{DI}}`$ corresponding to each frequency; for example, in the PLANCK mission ($`D=1.5m`$), observations will be carried out in nine different frequencies ranging from $`30GHz`$ to $`857GHz`$ and, consequently, the size $`\mathrm{\Delta }_{_{DI}}`$ ranges from $`30^{}`$ to $`1^{}`$. Furthermore, in the PLANCK case, the line of sight will move on a big circle in the sky each minute; hence, if two successive temperature assignations on the circle are performed at an angular distance $`\mathrm{\Delta }\alpha `$, the time –in seconds– between these asignations is $`\mathrm{\Delta }t=2.78\times 10^3\mathrm{\Delta }\alpha `$. The angle $`\mathrm{\Delta }\alpha `$ must be chosen in such a way that (1) no large overlaping of contiguous beam positions occurs and (2) time $`\mathrm{\Delta }t`$ is greater that the response time of the bolometers. For the chosen period of one minute and $`\mathrm{\Delta }\alpha 2\theta _{_{FWHM}}`$, there is no overlaping and condition (2) is satisfied for the PLANCK bolometers. A certain pixel size is only admissible if technology plus observational strategy ensure that each pixel is observed a large enough number of times during the mission. Let us estimate this number for PLANCK. For a pixel size $`\mathrm{\Delta }`$, the total number of pixels is $`N=1.5\times 10^8\mathrm{\Delta }^2`$ and admitting uniform coverage during a year (for qualitative estimates) each pixel is observed for a time $`\mathrm{\Delta }t_p=0.2\mathrm{\Delta }^2s`$; therefore, the number of observations per pixel is $`N_p=\mathrm{\Delta }t_p/\mathrm{\Delta }t=72\mathrm{\Delta }^2/\mathrm{\Delta }\alpha `$. Finally, for $`\mathrm{\Delta }\alpha =2\theta _{_{FWHM}}`$, one easily see that the size necessary to obtain $`N_p`$ observations by pixel (during a year of PLANCK mission) is $`\mathrm{\Delta }^{}=\frac{1}{6}(N_p\theta _{_{FWHM}})^{1/2}`$. Since the $`\theta _{_{FWHM}}`$ values for PLANCK detectors range from $`30^{}`$ to $`4.5^{}`$, assuming $`N_p100`$, we see that $`\mathrm{\Delta }^{}`$ ranges from $`9.5^{}`$ to $`1.3^{}`$. This means that, in order to have a number of observations by pixel greater than 100, the pixel size must be greater than $`\mathrm{\Delta }^{}(N_p=100)=\mathrm{\Delta }_{100}`$.
Given a frequency, there is an optimum pixel size, $`\mathrm{\Delta }_{_{OP}}`$, which will be assumed to be the maximum of the three above sizes $`\mathrm{\Delta }_{DE}`$, $`\mathrm{\Delta }_{DI}`$ and $`\mathrm{\Delta }_{100}`$. The value of $`\mathrm{\Delta }_{_{OP}}`$ depends on the frequency.
In order to separate the foregrounds and the cosmological signal in multifrequency experiments, various frequencies and a unique pixelization must be used and, consequently, the best pixelization would be the maximum of the $`\mathrm{\Delta }_{_{OP}}`$ optimal sizes corresponding to the involved frequencies. This maximum corresponds to the lowest frequency under consideration; for example, in the PLANCK case, if all the frequencies from $`30GHz`$ to $`857GHz`$ are considered, the minimun admissible pixel appears to have a size $`\mathrm{\Delta }=\mathrm{\Delta }_{_{DI}}(\nu =30)30^{}`$, for which, the spectrum can be only estimated up to $`\mathrm{}_{max}=360`$. Of course, we could consider only the frequencies greater than $`53GHz`$ (with some loss of information) and, then, the minimum pixel is $`\mathrm{\Delta }=\mathrm{\Delta }_{_{DI}}(\nu =53)17^{}`$ and $`l_{max}=630`$ and so on.
Let us reconsider the radiometer working at $`\nu 100GHz`$ with $`\theta _{_{FWHM}}=10^{}`$, which was projected to be inside PLANCK satellite. For this radiometer one easily find $`\mathrm{\Delta }_{_{DE}}5^{}`$, $`\mathrm{\Delta }_{_{DI}}=8.4^{}`$ and $`\mathrm{\Delta }_{100}3^{}`$; hence, $`\mathrm{\Delta }_{_{OP}}=\mathrm{\Delta }_{_{DI}}=8.4^{}`$ ($`\mathrm{}_{max}1290`$), for this optimum pixelization, the level of noise is $`5.3\times 10^6`$ and the angular spectrum can be obtained for $`\mathrm{}1290`$ (top panel of Fig. 5 and entries 1 and 2 of Table 4). We now consider the bolometer working at $`\nu 143GHz`$ with $`\theta _{_{FWHM}}=10.3^{}`$, which was also proposed to measure CMB anisotropy from the PLANCK satellite (phase A study). For this detector we easily find $`\mathrm{\Delta }_{_{DE}}5^{}`$, $`\mathrm{\Delta }_{_{DI}}5.87^{}`$ and $`\mathrm{\Delta }_{100}5^{}`$; hence, $`\mathrm{\Delta }_{_{OP}}=\mathrm{\Delta }_{_{DI}}=5.87^{}`$ ($`\mathrm{}_{max}1840`$). For the pixelization $`\mathrm{\Delta }=5.87^{}`$ the level of noise is $`1.28\times 10^5`$ and we have verified that FTS–deconvolution leads to the right spectrum for $`\mathrm{}1840`$ (bottom panel of Fig. 5 and entries 3 and 4 of Table 4).
For the bolometer working at $`857GHz`$ with $`\theta _{_{FWHM}}=4.4^{}`$, we get $`\mathrm{\Delta }_{_{DI}}=1^{}`$ and $`\mathrm{\Delta }_{100}1^{}`$, but the estimate of $`\mathrm{\Delta }_{DE}`$ is problematic as a result of the high level of noise of this bolometer. Perhaps, in this case, maximum entropy or wavelets could give good results; for example, wavelets could be used to lower the noise before beam deconvolution.
## 6 DISCUSSION AND CONCLUSIONS
We expect that beam deconvolution will be important in order to study some aspects of the observational maps given by experiments as PLANCK. As an example, let us argue that the study of the statistical properties of a given observational map should be performed after deconvolution. In fact, various methods can be used to know if the maps are Gaussian or they obey other statistics; among them, the estimation of the correlation function of pixels where the signal is above a certain threshold (excursion sets, Kaiser 1984) and the local analysis of the spots distributed in the map (Bond and Efstathiou, 1987). Since the beam smoothes the map, it alters the correlations of excursion sets and the structure and distribution of the spots; hence, the above methods for analyzing statistics should be applied after a good deconvolution.
The separation of the cosmic signal and the foregrounds requires a unique appropriate pixelization. In the PLANCK case, we have seen that the optimal size for this pixelization, $`\mathrm{\Delta }_{_{OP}}`$, coincides with the size $`\mathrm{\Delta }_{_{DI}}\theta _{_{FWHM}}`$ corresponding to the lowest frequency under consideration; nevertheless, other studies can be imagined (statistical analysis et cetera) which could be performed on the maps corresponding to a given frequency (without previous separation).
For small enough (but feasible) values of $`\sigma __N`$ plus a certain beam (either spherical with a $`\theta _{_{FWHM}}`$ or asymmteric), our codes allow us to find the most appropriate pixel size for deconvolution $`\mathrm{\Delta }=\mathrm{\Delta }_{_{DE}}`$. For the corresponding pixelization, FTS–deconvolution leads to a good estimation of the angular power spectrum in the most wide $`\mathrm{}`$–interval. The size $`\mathrm{\Delta }_{_{DE}}`$ must be compared to $`\mathrm{\Delta }_{100}`$ and $`\mathrm{\Delta }_{_{DI}}`$ to choose the most appropriate pixelization $`\mathrm{\Delta }=\mathrm{\Delta }_{_{OP}}`$. In the case $`\mathrm{\Delta }_{_{OP}}>\mathrm{\Delta }_{_{DE}}`$, the study about beam reversion presented in § 3 to § 5 proves that S–deconvolution can be performed using very simple methods. For levels of noise much higher that those of previous sections, further study is necessary; maximum entropy, wavelets or other methods should be tried out.
The goodness of a certain pixelization against beam S-deconvolution has appeared to be weakly dependent on the particular mathematical method used to reverse the beam average. This fact suggests that deconvolution procedures different from those of this paper could alter its results. Altough this suggestion should be a motivation for studying new methods to get approximate solutions of Eqs. (6) (S–deconvolutions), the structure of the system of linear equations to be solved is always the same and, consequently, all the mathematical methods could exhibit similar limitations to solve it. Indeed, we believe that new deconvolution methods could lead to some modifications of the results of this paper, but not to very different values of $`\mathrm{\Delta }_{_{DE}}`$.
Our estimates show that, in the absence of any problem with beam asymmetry, current technology could lead to a good estimate of the angular power spectrum (of the total signal including foregrounds) in a $`\mathrm{}`$–interval which depends on frequency. In the particular case of two instruments on board of PLANCK, we have used optimum pixelization and forty $`20^{}\times 20^{}`$ maps to find that the spectra are recovered from $`\mathrm{}=200`$ to $`\mathrm{}1300`$ in the case of a radiometer and from $`\mathrm{}=200`$ to $`\mathrm{}1800`$ for a certain bolometer. Since the sky can be divided into $`100`$ of these maps and we only need about $`40`$ for a good estimate of the spectrum (for large $`\mathrm{}`$ values), we can select the best forty maps; namely, the maps having minimum contaminations. The uncertainty produced by this partial coverage appears to be a little greater than the errors produced by the implemented deconvolution procedures (this means that these procedures are good enough for us).
The problem with the deviations of the beam structure with respect to spherical symmetry deserves much attention. As discussed above, FTS–deconvolution is compatible with beam asymmetry if the beam orientation is preserved from measurement to measurement. If the experiment is designed in such a way that the beam does not rotate, our codes for FTS–deconvolution work (see § 3.4.2 and the middle and bottom right panels of Fig. 2); however, if the beam rotates, FTS–deconvolution does not apply and, moreover, operative methods for making beam S–deconvolution are not known; hence, if the assymetry is high enough, beam S–deconvolution is not feasible (so far). In short, excepting the case of negligible deviations with respect to spherical symmetry, any effort directed to maintain unaltered the beam orientation during observations seems to be of great interest. Unfortunately, in spatial projects as Planck, the design of the observational strategy does not preserve this orientation.
Given an asymmetric rotating beam, it would be interesting to study the whole effect of asymmetry plus changing orientation. Even if S–deconvolution is not feasible, the estimation of this whole effect could be a further direct application of the techinques used in this paper. The following method seems to be appropriate: (i) average the asymmetric beam –on appropriate shells– to get a new associated one with spherical symmetry, (ii) simulate forty S1 maps, (iii) smooth the S1 maps with the assymmetric beam taking into account the orientation change produced by the observational strategy; thus, we obtain forty S2 maps, (iv) deconvolve the S2 maps with the spherically symmetric beam of reference to get the S3 maps, and (v) estimate the spectra $`E_1\mathrm{}`$, $`E_2\mathrm{}`$ and, $`E_3\mathrm{}`$. If $`E_1\mathrm{}`$ and $`E_3\mathrm{}`$ are similar enough, the assymetry can be neglected, on the contrary, the differences between these two spectra can be considered as a measure of the whole effect of asymmetry plus rotation. This study should be developed for realistic beams and observational strategies, which is out of the scope of this paper.
ACKNOWLEDGMENTS. This work has been partially supported by the Spanish DGES (project PB96-0797). Some calculations were carried out on a SGI Origin 2000s at the Centro de Informática de la Universidad de Valencia.
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FIGURE CAPTIONS
FIG. 1.– Top (middle) panel shows the relation between the size of a circular beam and that of the pixels for the most dense pixelization compatible with LS (FTS) deconvolution. Bottom panel shows the same for an asymmetric beam and FTS–deconvolution
FIG. 2.– Each panel shows quantity $`\mathrm{}(\mathrm{}+1)E_{\mathrm{}}\times 10^{10}`$ versus $`\mathrm{log}\mathrm{}`$. No noise is present and $`\mathrm{\Delta }=4.6875^{}`$. All panels, excepting the bottom right one, contain three lines: pointed line gives the spectrum before beam smoothing, dashed line is the spectrum after smoothing and, solid line corresponds to the S–deconvolved spectrum. Left: top, middle and bottom panels correpond to FTS–deconvolutions with beams having $`\theta _{_{FWHM}}=10^{}`$, $`\theta _{_{FWHM}}=11^{}`$, and $`\theta _{_{FWHM}}=11.5^{}`$, respectively. Right: top (middle) panel shows the same as the left panels for LS–deconvolution and $`\theta _{_{FWHM}}=8.8^{}`$ (for FTS–deconvolution and the asymmetric beam defined in the text). The dotted (solid) line of the bottom right panel gives the spectrum after smoothing for a spherically symmetric beam with $`\sigma =10^{}`$ (for the asymmetric beam of the text).
FIG. 3.– The quantities represented are the same as in all the panels of Fig. 2. Solid (pointed) line is the modified spectrum extracted from a coverage of forty $`20^{}\times 20^{}`$ maps (the theoretical modified spectrum).
FIG. 4.– The same as in the left panels of Fig. 2. Left: the level of uncorrelated noise per pixel is $`\sigma __N=16.4\mu K`$. Top, middle and bottom panels correpond to $`\theta _{_{FWHM}}=8.8^{}`$, $`\theta _{_{FWHM}}=9.5^{}`$, and $`\theta _{_{FWHM}}=10^{}`$, respectively. Right: the same as in left panels for $`\sigma __N=54.6\mu K`$. Top, middle and bottom panels correpond to $`\theta _{_{FWHM}}=8^{}`$, $`\theta _{_{FWHM}}=8.5^{}`$, and $`\theta _{_{FWHM}}=9^{}`$, respectively. Pixel size is $`\mathrm{\Delta }=4.6875^{}`$ in all cases.
FIG. 5.– The same as in Fig. 4. Top panel shows the results of FTS–deconvolution for a radiometer of the PLANCK mission described in the text. The pixel size is $`\mathrm{\Delta }=8.4^{}`$. The bottom panel shows the same for a bolometer of PLANCK and $`\mathrm{\Delta }=5.87^{}`$. |
warning/0001/astro-ph0001157.html | ar5iv | text | # The Stellar Content of Obscured Galactic Giant H II Regions: II. W42
## 1 INTRODUCTION
In this paper we continue the exploration of the stellar content of obscured Galactic giant H II regions begun by Blum et al. (1999, hereafter Paper I). $`J`$, $`H`$, and $`K`$ images are used to make a broad assessment of the stellar content of obscured star forming regions in the Milky Way. Infrared spectroscopy follows, providing details of the brightest cluster members which can be used to make distance, mass, and luminosity estimates. The spectra are placed in proper context by comparison to new infrared spectral classification systems for massive stars (Hanson et al., 1996; Blum et al., 1997; Figer et al., 1997; Hanson et al., 1998). The known hot star content of the Galaxy is rapidly expanding on the strength of the sophisticated infrared detector gains of the last 10 years. Hanson et al. (1997) have recently published a very detailed account of the most massive stars in the (relatively) nearby giant H II region M17. The present project (Paper I, this work, and future work) seeks to provide a large sample of massive star clusters with which to study the young and massive stellar content in the Galaxy. This sample builds on the detailed visual studies of the Galactic OB associations (Massey et al., 1995) and provides a bridge to studies of young stellar objects in star forming regions. Finally, the investigation of a large sample of clusters in Galactic giant H II regions will be important in understanding the massive star clusters in the Galactic center (Cotera et al., 1996; Figer et al., 1999a, b) which may have formed under different conditions than are typical in the disk of our Galaxy (Morris & Serabyn, 1996).
W42 is located in the fourth Galactic quadrant at l,b $`=`$ 25.4, 0.2. Lester et al. (1985) determined W42 to be at the “near” kinematic distance (3.7 kpc for $`R_{}=`$ 8 kpc) and hence somewhat less luminous than earlier estimates (Smith et al., 1978). In this series of papers we shall follow a suggestion of Dr. Robert Kennicutt (private communication) that “giant” means that more than $`10^{50}`$ Lyman continuum ($`=`$ Lyc) photons are inferred to be emitted per second from the H II region. This is about ten times the luminosity of the Orion nebula and roughly the number emitted from the hottest single O3-type star. As these stars are not found in isolation, there is an implication that a “giant” H II region contains some minimum of multiple O-type stars. In light of our new distance estimate (see §4), W42 probably falls below this limit being perhaps a few times more luminous than the Trapezium in Orion. Our target list was originally based on eleven of the most luminous giant H II regions from the study of (Smith et al., 1978) who tabulated their Lyc output derived from radio continuum measurements and kinematic distance estimates. W42 is not seen in visual images, and we estimate a foreground extinction of $`A_V`$ $`=`$ 10 mag (see below).
## 2 OBSERVATIONS AND DATA REDUCTION
$`J`$ ($`\lambda 1.3`$ $`\mu `$m, $`\mathrm{\Delta }\lambda 0.3`$ $`\mu `$m), $`H`$ ($`\lambda 1.6`$ $`\mu `$m, $`\mathrm{\Delta }\lambda 0.3`$ $`\mu `$m), and $`K`$ ($`\lambda 2.2`$ $`\mu `$m, $`\mathrm{\Delta }\lambda 0.4`$ $`\mu `$m) images of W42 were obtained on the nights of 29 August 1998 and 01, 02 May 1999 with the f/14 tip–tilt system on the Cerro Tololo Interamerican Observatory (CTIO) 4m Blanco telescope using the two facility imagers CIRIM (1998 data) and OSIRIS, the Ohio State InfraRed Imager/Spectrometer<sup>1</sup><sup>1</sup>1OSIRIS is a collaborative project between the Ohio State University and CTIO. OSIRIS was developed through NSF grants AST 9016112 and AST 9218449. (1999 data). Spectroscopic data were obtained on the night of 02 May 1999 using the f/14 tip–tilt system at the Blanco telescope with OSIRIS. CIRIM and OSIRIS are described in the instrument manuals found on the CTIO web pages (www.ctio.noao.edu). For OSIRIS, see also DePoy et al. (1993). The tip-tilt system is described by Pérez & Elston (1998). The tip–tilt system uses three piezo–electric actuators to move the secondary mirror at high frequency in a computer controlled feed–back loop which corrects the natural image centroid motion. OSIRIS employs 0.16<sup>′′</sup> pixels and CIRIM 0.21<sup>′′</sup> pixels.
All basic data reduction was accomplished using IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories.. Each image/spectrum was flat–fielded using dome flats and then sky subtracted using a median combined image of five to six frames. For W42 itself, independent sky frames were obtained five to ten arcminutes south of the cluster. Standard stars used the median combination of the data for sky.
### 2.1 Images
The OSIRIS 1999 May images were obtained under photometric conditions and in $``$ 0.5<sup>′′</sup> to 0.6<sup>′′</sup> FWHM seeing (with the tip–tilt correction). Total exposure times were 270 s, 135 s, and 135 s at $`J`$, $`H`$, and $`K`$, respectively. The individual $`J`$, $`H`$, and $`K`$ frames were shifted and combined (Figure 1), and these combined frames have point sources with FWHM of $``$ 0.6<sup>′′</sup>, 0.7<sup>′′</sup>, and 0.6<sup>′′</sup> at $`J`$, $`H`$, and $`K`$, respectively. DoPHOT (Schecter et al., 1993) photometry was performed on the combined images. The flux calibration was accomplished using standards 9170 and 9172 from Persson et al. (1998) which are on the Las Campanas Observatory standard system (LCO). The LCO standards are essentially on the CIT/CTIO system (Elias et al., 1982), though color transformations exist between the two systems for redder stars. The standards were taken just before the W42 data and within 0.16 airmass of the airmass for W42; no corrections were applied for these small differences in airmass. Aperture corrections using 11 pixel radius apertures were used to put the instrumental magnitudes on a flux scale.
Stars brighter than about 10th magnitude are expected to be a few percent non-linear on the OSIRIS images. We have included the 1998 CIRIM data for such stars since the count levels for the CIRIM images were in the fully linear regime (no linearity correction needed). The zero point to the CIRIM photometry was determined by comparing the instrumental magnitudes of stars in common to the OSIRIS and CIRIM images. The CIRIM $`J`$band photometry includes a color correction term in order to make a transformation on to the CIT/CTIO system (PAPER I). The OSIRIS data was placed on the CIT/CTIO system by making small corrections ($`<`$ 10 $`\%`$) to the $`J`$, $`H`$ and $`K`$ magnitudes based on linear fits to the magnitude differences as a function of color for stars in common to the CIRIM and OSIRIS images.
Uncertainties for the final $`JHK`$ magnitudes include the formal DoPHOT error added in quadrature to the error in the mean of the photometric standards (including the transformation to OSIRIS magnitudes for the CIRIM data), and the error in aperture corrections used in transforming from the DoPHOT photometry to OSIRIS magnitudes. The latter errors dominate and were derived from the scatter in the measurements of four to seven relatively uncrowded stars on the mosaic frames. The sum (in quadrature) of the aperture correction and standard star uncertainties is $`\pm `$ 0.018, $`\pm `$ 0.023, $`\pm `$ 0.019 mag in $`J`$, $`H`$, and $`K`$, respectively. The DoPHOT errors ranged from approximately $`\pm `$ 0.01 mag to an arbitrary cut–off of 0.2 mag (stars with larger errors were excluded from further analysis).
The flat–field illumination was not uniform. A smooth gradient with full range of about 10$`\%`$ was present. Corrections for this gradient were made based on observations of a standard star taken over a 49 position grid covering the array.
### 2.2 Spectra
The spectra of three of the brightest four stars in the center of W42 were obtained with a 0.48<sup>′′</sup> wide slit (oriented EW) in $`<`$ 0.7<sup>′′</sup> FWHM seeing and divided by the spectrum of HR 6813 (A1V) to remove telluric absorption features. Br$`\gamma `$ absorption in HR 6813 was removed by eye by drawing a line across it between two continuum points. One dimensional spectra were obtained by extracting and summing the flux in a $`\pm `$ 2 pixel aperture (0.64<sup>′′</sup> wide). The extractions include background subtraction from apertures centered $`<`$ 1.0<sup>′′</sup> on either side of the object.
The wavelength calibration was accomplished by measuring the positions of bright OH<sup>-</sup> lines from the $`K`$band sky spectrum (Olivia & Origlia, 1992). Lines are identified by their relative differences between one and another. The measured dispersion is 0.0003683 $`\mu `$m pix<sup>-1</sup>. The spectral resolution at 2.2 $`\mu `$m is $`\lambda /\mathrm{\Delta }\lambda `$ 3000.
## 3 RESULTS
A spectacular stellar cluster is revealed at the heart of W42 in our near infrared images (Figure 1 and Figure 2). Apparently, the cluster has just emerged from the edge of the molecular cloud from which it formed. This is confirmed below through $`K`$band spectroscopy which shows that the central massive star has largely cleared away its birth cocoon, but two of the next brightest stars have not. Lester et al. (1985) observed W42 in the mid and far infrared. Their analysis revealed two distinct sites of star formation toward W42. One, G25.4SE is located at 3.7 kpc from the sun (accounting for a sun to Galactic center distance of 8 kpc), while the other, G25.4NW is probably located at about 9.6 kpc. Lester et al. found a 10 $`\mu `$m unresolved source located at RA (2000) $`=18^\mathrm{h}38^\mathrm{m}15^\mathrm{s}.2`$, Dec (2000) $`=06`$$`47^{}50^{\prime \prime }`$ which is coincident with G25.4SE. This mid IR source is associated with the stellar cluster we have observed. Comparison of the position of the bright foreground star in the SE corner of our $`K`$band image to the same star on the image from the Digitized Sky Survey (DSS) <sup>3</sup><sup>3</sup>3 Based on photographic data obtained using The UK Schmidt Telescope. The UK Schmidt Telescope was operated by the Royal Observatory Edinburgh, with funding from the UK Science and Engineering Research Council, until 1988 June, and thereafter by the Anglo-Australian Observatory. Original plate material is copyright (c) the Royal Observatory Edinburgh and the Anglo-Australian Observatory. The plates were processed into the present compressed digital form with their permission. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166. results in a position for the ($`K`$band) bright central star in the cluster (W42 #1; see §3.2) of RA (2000) $`=18^\mathrm{h}38^\mathrm{m}15^\mathrm{s}.3`$, Dec (2000) $`=06`$$`47^{}58^{\prime \prime }`$.
### 3.1 Images
The $`HK`$ color$``$magnitude diagram (CMD) for the region toward W42 is shown in Figure 3. A cluster sequence is evident at 0.6 $`<`$ $`HK`$ $`<`$ 1.5 along with stars with much redder colors and a probable foreground sequence. The $`JH`$ vs. $`HK`$ color–color diagram is presented in Figure 4. As expected from the morphology and range of colors in Figure 2, the strong effects of differential reddening can be seen in the color–color diagram. Typical reddening lines are shown for M giants (Frogel et al., 1978) and early O stars (Koorneef, 1983a). The latter color was transformed as described in §3.1.3. The cluster stars defined below in §3.1.1 and shown as open circles in Figure 4 may have a slight unexplained systematic offset relative to the reddening line for normal O stars. This possible offset will not affect the conclusions in this paper regarding the cluster stars. The relationship for the intrinsic colors of classical T Tauri stars (pre–main sequence stars) (Meyer et al., 1997) is also shown for reference. The adopted reddening law is from Mathis (1990).
#### 3.1.1 The Central Cluster
By separating the stars in Figure 3 based on radial position, we can better define the central cluster relative to the surrounding field. In Figure 5, we plot the radial surface density of stars centered on the position of the bright central star, W42 #1 (§3.2). We also plot the radial surface density for stars with $`K`$ $``$ 14 mag, for which the number counts are more nearly complete (see below). The surface density becomes approximately uniform at a radius of $`R30^{\prime \prime }`$, continuing up to 50<sup>′′</sup>. It then begins to fall rapidly at the edge of the array (dashed vertical line) as expected due to the rectangular shape of the field. Taking these radii as representative, we divide the CMD into regions with 0$`{}_{}{}^{\prime \prime }R<`$ 30<sup>′′</sup>, 30$`{}_{}{}^{\prime \prime }R<`$ 50<sup>′′</sup>, and $`R>`$ 50<sup>′′</sup>, as shown in Figure 6. These represent the central cluster, a background annulus, and the region of the array where edge effects are important to the radial number counts.
The range of color and brightness in Figure 6 overlaps for stars in the cluster and background regions. A distinct blue sequence can be seen in the 0$`{}_{}{}^{\prime \prime }R<`$ 30<sup>′′</sup> cluster CMD which appears to merge smoothly with other stars in the 30$`{}_{}{}^{\prime \prime }R<`$ 50<sup>′′</sup> CMD. Clearly the cluster can not be completely extracted from the surrounding field based on radial position alone. In order to further enhance the cluster sequence, we defined a background CMD using the 30$`{}_{}{}^{\prime \prime }R<`$ 50<sup>′′</sup> region and accounting for the area difference between this annulus and the central region. The background CMD was binned in 0.5 mag color–magnitude bins. We then randomly selected and subtracted stars from the cluster in equal numbers from bins matching the background CMD. The resulting CMD is shown in Figure 7. The corresponding colors for these stars are plotted as open circles in Figure 4.
#### 3.1.2 Extinction to the Cluster
Most of the brightest stars in the central cluster fall along an essentially vertical track as expected for hot stars. The average color of the brightest seven of these stars is $`HK`$ $`=`$ 0.637 which corresponds to an extinction at 2.2 $`\mu `$m of $`A_K=`$ 1.07 mag ($`A_V`$ $``$ 10 mag) using the interstellar reddening curve of Mathis (1990) and an intrinsic $`HK`$ $`=`$ $`0.04`$ (Koorneef, 1983a). Many other stars appear more reddened, typically up to $`HK`$ $``$ 2 with some as high as 3.8. A star with $`HK`$ $`=`$ 2 and intrinsic $`HK`$ $`=`$ 0 would have $`A_K=`$ 3.2 mag ($`A_V`$ $``$ 32 mag).
#### 3.1.3 $`K`$ vs. $`HK`$ CMD for the ZAMS
It will be useful in discussing the W42 cluster below to have an estimate for the zero age main sequence (ZAMS) transformed into the $`K`$ vs. $`HK`$ plane. We have constructed such an estimate using the model results of Schaller et al. (1992). The models give bolometric luminosities ($`M_{\mathrm{bol}}`$) and effective temperatures ($`T_{\mathrm{eff}}`$) for stars of a given mass as they begin their evolution on the ZAMS. Using relationships for spectral type vs. $`T_{\mathrm{eff}}`$ (Vacca et al., 1996; Johnson, 1966), visual bolometric correction (BC<sub>V</sub>) vs. effective temperature (Vacca et al., 1996; Malagnini et al., 1986), $`VK`$ vs. spectral type (Koorneef, 1983a), and $`HK`$ vs. spectral type (Koorneef, 1983a), we transform the $`M_{\mathrm{bol}}`$ and $`T_{\mathrm{eff}}`$ to the $`K`$ vs. $`HK`$ CMD. A small correction has been made in an attempt to place the $`HK`$ colors from Koorneef (1983a) onto the CIT/CTIO system. The Koorneef (1983a) colors are basically referred to the system of Johnson (1966) with newly defined colors involving $`H`$ (Johnson used no $`H`$ filter). In addition, Koorneef (1983b) found that interpolating an $`H`$ magnitude from the observed $`J`$ and $`K`$ values gives a very good estimate of the observed $`H`$ magnitude. Carter (1990) found no discernible color correction between the SAAO system and that of Johnson (1966). We have therefore used the transformation of Carter (1990) between the SAAO system and CIT/CTIO to transform the $`HK`$ colors given by Koorneef (1983a). These corrections are at most one percent and hence negligible compared to the measurement uncertainties and scatter due to differential reddening.
The ZAMS CMD is shown in Figure 7 for a particular distance (2.2 kpc) which is discussed in §4. The values plotted in Figure 7 are listed in Table 1. These values may be compared to those given by Hanson et al. (1997). There is generally good agreement as a function of $`T_{\mathrm{eff}}`$. This is expected since the same models and colors were used in both cases. There are systematic differences at the $``$ $`\pm `$ 0.3 mag level for $`M_K`$ as a function of spectral type since the spectral type vs. $`T_{\mathrm{eff}}`$ relation adopted here (Vacca et al., 1996; Johnson, 1966) is different than that used by Hanson et al. (1997) who used the relation given by Massey et al. (1989). These differences do not affect the conclusions of §4.3 since the resulting $`M_K`$ for the spectral type in question (O5–O6) is within 0.1 mag in both cases.
#### 3.1.4 The $`K`$band Luminosity Function
The $`K`$band luminosity function (KLF) is shown in Figure 8. Neglecting the last five bins where the counts appear to be incomplete and the first 3 bins where the counts may be better described as uniform, the $`K`$band counts are well fit by a power–law with index 0.4 (log<sub>10</sub>N $`=`$ 0.4 $`\times `$ $`K`$ \+ constant).
We estimated the completeness of the KLF by performing artificial star experiments. Initially, we attempted to add stars to the original $`K`$band image, but even adding 10$`\%`$ of the observed KLF resulted in recovered luminosity functions much less complete than the observed one. This is due to the high spatial density of the cluster. To avoid this problem, we constructed complete artificial frames and analyzed them in the same way as the original frame. Anticipating the result, we used the spatial density distribution of the stars on the $`K`$band image with $`K<`$ 14 mag (Figure 5) to generate the positions of stars on the artificial frames. The input luminosity function was constructed from two components, a uniform distribution from 8.5 $`K`$ 10.5, and a power–law for 10.5 $`<K`$ 19.5. The total number of stars was set by the flux on the $`K`$band image, less a uniform background component determined from the sky frame. This should be conservative regarding the number of stars since the strong nebular flux is included. This latter aspect may be balanced somewhat by the fact that the artificial frames did not include variable extinction and nebular emission. The total $`K`$band flux in Figure 1$`c`$ is 6.0 mag, not including the bright saturated star in the SE corner of the frame. The total flux of the average actual input luminosity function (Figure 9) was 5.9 mag which can be compared to the total flux in observed stars (Figure 8), $`K=`$ 6.18 mag.
We constructed 10 artificial frames by randomly sampling the spatial distribution and luminosity functions. DoPHOT was run on each frame and the recovered stars matched to the input lists. The average recovered luminosity function, input luminosity function, and completeness fraction are shown in Figure 9. The general shape and distribution of the recovered stars in Figure 9 suggest the experiments are a fair test of completeness of the original frame. From the same input and recovered star lists, we also show the luminosity functions and completeness fraction for the stars located at 0$`{}_{}{}^{\prime \prime }R<`$30<sup>′′</sup>. See Figure 10.
Tests were also made to check whether somewhat steeper power–laws were also consistent with the observed KLF. This might be the case if the crowding were sufficient to “hide” many fainter stars. Artificial star experiments analogous to those above (including the same uniform component) but with a power–law component of 0.5 are not consistent with the data. A single power–law component of this steepness produces more star light than is necessary to account for the observed total flux and produces too many stars in the recovered luminosity function at fainter magnitudes if it is required to produce approximately the correct numbers at brighter magnitudes.
### 3.2 Spectra
The spectra of the targets W42 #1, #2, and #3 are shown in Figure 11. The final signal–to–noise in the three spectra is typically 80–105, 75–95, and 55–78 for W42 #1, W42 #2, and W42 #3, respectively and is higher on the red end than the blue end for all three. The brightest star in the central cluster, W42 #1 ($`K=`$ 8.8 mag), shows characteristic O star features (Hanson et al., 1996). These include C IV (2.069 $`\mu `$m and 2.078 $`\mu `$m) emission, N III (2.1155 $`\mu `$m) emission, Br$`\gamma `$ (2.1661 $`\mu `$m) absorption, and He II (2.1891 $`\mu `$m) absorption. Comparison to the standards presented by Hanson et al. (1996) results in $`K`$band spectral type of kO5-O6. These stars typically have MK spectral types of O5 to O6. The present classification system laid out by Hanson et al. (1996) does not have strong luminosity class indicators. Still, the He I (2.06 $`\mu `$m) and Br$`\gamma `$ features can be used to approximately distinguish between dwarf+giants on the one hand, and supergiants on the other: the supergiants tend to have emission or weak absorption in these lines. The spectrum of W42 #1 shown in Figure 11 has been background subtracted with nearby apertures to account for the nebular emission seen in projection toward the star. The apparent absorption feature at the position of Br$`\gamma `$ and absence of a feature at He I (2.06 $`\mu `$m) (which might indicate a poor subtraction of the nebular contribution), suggests that W42 #1 is a dwarf or giant star.
In contrast to W42 #1, the spectra of W42 #2 and #3 show only emission features at He I (2.06 $`\mu `$m) and Br$`\gamma `$. These spectra have also been background subtracted with nearby ($`<`$ 1.0<sup>′′</sup>) background apertures. We believe the strong emission remaining after this subtraction is related to the local environment in these stars (see §4.1).
## 4 DISCUSSION
The dense stellar cluster evident in Figure 1 surrounded by intense nebulosity leaves no doubt that this is a young object still emerging from its birth environment at the edge of its parent molecular cloud. The cluster appears to have emerged by clearing the foreground material to the West; darker regions with fewer stars remain toward the East. There is a suggestion of photoevaporated regions on the edge of the cloud to the East similar to those seen in M16 (Hester et al., 1996) and NGC 3603 (Brander et al., 1999). This is particularly evident in Figure 1$`a`$. This picture is consistent with the CO line maps of Lester et al. (1985), who found the peak brightness temperature of the associated molecular cloud to be offset to the East of the W42 (G25.4SE) H II region (see their Figure 8). Higher spatial resolution images and in nebular lines are in order to further study the interaction and impact of the ionizing cluster on the molecular cloud interface.
The process of clearing the local environment is still on–going, clumps of dark material are seen in projection against the photoionized H II region (e.g. SW of the cluster center). In §3.1.2, the extinction to the brightest stars on the cluster sequence was found to be $`A_K`$ $`=`$ 1.07 mag. This value can be taken as representative of the foreground extinction to the cluster, and indicates a $`V`$ magnitude of approximately 18 for W42 #1. This is consistent with the non–detection of the cluster on the DSS plates. A few stars remain after subtracting the ”background” component from the CMD (§3.1.1), but these are more likely in the foreground since the extinction toward them is as low as $`A_K`$ $``$ 0.3 mag.
### 4.1 Young Stellar Objects
There are a host of very red objects indicated in Figures 1, 6, and 7. Such colors are suggestive of young stellar objects (YSO) in the context of an embedded stellar cluster emerging from its parent cloud. The observed colors of known populations of pre–main sequence stars (PMS) and more heavily embedded protostars appear to lie on/occupy rather well defined sequences/regions in observational color space; and their colors, which are redder than usual stellar photospheres have been successfully modeled as arising from excess emission produced by circumstellar disks (Lada & Adams, 1992; Meyer et al., 1997) and envelopes (Hartmann et al., 1993; Pezzuto et al., 1997). Like normal stars, PMS objects may be seen through foreground extinction, further reddening their colors. Deeply embedded protostars may have colors with or without excess emission, but typically suffer large extinction due to the dense envelope surrounding them (Lada & Adams, 1992). Obviously, protostars too may suffer additional foreground extinction.
The $`JH`$ vs. $`HK`$ color–color diagram is most useful in assessing whether any particular star may have excess emission. The diagram distinguishes between normal stellar colors which are seen through a column of dust, hence making them redder, and a contribution which is due to emission (reprocessed stellar light from circumstellar envelope and/or accretion luminosity from a circumstellar disk). Lada & Adams (1992) and Meyer et al. (1997) have shown that disks can provide a source of excess emission to the normal stellar colors. In Figure 4, we have plotted the classical T Tauri sequence (CTTS) from Meyer et al. (1997) along with reddening lines for the interstellar reddening law of Mathis (1990). The former locus may be understood as arising from disk luminosity and projection effects (e.g, Lada & Adams 1992 and references therein). The near infrared colors of weak lined T Tauri stars (WTTS) are similar to normal stars (Meyer et al., 1997). These stars exhibit spectra which suggest they have disks, but the disk contribution to the near infrared colors is negligible. The Herbig AeBe stars which are higher mass analogs to the CTTS show similar colors to the CTTS but extend to generally larger $`HK`$ excess and with a larger contribution to the colors from circumstellar extinction (Lada & Adams, 1992). The colors of these objects have been fit by circumstellar envelope models in which the excess emission arises from reprocessed stellar light via dust heating and re-radiation (Hartmann et al., 1993; Pezzuto et al., 1997), or alternatively, from circumstellar disks with central holes (Lada & Adams, 1992). We also show the reddening lines for a main sequence O star and M giant (see S3.1) which can be taken as approximate guides for the expected colors for normal stars along sight lines to the inner Galaxy. The black squares in Figure 4 represent all the stars in Figure 1 for which $`J`$, $`H`$, and $`K`$ colors were measured. The open circles represent the stars in the ”background subtracted” central cluster (§3.1.1).
The W42 cluster stars (including W42 #1) occupy a rather tight sequence at modest reddening in Figure 4 as expected for young massive stars seen through different columns of obscuring material. At larger reddening ($`HK`$ , $`JH`$ $`=`$ 1,2), there is a larger dispersion of colors. There are also a number of stars which lie significantly beyond the reddening band for normal stars in the region of PMS stars and protostars as discussed above. Two of these objects are W42 #2 and #3. Spectra for these objects are shown in Figure 11. Neither spectrum shows stellar absorption features, though both exhibit emission at Br$`\gamma `$ or He I (2.06 $`\mu `$m) after background subtraction (§2). Accounting for the foreground extinction to W42 #1, both W42 #2 and #3 lie in a region of the color–color plot occupied by luminous protostars and Herbig AeBe stars (Lada & Adams, 1992). Lada & Adams (1992) suggest that Herbig AeBe stars might be recently ”uncovered” protostars given the large overlap in the color–color plot for the two objects. We believe that the position of W42 #2 and #3 in the color–color plot combined with their emission–line spectra, strongly suggests they are also luminous YSOs. The absence of absorption features is consistent with veiling due to the excess emission seen in Figure 4. In the case of W42 #2, the residual Br$`\gamma `$ emission may be due to a circumstellar disk; spherical or envelope distributions to the circumstellar material are also possible. Higher spectral resolution data on the Br$`\gamma `$ line in this object will be required to rule out models of one nature or the other. In W42 #3, the presence of He I emission but no Br$`\gamma `$ could be due to imperfect subtraction. W42 #3 has closer neighbors making the subtraction of the background more difficult. If the emission seen in projection toward W42 #3 is due to a compact H II region, then the He I emission may be coming from a region closer to the star than the Br$`\gamma `$, the latter naturally subtracting off better on these angular scales. We plan to obtain $`H`$band spectra of W42 #2 and #3 where the excess emission should be less allowing for an improved picture of the nature of the embedded objects which give rise to the emission spectra seen at $`K`$.
There are several other stars with similar or slightly lower brightness indicated in Figure 7, but with apparently normal colors for hot stars. This suggests a mixture of massive objects with and without circumstellar (possibly disk) signatures as seen in M17 (Hanson et al., 1997). Like M17, the YSOs indicate a very young age for the cluster.
The situation in W42 may be compared to other stellar clusters in giant H II regions. For M17, Lada et al. (1991) report that the vast majority of cluster stars they studied have infrared excesses. This result should be confirmed with higher spatial resolution images; the photometry presented by Lada et al. (1991) was performed by summing the flux in 4.8<sup>′′</sup> diameter apertures on their 0.8<sup>′′</sup> pix<sup>-1</sup> images. Lada et al. (1991) note that their data set was analyzed with profile fitting by Gatley, et al. (1991) who obtain similar results, and Hanson et al. (1997) clearly show that at least some of the massive stars in M17 have disk–like spectroscopic features. Gatley, et al. (1991) find far fewer stars with excess emission in the Orion Trapezium cluster and attribute this to a mixture of older and younger stars. However, Zinnecker et al. (1993) reach a different conclusion regarding the ages of the Orion cluster stars, and we will discuss this further in the next section. In W43 (Paper I), only a handful of objects appear in the excess emission region of the color–color plot and none of the three spectroscopically classified hot stars do. The brightest object in W43 exhibits Wolf–Rayet features in its $`K`$band spectrum suggesting an older age relative to W42. This is similar to NGC 3603, R136, and the Arches (see the discussion in Paper I and references therein), and the implication is that while such stars may still be core hydrogen burning, they are not on the ZAMS.
These comparisons need to account for the fact that the extinction is generally greater in W42 and (much so) in W43, so that many objects are not detected at $`J`$ or $`H`$ and hence do not appear in the color–color plot. Clearly, high spatial resolution, homogeneous data sets, each analyzed in detail for for completeness at $`J`$, $`H`$, and $`K`$ would be useful in assessing the intrinsic fractions of stars with excess emission in embedded clusters in giant H II regions.
Without spectra, it is not possible to identify the nature of the remaining objects in the color excess region of Figure 4, though in the context of this young cluster, it is likely that some are PMS stars or protostars. It is also clear from Figure 4 that some YSOs can occupy the same region of color space as normal stars. Some of the cluster stars in Figure 4 near ($`HK`$ , $`JH`$ $`=`$ 1,2) where the dispersion in color is larger might be a combination of WTTS with essentially normal stellar colors and CTTS/Herbig AeBe stars with smaller excess emission, the actual nature depending also on the luminosity of any particular object. Alternately, these could be normal main sequence stars. Finally, some could be normal K or M giants in the background, though it is unlikely that many are, given the large number of stars in close association with the central region of the cluster and the fact that we have subtracted a background component from the CMD.
Inspection of Figures 3 and 4 indicates that W42 #3 has varied in brightness and color between the 1998 and 1999 observations presented here, becoming fainter at $`K`$ and redder. Color and brightness variations in YSOs can be explained by a combination of changes in the accretion luminosity produced by the circumstellar disk and changes in the extinction toward the source. Skrutskie et al. (1996) monitored a sample of YSOs for color changes, finding the slopes of the colors (measured over day to year timescales) in the $`JH`$ vs. $`HK`$ color–color diagram were intermediate between the CTTS locus (excess color arising from disk luminosity) and pure variations along the interstellar reddening lines such as might occur when clumps of obscuring material pass in front of the line of sight to the YSO. Skrutskie et al. (1996) found some YSO variations were consistent with one or the other effect dominating. The color variation for W42 #3 appears most consistent with a change in the local extinction.
### 4.2 The $`K`$band Luminosity Function
The KLF has been used as the basis for mass function determinations in young embedded or obscured massive star clusters. Lada et al. (1991) found that the M17 cluster KLF had a slope consistent with the Salpeter (1955) initial mass function (IMF) value if the M17 stars were normal main sequence type, but noting that this was difficult to reconcile with their finding that essentially all the M17 cluster stars had infrared excesses. Only in the case where such excesses arise in passive disks could their mass function determination still follow from the KLF. Gatley, et al. (1991) present a KLF for the Trapezium cluster in Orion (M42). They found the slope of the KLF to be inconsistent with that for the Salpeter (1955) initial mass function (IMF), hypothesizing that there are substantial numbers of older stars in M42 as well as very young ones. However, Zinnecker et al. (1993) found that the KLF in M42 was consistent with a population of only very young stars if PMS evolutionary tracks were used instead of assuming the cluster stars were on the main sequence. Zinnecker et al. (1993) show that deuterium burning on the PMS can cause peaks in the KLF which are a function of age thus producing luminosity functions different from expected if main sequence mass–luminosity relations are assumed. More recently, the detailed optical imaging/spectroscopic investigation of Hillenbrand (1997) exquisitely details the young main sequence and PMS stellar population in the Orion Trapezium cluster, clearly demonstrating the overall youth of the cluster.
Figer et al. (1999b) presented a mass function determination for the Arches cluster located near the Galactic center (GC). They find that the Arches cluster may have upwards of 100 O stars and an age of $``$ 2 Myr. Using their observed KLF and a mean value for $`A_K`$, these authors determined the mass function by relating the $`K`$ magnitudes to stellar mass using the Meynet et al. (1994) stellar evolutionary models. Figer et al. (1999b) find a mass function in the Arches which is significantly flatter than Salpeter (1955). They attribute this result to the different pre–conditions for star formation in the central few hundred parsecs of the Galaxy which they argue should favor higher mass stars. This is in contrast to the situation elsewhere in the Galaxy and in the Large Magellanic Cloud (LMC) (Massey & Hunter, 1998). Massey & Hunter (1998) find from a combination of spectroscopy and imaging that the IMF is similar for OB associations and dense massive star clusters (most notably R136 in the LMC) and in agreement with a Salpeter (1955) IMF.
The preceding discussion indicates that care must be taken in transforming the KLF into a mass function, including the effects of PMS evolution. In this sense, the determination of a mass function for W42 is premature given the implied young age and lack of transformations between PMS models and near infrared colors (which should include the effects of associated disk emission). However, a comparison to the different KLFs is warranted.
In §3.1 we derived the KLF for the inner 30<sup>′′</sup> central cluster in W42 and showed that it should be complete to $`>`$ 80 $`\%`$ for $`K`$ 15 mag. The measured slope of the cumulative counts in the central 30<sup>′′</sup> KLF ($`K`$ 15, correcting for incompleteness) is 0.38 $`\pm `$ 0.016. If we consider only stars in the background subtracted CMD (Figure 7, the measured slope ($`K`$ 15) is 0.36 $`\pm `$ 0.011. The slope in our background annulus (30<sup>′′</sup> to 50<sup>′′</sup>, $`K`$ 15, correcting for incompleteness) is 0.51 $`\pm `$ 0.017. These values may be compared to the results of Lada et al. (1991), Zinnecker et al. (1993), and Figer et al. (1999b). For convenience , we compare the power–law slope to the cumulative KLF which Lada et al. (1991) also calculated for their data. In M17, Lada et al. (1991) report a slope of 0.26, which they claim is consistent with the Salpeter (1955) IMF. We have calculated approximate slopes for Zinnecker et al. (1993) and Figer et al. (1999b) from figures of the published luminosity functions; these are 0.39 $`\pm `$ 0.016 and 0.28 $`\pm `$ 0.013, respectively. Figer et al. (1999b) conducted completeness detailed tests and corrected the published counts. Lada et al. (1991) conducted more rudimentary checks on completeness and report the KLF for the magnitude range they believe to be complete. Zinnecker et al. (1993) and Gatley, et al. (1991) conducted no completeness tests using artificial stars, to the best of our knowledge, though Zinnecker et al. (1993) claim the KLF is complete. All the KLFs discussed here are defined for the central clusters in the respective star forming regions, and none have been corrected for reddening. Thus our comparison is only for the slopes of the KLFs and does not account for possible differences resulting from differential reddening.
The KLF in W42 appears to be more similar to that in Orion (i.e. relatively steep) than in M17 or the Arches. Given the likely presence of YSOs as discussed above, we may be seeing the effects of PMS evolution on the number counts as is the case for the Trapezium. Both the Arches and M17 clusters clearly have more massive stars than are present in the central cluster of W42 or the Trapezium (Cotera et al., 1996; Hanson et al., 1997). The apparent KLFs in these latter two clusters appear at least superficially similar, though Figer et al. (1999b) and Lada et al. (1991) reach rather different conclusions when converting to mass functions. Figer et al. (1999b) derive an IMF which is much flatter ($``$ factor of 2) than Salpeter (1955), while Lada et al. (1991) find a mass function which is possibly consistent with Salpeter (1955).
### 4.3 Distance
In §3.2 we classified the spectrum of the bright central star, W42 #1, as an early to mid O type ($``$ O5–O6). Several lines of evidence presented above suggest that the W42 cluster is quite young. If we take W42 #1 to be on the ZAMS then its apparent brightness would give a distance to the cluster of 2.2 kpc, considerably closer than the radio distance (3.7 kpc, see §3). In §3.2, we argued that the spectrum of W42 #1 was most similar to those for the dwarf or giant luminosity classes. Using the average $`M_V`$ from Vacca et al. (1996) for dwarf stars and colors from §3, gives 2600$`{}_{}{}^{1000}{}_{700}{}^{}`$ kpc. If the giant star $`M_V`$ from Vacca et al. (1996) is used, then the distance estimate becomes 3400$`{}_{}{}^{1200}{}_{900}{}^{}`$ pc. The uncertainty in the distance estimate is completely dominated by the luminosity class assumed and the scatter in the intrinsic $`M_V`$ of O stars ($`\pm `$ 0.67 mag). The uncertainties in the reddening and apparent magnitude are negligible in comparison: a few percent for $`m_K`$ and $`<`$ 10 $`\%`$ for $`A_K`$, where the largest part of the uncertainty is in the choice of reddening law (see Mathis 1990). These distance estimates are in agreement with Lester et al. (1985) who argued that W42 (G25.4SE) was at the near distance given by the radio recombination line velocity and Galactic rotation model. Depending on the true luminosity of W42 #1, the cluster may even be somewhat closer. Smith et al. (1978) estimated the Lyc luminosity of W42 to be 8.2$`\times 10^{50}s^1`$ assuming the far kinematic distance (13.4 kpc). Adopting the ZAMS distance (2.2 kpc) as indicated by the young nature of the cluster (stars W42 #2 and #3) considerably reduces the expected ionizing flux to 2$`\times 10^{49}s^1`$.
In Figure 7, we have plotted a zero age main sequence (ZAMS) derived from the models of Schaller et al. (1992) (Z $`=`$ 0.02); see §3.1.3. We have placed the ZAMS in Figure 7 at 2200 pc, assuming that the W42 #1 is on the ZAMS. In this case, the (incomplete) cluster sequence reaches early K type stars. For the Trapezium (Zinnecker et al., 1993) claim that the KLF shows an intrinsic peak well above the sensitivity of their $`K`$band images. This peak is due to the PMS evolution of the lower mass stars and is sensitive to age on the PMS. For the close distance to W42 implied if W42 #1 is on the ZAMS, the similarity found above for the KLF of W42 compared to the Trapezium then suggests a similar effect on the KLF may be at work. I.e., the fainter magnitudes may correspond to lower mass PMS stars. In the present case, we are not claiming to see a real turn over in the KLF (as Zinnecker et al. do in Orion) since this occurs where the $`K`$band counts are demonstrably incomplete due to crowding.
For a distance of 2.2 kpc, the cluster has a radius of 0.32 pc at 30<sup>′′</sup>. The surface density in the inner 10<sup>′′</sup> (Figure 5) is then within a factor of two for that in the Trapezium (McCaughrean & Stauffer, 1994, $``$ 6500 pc<sup>-2</sup>) not accounting for the substaintial incompleteness in the present case. W42 is thus quite dense, and this is another indication of its young age (McCaughrean & Stauffer, 1994; Figer et al., 1999b).
The ZAMS clearly demonstrates the lack of any sensitivity to temperature in the hot stars at near infrared wavelengths. This means that traditional methods of cluster fitting (in the observational color–magnitude plane) for age or distance can not be used and only through suitable numbers of infrared spectra (resulting in $`M_{\mathrm{bol}}`$ and $`T_{\mathrm{eff}}`$) will the cluster properties truly emerge. The lack of temperature sensitivity in the near infrared colors does, however, lead to accurate estimates of the foreground extinction.
## 5 SUMMARY
We have presented high spatial resolution $`J`$, $`H`$, and $`K`$ images of the massive star cluster at the heart of the giant H II region W42. Our $`K`$band spectra of three of the brightest four stars in the central 30<sup>′′</sup> of the cluster indicate a very young population. The brightest star W42 #1 is classified as kO5–O6 based on the system of Hanson et al. (1997). Such stars are typically associated with MK O5–O6.5 dwarfs. W42 #2 and W42 #3 show no stellar absorption features. This fact, combined with their position in the excess emission band of the $`JH`$ vs. $`HK`$ color–color plot, leads us to classify them as YSOs.
The KLF was computed and compared to other massive star clusters. The KLF in W42 appears more similar to that of the Trapezium than to the more massive clusters in M17 and the Arches. The steepness of the KLF in the Trapezium has been attributed to the PMS stars present there, and it is possible that a similar effect is present in W42 given our spectroscopic and imaging evidence for YSOs.
Our spectrum of W42 #1 confirms the results of Lester et al. (1985) that W42 (G25.4SE) can not be at the far kinematic distance. Earlier estimates of the Lyc output from this giant H II region must therefore be revised downward (by $``$ an order of magnitude). If W42 #1 is on the ZAMS, as we argue based on the presence of young stellar objects, then the cluster is even closer than the near kinematic distance (3.7 kpc). We estimate 2.2 kpc. In this case, W42 should not be considered a giant H II region as defined in §1. Our spectrum of W42 #1 is not sensitive to luminosity class, though it suggests a dwarf or giant classification.
PSC appreciates continuing support from the National Science Foundation. We wish to acknowledge the continuing excellent support received from the CTIO mountain staff, particularly from night assistants Hernán Tirado, Patricio Ugarte, and Alberto Zuñiga. |
warning/0001/hep-th0001197.html | ar5iv | text | # References
SU-ITP-00/04
hep-th/0001197
January 2000
A SMALL COSMOLOGICAL CONSTANT FROM
A LARGE EXTRA DIMENSION
Nima Arkani-Hamed<sup>a,b,</sup><sup>1</sup><sup>1</sup>1arkani@thsrv.lbl.gov, Savas Dimopoulos<sup>c,</sup><sup>2</sup><sup>2</sup>2savas@stanford.edu, Nemanja Kaloper<sup>c,</sup><sup>3</sup><sup>3</sup>3kaloper@epic.stanford.edu
and
Raman Sundrum<sup>c,</sup><sup>4</sup><sup>4</sup>4sundrum@leland.stanford.edu
<sup>a</sup> Department of Physics, University of California, Berkeley, CA 94530
<sup>b</sup>Theory Group, Lawrence Berkeley National Laboratory, Berkeley, CA 94530
<sup>c</sup> Department of Physics, Stanford University, Stanford, CA 94305
ABSTRACT
We propose a new approach to the Cosmological Constant Problem which makes essential use of an extra dimension. A model is presented in which the Standard Model vacuum energy “warps” the higher-dimensional spacetime while preserving $`4D`$ flatness. We argue that the strong curvature region of our solutions may effectively cut off the size of the extra dimension, thereby giving rise to macroscopic $`4D`$ gravity without a cosmological constant. In our model, the higher-dimensional gravity dynamics is treated classically with carefully chosen couplings. Our treatment of the Standard Model is however fully quantum field-theoretic, and the $`4D`$ flatness of our solutions is robust against Standard Model quantum loops and changes to Standard Model couplings.
The extremely small value of the cosmological constant poses the most severe naturalness problem afflicting fundamental physics (see Ref. for a review). The problem stems from the fact that in General Relativity, all forms of energy necessarily act as gravitational sources which curve spacetime. Generically, the energy density of the vacuum as represented by the cosmological constant, would yield an unacceptably high curvature. An old but still exciting idea for ameliorating this situation is to assume that we live in a fundamentally higher-dimensional spacetime, which is indeed greatly curved by vacuum energy . However, it is now possible that the curvature spills into the extra dimensions, which can be larger than the fundamental Planck length . This idea has never been realized in any way which does not require the same level of fine-tuning as in four dimensions, although the way the fine-tuning shows up does change in interesting ways . A very interesting recent approach to the problem has been explored in .
In this paper, we report on some limited progress in this direction. Working within a five-dimensional “brane universe” effective field theory in which the Standard Model (SM) is confined to a 3-brane, we will show that one can carefully choose the higher-dimensional gravitational dynamics so that the SM vacuum energy is effectively converted into a current which is carried off the brane into the “bulk”. The gravitational back-reaction of this current warps the higher-dimensional spacetime, but in a manner consistent with four-dimensional Poincare invariance of the vacuum solution and inconsistent with four-dimensional de-Sitter or Anti-de-Sitter symmetry. However, our vacuum solution necessarily contains a singular region parallel to the SM brane, signaling high sensitivity to short-distance gravity and the breakdown of effective field theory in that region. Nevertheless we will argue that the singular region may effectively cut off the higher dimensional spacetime as in (see also ), thereby rendering the gravitational dynamics macroscopically four-dimensional, with vanishing effective cosmological constant.
In this paper we will explicitly consider only classical (higher-dimensional) gravity, but with quantum SM effects included. The gravitational couplings must be chosen very precisely for the mechanism to work, although having chosen them the mechanism is completely stable under SM radiative corrections and changes in SM couplings. To be precise about the size of effects neglected by the truncation to classical gravity, it is useful to distinguish the contributions to the cosmological constant which do and do not involve quantum gravity, whose fundamental scale is $`M_{}`$. In an effective field theory framework valid up to an energy $`\mathrm{\Lambda }<M_{}`$, we can expand the cosmological constant as
$$\lambda =𝒪(\mathrm{\Lambda }^4)+𝒪(\frac{\mathrm{\Lambda }^6}{M_{}^2})+\mathrm{}.$$
(1)
The first term obtains a contribution from the SM vacuum energy, treated fully quantum field theoretically, but is insensitive to quantum gravity corrections. The subsequent terms are sensitive to quantum gravity effects, described in Feynman diagrams by virtual graviton exchanges suppressed by powers of $`1/M_{}^2`$. Note that the formally larger first term is under the best theoretical control, while the subleading terms are shrouded to some extent by the mysteries of quantum gravity. Our truncation limits us to understanding the screening of the $`𝒪(\mathrm{\Lambda }^4)`$ SM vacuum energy, deferring the problem of $`𝒪(\frac{\mathrm{\Lambda }^6}{M_{}^2})`$ effects. This does not mean that these subleading effects are phenomenologically negligible since the smallest we can take $`\mathrm{\Lambda }`$ is our experimental cutoff of $`𝒪`$(TeV), so that even $`\frac{\mathrm{\Lambda }^6}{M_{}^2}`$ is much larger than the observed cosmological constant. Nevertheless we consider it a significant advance to show how the effects of the SM vacuum energy alone can be nullified.
Our set-up is as follows. We consider a five-dimensional spacetime, or “bulk”, where the fifth dimension is a half-line with coordinate $`y`$, while the usual four dimensions have coordinates $`x^\mu `$. Orbifold boundary conditions as in will be used to realize the half-line in terms of a full real line with the identification of $`y`$ with $`y`$. The five-dimensional degrees of freedom will therefore be taken to be symmetric about $`y=0`$. We will consider the boundary of the half-line, $`y=0`$, to be the location of an “end-of-the-world” 3-brane, to which the SM quantum field theory (or some extension thereof) is confined. The SM degrees of freedom interact with classical five-dimensional bulk gravity as well as a classical five-dimensional scalar field $`\varphi `$. Similar models have been studied in different contexts in . In the bulk, the scalar field will only be coupled minimally to bulk gravity, and on the brane it will be conformally coupled to the SM. Similar couplings have been employed in $`4D`$ attempts to solve the cosmological constant problem . We will carefully choose the conformal coupling constant such that for a given value of the vacuum energy on the brane there exists a flat solution. However, as we will show, once this choice is made, the ensuing flat solution remains unaffected by quantum loops on the brane, showing that the vanishing of the $`4D`$ cosmological constant is both radiatively stable and insensitive to the values of the SM parameters. Up to two derivatives, the action of our model is
$$S=d^4x𝑑y\sqrt{g_5}\left(\frac{R}{2\kappa _5^2}\frac{3}{2}(\varphi )^2\right)d^4x\sqrt{det(g_4(0)e^{\kappa _5\varphi (0)})}_{SM}(H,g_{\mu \nu }(0)e^{\kappa _5\varphi (0)}).$$
(2)
The constant $`\kappa _5^2`$ is related to the $`5D`$ Planck scale $`M_{}`$ by $`\kappa _5^2=M_{}^3`$. Here $`M,N`$ run over the $`5D`$ coordinates, while $`\mu ,\nu `$ run over the $`4D`$ brane coordinates. Note the special normalization of the bulk scalar $`\varphi `$, or equivalently the special coupling to the brane for the canonically normalized scalar field. The SM Lagrangian $`_{SM}`$ is localized to the brane at $`y=0`$. It depends on the SM fields, $`H(x)`$ and any necessary ultraviolet regulators, all minimally coupled to the Weyl-rescaled induced metric $`g_{\mu \nu }(y=0)e^{\kappa _5\varphi (y=0)}`$. Note that we have set the bulk cosmological constant and the operators $`f(\kappa _5\varphi )R`$ equal to zero. With the Weyl transformation $`\overline{g}_{MN}=g_{MN}\mathrm{exp}(\kappa _5\varphi )`$ we can cast the action (2) into the “string frame”, where it takes the form
$$S=d^4x𝑑y\sqrt{\overline{g}_5}e^{3\kappa _5\varphi /2}\frac{\overline{R}}{2\kappa _5^2}d^4x\sqrt{\overline{g}_4(0)}_{SM}(H,\overline{g}_{\mu \nu }(0)).$$
(3)
It is now clear from this action that our choice of the conformal coupling is completely unaffected by quantum loops on the brane. The brane action, (including ultraviolet regulators), does not even contain the scalar field $`\varphi `$, so $`4D`$ general covariance guarantees that the SM renormalization will not change the form of the actions (3) and (2). For simplicity, we will use the canonical action (2) in what follows.
To incorporate the effects of SM quantum loops we can integrate out the SM degrees of freedom, and work in terms of the full 1PI effective action. Hence the brane action in (2) is replaced by the effective action,
$$S_{brane}\mathrm{\Gamma }_{eff}^{SM}(H,g_{\mu \nu }(0)e^{\kappa _5\varphi (0)}).$$
(4)
The equations of motion are now straightforward to obtain. Varying the action (2) and using (4) we find
$`R^{MN}{\displaystyle \frac{1}{2}}g^{MN}R`$ $`=`$ $`3\kappa _5^2\left(^M\varphi ^N\varphi {\displaystyle \frac{1}{2}}g^{MN}(\varphi )^2\right)+{\displaystyle \frac{2\kappa _5^2}{\sqrt{g_5}}}{\displaystyle \frac{\delta \mathrm{\Gamma }_{eff}^{SM}}{\delta g_{\mu \nu }}}\delta ^M{}_{\mu }{}^{}\delta _{}^{N}{}_{\nu }{}^{}\delta (y),`$ (5)
$`3\mathrm{}_5\varphi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{g_4}}}{\displaystyle \frac{\delta \mathrm{\Gamma }_{eff}^{SM}}{\delta \varphi }}\delta (y),`$ (6)
$`{\displaystyle \frac{\delta \mathrm{\Gamma }_{eff}^{SM}}{\delta H}}`$ $`=`$ $`0.`$ (7)
Equation (7) encapsulates the full SM quantum field theory in the curved background.
We now seek solutions with $`4D`$ Poincare symmetry. With this symmetry the SM effective action is given by just the effective potential,
$$\mathrm{\Gamma }_{eff}^{SM}=d^4x\sqrt{g_4}V_{eff}(H)e^{2\kappa _5\varphi },$$
(8)
and hence the SM equations of motion take the simple form
$$\frac{V_{eff}}{H}=0.$$
(9)
We will denote an extremal value of the potential by $`V_{extremal}`$. Further, the structure of the $`5D`$ metric is restricted by the Poincare symmetry to the form
$$ds^2=a^2(y)\eta _{\mu \nu }dx^\mu dx^\nu +dy^2,$$
(10)
where $`a(y)`$ is the warp factor. The equations of motion (5)-(6) then take a particularly simple form. Defining the shifted scalar field $`\stackrel{~}{\varphi }=\varphi +\frac{1}{2\kappa _5}\mathrm{ln}(\frac{V_{extremal}}{M_{}^4})`$, the field equations in the bulk become
$$\frac{a^2}{a^2}=\frac{\kappa _5^2\stackrel{~}{\varphi }^2}{4},\stackrel{~}{\varphi }^{\prime \prime }+4\frac{a^{}}{a}\stackrel{~}{\varphi }^{}=0,\frac{a^{\prime \prime }}{a}=\frac{3\kappa _5^2\stackrel{~}{\varphi }^2}{4},$$
(11)
while the $`\delta `$-function sources and the $`yy`$ symmetry imply the matching conditions for the first derivatives
$$a^{}(0)=\frac{M_{}}{6}e^{2\kappa _5\stackrel{~}{\varphi }(0)}a(0),\stackrel{~}{\varphi }^{}(0)=\frac{M_{}^{5/2}}{3}e^{2\kappa _5\stackrel{~}{\varphi }(0)}.$$
(12)
Here the prime refers to the derivative with respect to $`y`$. Note that with the definition of the variable $`\stackrel{~}{\varphi }`$, the gravitational field equations are completely independent of the extremum value of the SM effective potential $`V_{extremal}`$. This will be crucial for ensuring the success of our mechanism.
It is straightforward to solve the system (11)-(12). The explicit solutions with Poincare symmetry are
$`ds_5^2`$ $`=`$ $`(1{\displaystyle \frac{2M_{}}{3}}e^{2\kappa _5\stackrel{~}{\varphi }_0}|y|)^{1/2}\eta _{\mu \nu }dx^\mu dx^\nu +dy^2,`$
$`\stackrel{~}{\varphi }`$ $`=`$ $`\stackrel{~}{\varphi }_0{\displaystyle \frac{1}{2\kappa _5}}\mathrm{ln}\left(1{\displaystyle \frac{2M_{}}{3}}e^{2\kappa _5\stackrel{~}{\varphi }_0}|y|\right),`$
$`\varphi `$ $`=`$ $`\varphi _0{\displaystyle \frac{1}{2\kappa _5}}\mathrm{ln}\left(1{\displaystyle \frac{2V_{extremal}}{3M_{}^3}}e^{2\kappa _5\varphi _0}|y|\right),`$ (13)
where $`\varphi _0`$ is an integration constant. We have checked explicitly that the solutions with $`4D`$ Poincare symmetry are the only allowed solutions which are $`4D`$ maximally symmetric. Making a more general ansatz for the metric with the symmetries of $`4D`$ (Anti-) de-Sitter space with curvature $`h`$, it is straightforward to show that if the scalar field coupling to the brane is $`\mathrm{exp}(\zeta \kappa _5\varphi )`$, the $`4D`$ curvature is
$$h^2=a^2(0)\frac{\kappa _5^2\stackrel{~}{\varphi }^2(0)}{4}a^2(0),$$
(14)
which vanishes by virtue of our choice of the conformal coupling coefficient $`\zeta =2`$ in eq. (2). This coupling prohibits both de-Sitter symmetry and Anti-de-Sitter symmetry on the brane, regardless of the details of the SM physics. The solutions (13) give a flat Minkowski space, for any value of $`V_{extremal}`$ even after all quantum corrections are included. Further, altering the parameters of SM, such as for example the electron mass and the fine structure constant, will not break the Poincare symmetry.
The fact that we have found a solution with $`4D`$ Poincare invariance, robust against SM loops and couplings, does not complete our task. For instance, we can always find $`4D`$ flat solutions for a 3-brane carrying a range of tensions in $`6D`$, the tension just inducing a deficit angle in the bulk. The problem is that gravity remains six-dimensional at long distances. Therefore, we must ask whether our $`5D`$ set-up leads to $`4D`$ gravity at long distances. If it does, then our Poincare invariant solution demonstrates that the effective $`4D`$ gravitational dynamics has vanishing cosmological constant.
An important feature of our solution which gives hope for ensuring macroscopic $`4D`$ gravity is the appearance of a naked curvature singularity at finite proper distance,
$$y_s=\frac{3}{2M_{}}e^{2\kappa _5\stackrel{~}{\varphi }_0}=\frac{3M_{}^3}{2V_{extremal}}e^{2\kappa _5\varphi _0},$$
(15)
away from the brane. Similar singularities have appeared in the work of . The appearance of our singularity can be understood by noting the remarkable analogy between our equations and Einstein’s equations in cosmological Friedmann-Robertson-Walker (FRW) spacetimes (with four spatial dimensions), if we interpret the coordinate $`y`$ as a cosmological “time” and the warp factor as the FRW “scale factor”. The bulk equations (11) then coincide with the cosmological equations of the FRW universe dominated by a massless scalar field. This should not be a surprise, since the ansatz of (10) is a Wick rotation of the FRW ansatz. The analogy with cosmological dynamics immediately shows that away from the brane only two possibilities can occur. One possibility is that the warp factor $`a`$ monotonically increases forever and the scalar field dissipates away, which means that the topology of the extra dimension is not compact and manifestly has infinite volume. This case however requires negative energy on the brane. Our solution follows from $`V_{extremal}>0`$, namely the warp factor monotonically decreases to zero at a finite distance from the brane, where the curvature and the scalar field diverge. This FRW analogy can be thought of in two ways related by time reversal: either the brane specifies initial conditions which evolve into a “Big Crunch”, or a “Big-Bang” evolved with final conditions specified by the brane.
Clearly our solution cannot be trusted for $`y>y_s`$. A careful discussion of the singularity becomes crucial in deciding whether gravity becomes four-dimensional at distances larger than $`y_s`$. That we need to do this at all is already striking: the singularity of our solution forces the short-distance properties of quantum gravity to become relevant to whether we recover the inverse square law for gravity at long distances! We will assume that the singularity is smoothed out by the true short-distance theory of gravity. Since we do not know the details of this theory, we will not be able to make any rigorous claims about whether the resolution of the singularity does what we want. Nevertheless, we see that even without a detailed understanding of the physics which smooths out the singularity, there are two qualitatively different possibilities. The first is an analog of a “Big Crunch/Big Bang” transition. That is, the theory extends across the singularity into a region where a weakly coupled Einstein gravity description is valid again, and the warp factor $`a(y)`$ diverges as $`y\mathrm{}`$. In this case the singularity does not end space, and clearly there is no $`4D`$ gravity at long distances. The second and more attractive possibility is that space ends with a finite volume at the singularity (the analog of time starting at the Big Bang). For instance, short-distance effects might smoothly join the region of high curvature to a highly curved $`AdS_5`$ geometry; pushing the singularity to infinite proper distance while preserving a finite volume space. In this case, the extra dimension is effectively compactified in the manner of , and we should get $`4D`$ gravity at long distances. Alternatively, any form of spacetime description might become impossible beyond $`y_s`$. In either of these cases, the zero-mode tensor fluctuations
$`ds_5^2`$ $`=`$ $`(1{\displaystyle \frac{2M_{}}{3}}e^{2\kappa _5\stackrel{~}{\varphi }_0}|y|)^{1/2}\overline{g}_{\mu \nu }(x)dx^\mu dx^\nu +dy^2,`$ (16)
would correspond to a massless $`4D`$ graviton, with a finite $`4D`$ Planck scale,
$$M_{Pl}^2=M_{}^3_0^{y_s}a^2(y)=M_{}^2e^{2\kappa _5\stackrel{~}{\varphi }_0}=\frac{M_{}^6}{V_{extremal}}e^{2\kappa _5\varphi _0}.$$
(17)
Similarly the integration constant $`\stackrel{~}{\varphi }_0`$ is promoted to a $`4D`$ scalar zero-mode $`\stackrel{~}{\varphi }_0(x)`$<sup>1</sup><sup>1</sup>1Of course, a massless, gravitationally coupled scalar field is experimentally excluded. In a realistic theory, this scalar will have to pick up a mass of at least $``$(mm)<sup>-1</sup>..
Summarizing: we have found a solution with a 3-brane in $`5D`$, interacting with bulk gravity and a massless scalar. The bulk couplings are carefully chosen to guarantee $`4D`$ flat space solutions, while prohibiting other maximally symmetric solutions with $`4D`$ (Anti-) de Sitter symmetry. The solutions have a curvature singularity at a finite distance from the brane, analogous to the Big Bang or Big Crunch singularity of FRW cosmology. Just as it is possible to think of the crunch singularity as “ending time”, it is natural to assume that our singularity ends space, effectively compactifying the extra dimension and giving $`4D`$ gravity at long distances. If this assumption is correct, our model represents a partial solution to the Cosmological Constant Problem, where SM fine-tuning is avoided. Note that as the value of the SM vacuum energy, the warp factor adjusts itself to maintain the vanishing of the $`4D`$ cosmological constant. This is similar in spirit to the ideas of . Unlike the proposal of , in our model the SM fields reside on the “Planck” brane, rather than in the bulk.
But how can we be sure that the same short-distance physics which resolves the singularity does not upset our mechanism ensuring the appearance of a Poincare invariant solution, and hence vanishing $`4D`$ cosmological constant, without SM fine-tuning? A way of “resolving” the singularity which does not work is by simply ending space with a second brane. For instance, a standard orbifold compactification with a second brane cutting off the singular region with negative energy $`V_{neg}`$ does require the precise fine-tuning, $`V_{neg}=V_{extremal}`$ . Alternatively one can note that our bulk Lagrangian, with the special scalar couplings, arises precisely from compactifying one of the dimensions of a $`6D`$ theory, with the 3-brane coupling only to the induced metric, and $`\varphi `$ being the modulus of the extra dimension. It is straightforward to lift our $`5D`$ solution into the full $`6D`$ theory. The $`6D`$ space is compact and locally flat away from our brane, except at a conical singularity corresponding to our singularity at $`y_s`$. If this conical singularity is described from the outset as a brane in the theory, its tension must again be fine-tuned against the SM vacuum energy.
We will formulate the conditions for the short distance physics under which our mechanism works. First, we assume that in the absence of the SM 3-brane the full bulk dynamics admits $`4D`$ Poincare invariant solutions describing $`5D`$ spacetime ending in the singularity region. We will focus on these because they are the only ones compatible with the matching conditions on the brane. In the analog FRW picture, this amounts to saying that any spatially flat universe can end in a Big Crunch . We also assume that the values of $`\varphi `$ and $`\varphi ^{}`$ emerging from the brane are always compatible with the physics of the singularity. In the FRW analogy, this corresponds to saying that a few Planck times after the Big Bang, a scalar field without a potential can emerge with some $`\dot{\varphi }`$, but with any $`\varphi `$. The resolutions of the singularity discussed above failed because they violated these requirements. A sufficient but not necessary condition for the second assumption to be valid is if the short distance physics resolving the singularity is shift symmetric under $`\varphi \varphi +`$ constant. Also note that the contribution to the $`4D`$ cosmological constant from higher derivative bulk operators is Planck-suppressed relative to $`V_{extremal}`$ and therefore of the order we are already neglecting in this paper, as discussed in Eq. (1).
A physical picture of our mechanism follows from the observation that shift symmetry of the bulk action implies an associated conserved current,
$$J^M=\frac{1}{\sqrt{g_5}}\frac{\delta S_{bulk}}{\delta _M\varphi }.$$
(18)
$`4D`$ Poincare invariance implies that only $`J_y`$ can be non-zero. The fact that the $`\varphi `$-couplings to the SM on the brane explicitly break shift symmetry corresponds to a brane-localized source for $`J`$, thereby determining the current to be
$$J_y=3\varphi ^{}=M_{}^{5/2}e^{2\kappa _5\stackrel{~}{\varphi }}=\frac{V_{extremal}}{M_{}^{3/2}}e^{2\kappa _5\varphi }.$$
(19)
Thus we see that in our set-up the SM vacuum energy is converted into a current emerging from the 3-brane and ending in the singularity region. While the vacuum energy does not show up as $`4D`$ curvature, the gravitational backreaction of the associated current warps $`5D`$ spacetime and gives rise to the singularity into which it pours. As discussed, the singularity might be resolved into a highly curved $`AdS_5`$ geometry so that the current has room to continue infinitely far away from the brane, or possibly spacetime really ends at the singularity, in which case there may be some non-perturbative breakdown of shift symmetry which allows the current to end.
We conclude by making some comments on the long-distance $`4D`$ effective theory in our set-up. While we do not know the details of short-distance gravity and are forced to speculate on the nature of the singularity, the physics must yield a consistent $`4D`$ effective field theory. In particular, this theory must reflect the property of our mechanism that all the extrema of $`V(H)`$ lead to $`4D`$ flat solutions. Therefore, the naive guess that the effective theory simply contains the term
$$d^4x\sqrt{g_4}\left(V_{eff}(H)V_{extremal}\right)$$
(20)
cannot work because of the possibility of multiple extrema $`V(H_1),V(H_2),\mathrm{}`$. If we subtract the extremum from one vacuum, any other extremum would gravitate. If $`H`$ is the only light field in the theory, the Lagrangian (20) is the only one consistent with the flat space limit. However, the presence of other light gravitationally coupled fields $`\psi `$ offers a way out. For example, the true effective $`4D`$ potential may have the form $`V_{eff}(H,\psi )=F(\psi )V_4(H)G(\psi )`$. In the limit where $`M_{Pl}\mathrm{}`$, the $`\psi `$ decouple and the $`H`$ dynamics is governed by $`V_{eff}(H)`$. However, with finite $`M_{Pl}`$, it is possible for $`V_4(H,\psi )`$ to only have extrema at points $`(H_1,\psi _1),(H_2,\psi _2)`$ with vanishing potential, while $`(H_2,\psi _1)`$, $`(H_1,\psi _2)`$ are not extrema because they do not satisfy the $`\psi `$ equation of motion. It is straightforward to choose functions $`F`$ and $`G`$ with this property. In our model a natural candidate for the modulus is $`\varphi `$; furthermore, such a potential between $`\varphi `$ and the electroweak symmetry breaking sector can naturally induce a mass for $`\varphi `$ of order (TeV$`)^2/M_{Pl}`$(mm$`)^1`$ as required phenomenologically.
In this paper, we have presented a model with a $`3`$-brane in five dimensions, whose only maximally symmetric solutions are $`4D`$ Poincare invariant, independent of the SM parameters. Our solutions are forced into a strong curvature region which connects the fate of long distance gravity with its short distance properties. With specific assumptions on the nature of the singularity, we recover macroscopic $`4D`$ gravity with vanishing cosmological constant, in a manner consistent with a $`4D`$ effective field theory. A better understanding of the singularity would allow us to establish or exclude this idea.
Acknowledgements
We would like to thank Andy Cohen, Michael Dine, Bob Holdom, Joe Polchinski, Lisa Randall and Herman Verlinde for discussions. The work of N.A-H. has been supported in part by the DOE under Contract DE-AC03-76SF00098, and in part by NSF grant PHY-95-14797. The work of S.D., N.K. and R.S. has been supported in part by NSF Grant PHY-9870115.
Note added: While we were completing this paper, we were informed of the upcoming work which overlaps with the ideas presented here. |
warning/0001/cond-mat0001421.html | ar5iv | text | # Vortex lattice structures and pairing symmetry in Sr2RuO4
## 1 Introduction
The oxide Sr<sub>2</sub>RuO<sub>4</sub> has a structure similar to high $`T_c`$ materials and was observed to be superconducting by Maeno et al. in 1994 . It has been established that this superconductor is not a conventional $`s`$-wave superconductor: NQR measurements show no indication of a Hebel-Slichter peak in $`1/T_1T`$ , and $`T_c`$ is strongly suppressed below the maximum value of 1.5 K by non-magnetic impurities . More recent experiments indicate an odd parity gap function of the form $`𝐝(𝐤)=\widehat{z}(\eta _xk_x+\eta _yk_y)`$. The Knight shift measurements of Ishida et al. reveal that the spin susceptibility is unchanged upon entering the superconducting state; this is consistent with $`p`$-wave superconductivity (as predicted by Rice and Sigrist ). Furthermore, these measurements were conducted with the applied field in the basal plane. Since the orientation of the gap function (that is $`𝐝`$) is orthogonal to the spin projection of the Cooper pair , these measurements are consistent with the gap function aligned along the $`\widehat{z}`$ direction. The $`\mu `$SR experiments of Luke et. al. have revealed spontaneous fields in the Meissner state . This indicates that the superconducting order parameter must have more than one component . Naively, the observation of spontaneous fields leads to the conclusion that the superconducting gap function breaks time reversal symmetry ($`𝒯`$) and therefore must have the form $`𝐝(𝐤)=\widehat{z}(k_x\pm ik_y)`$. However, the muons probe inhomogeneities in the superconducting state since any bulk magnetic fields are screened in the Meissner state. Consequently, $`𝒯`$ need only be broken in the vicinity of the inhomogeneities. For example, it has been pointed out in Ref. that a gap function given by $`\stackrel{}{\eta }=(1,\pm 1)`$ can give rise to domain walls that break time reversal symmetry. In principle this can also lead to a $`\mu `$SR signal as seen by Luke et al.. Consequently, these two experiments lead to the gap function $`𝐝(𝐤)=\widehat{z}(\eta _xk_x+\eta _yk_y)`$, with the specific form of the order parameter $`\eta _x`$ and $`\eta _y`$ undetermined. In this talk, the vortex lattice structures arising in the phenomenological theory of this $`E_u`$ representation are examined. First, the theory is examined for an applied magnetic field in the basal plane along one of the two-fold symmetry axes. It is shown that multiple vortex lattice phases are generic to this representation. Then the weak-coupling limit of the theory is examined for the applied magnetic field along the $`c`$-axis. It is shown that a square vortex lattice results for parameters relevant to Sr<sub>2</sub>RuO<sub>4</sub>; consistent with the experimental observation of a square vortex lattice . The theory of the square vortex lattice is compared with the observed results and good agreement is found.
## 2 Free Energy
The dimensionless free energy for the $`E_u`$ representation of $`D_{4h}`$ is given by
$`f=`$ $`|\stackrel{}{\eta }|^2+|\stackrel{}{\eta }|^4/2+\beta _2(\eta _x\eta _y^{}\eta _y\eta _x^{})^2/2`$
$`+\beta _3|\eta _x|^2|\eta _y|^2+|\stackrel{~}{D}_x\eta _x|^2+|\stackrel{~}{D}_y\eta _y|^2`$
$`+\kappa _2(|\stackrel{~}{D}_y\eta _x|^2+|\stackrel{~}{D}_x\eta _y|^2)`$
$`+\kappa _5(|\stackrel{~}{D}_z\eta _x|^2+|\stackrel{~}{D}_z\eta _y|^2)`$
$`+\kappa _3[(\stackrel{~}{D}_x\eta _x)(\stackrel{~}{D}_y\eta _y)^{}+h.c.]`$
$`+\kappa _4[(\stackrel{~}{D}_y\eta _x)(\stackrel{~}{D}_x\eta _y)^{}+h.c.]+h^2/(8\pi ).`$
where $`\stackrel{~}{D}_j=_j\frac{2ie}{\mathrm{}c}A_j`$, $`𝐡=\times 𝐀`$, and $`𝐀`$ is the vector potential. The stable homogeneous solutions are easily determined. There are three phases: (a) $`\stackrel{}{\eta }=(1,i)/\sqrt{2}`$ ($`\beta _2>0`$ and $`\beta _2>\beta _3/2`$), (b) $`\stackrel{}{\eta }=(1,0)`$ ($`\beta _3>0`$ and $`\beta _2<\beta _3/2`$), and (c) $`\stackrel{}{\eta }=(1,1)/\sqrt{2}`$ ($`\beta _3<0`$ and $`\beta _2<0`$). Phase (a) is nodeless and phases (b) and (c) have line nodes. For some of the discussion in this paper the Ginzburg Landau coefficients are determined within a weak-coupling approximation in the clean limit. In this case, taking for the $`E_u`$ REP the gap function described by the pseudo-spin-pairing gap matrix: $`\widehat{\mathrm{\Delta }}=i[\eta _1v_x/\sqrt{v_x^2}+\eta _2v_y/\sqrt{v_x^2}]\sigma _z\sigma _y`$, where the brackets $``$ denote an average over the Fermi surface and $`\sigma _i`$ are the Pauli matrices, it is found that $`\beta _2=\kappa _2=\kappa _3=\kappa _4=(\nu +1)/(3\nu )`$ and $`\beta _3=4\nu /(3\nu )`$ where
$$\nu =\frac{v_x^43v_x^2v_y^2}{v_x^4+v_x^2v_y^2}.$$
(2)
Note that $`|\nu |1`$ and $`\nu =0`$ for a spherical or a cylindrical Fermi surface. Also note that $`\stackrel{}{\eta }(1,i)`$ is the stable homogeneous state for all $`\nu `$.
## 3 Magnetic field in the basal plane
For the $`E_u`$ model, symmetry arguments imply that the vortex lattice phase diagram contains at least two vortex lattice phases for magnetic fields applied along at least two of the four two-fold symmetry axes . To demonstrate this, consider the magnetic field along the $`\widehat{x}`$ direction ($`\widehat{x}`$ is chosen to be along the crystal $`𝐚`$ axis) and a homogeneous zero-field state $`\stackrel{}{\eta }(1,i)`$. The presence of a magnetic field along the $`\widehat{x}`$ direction breaks the degeneracy of the $`(\eta _x,\eta _y)`$ components, so that only one of these two components will order at the upper critical field \[e.g. $`\stackrel{}{\eta }(0,1)`$\]. It has been shown for type II superconductors with a single component order parameter that the solution is independent of $`x`$ so that $`\sigma _x`$ (a reflection about the $`\widehat{x}`$ direction) is a symmetry operation of the $`\stackrel{}{\eta }(0,1)`$ vortex phase. Now consider the zero-field phase $`\stackrel{}{\eta }(1,i)`$; $`\sigma _x`$ transforms $`(1,i)`$ to $`(1,i)e^{i\psi }(1,i)`$ where $`\psi `$ is phase factor. This implies that $`\sigma _x`$ is not a symmetry operator of the zero-field phase. It follows that there must exist a second transition in the finite field phase at which $`\eta _x`$ becomes non-zero. Similar arguments hold for the field along any of the other three two-fold symmetry directions in the basal plane. Consequently a zero-field state $`\stackrel{}{\eta }(1,i)`$ must exhibit two vortex lattice phases when the field is applied along any of the four two-fold symmetry axes. Similar arguments for $`\stackrel{}{\eta }(1,0)`$ or $`\stackrel{}{\eta }(1,1)`$. imply that there must exist at least two vortex lattice phases for only two of the four two-fold symmetry axes. For a zero-field state $`\stackrel{}{\eta }=(1,0)`$ fields along the $`(1,1)`$ or the $`(1,1)`$ directions will result in two vortex lattice phases (for the $`(1,0)`$ and the $`(0,1)`$ directions multiple vortex lattice phases may exist, but are not required by symmetry). For a zero-field state $`\stackrel{}{\eta }(1,1)`$, two vortex lattice phases will exist for fields along the $`(1,0)`$ or the $`(0,1)`$ directions.
An analysis of the free energy of Eq. 2 reveals additional information about the phase diagram . Assuming the large $`\kappa `$ limit (note $`\kappa 30`$ for this field orientation in Sr<sub>2</sub>RuO<sub>4</sub>), the vector potential can be taken to be $`𝐀=Hz(\mathrm{sin}\varphi ,\mathrm{cos}\varphi ,0)`$ ($`\varphi `$ is the angle in the basal plane that the applied magnetic field makes with the $`\widehat{x}`$ direction). The component of $`𝐃`$ along the field is set to zero. The upper critical field found by this method exhibits a four-fold anisotropy in the basal plane . For concreteness, the field is taken along the $`\widehat{x}`$ direction ($`\varphi =0`$). Introducing the raising and lowering operators $`\mathrm{\Pi }_\pm =q(\sqrt{\kappa _2/\kappa _5}D_y\pm iD_z)/\sqrt{2}`$ with $`q^2=\sqrt{\kappa _5/\kappa _2}/H`$, the gradient portion of the free energy can be written as
$`f_{grad}=`$ $`\sqrt{\kappa _5\kappa _2}H\{\eta _x^{}[1+2N]\eta _x+\eta _y^{}[({\displaystyle \frac{1}{2}}+{\displaystyle \frac{1}{2\kappa _2}})`$
$`(1+2N)+({\displaystyle \frac{1}{2\kappa _2}}{\displaystyle \frac{1}{2}})(\mathrm{\Pi }_+^2+\mathrm{\Pi }_{}^2)]\eta _y\}.`$
Assuming that $`\kappa _2<1`$ the first transition is given by the standard hexagonal Abrikosov vortex lattice solution for $`\eta _x`$. To find the second transition the lowest eigenstate of $`(\frac{1}{2}+\frac{1}{2\kappa _2})(1+2N)+(\frac{1}{2\kappa _2}\frac{1}{2})(e^{2i\theta }\mathrm{\Pi }_+^2+e^{2i\theta }\mathrm{\Pi }_{}^2)`$ must be found (where $`\theta `$ has been introduced so that the vortex lattice can be rotated with respect to the ionic lattice) and a vortex lattice solution for $`\eta _y`$ must be constructed from this. Note that the zeroes of the $`\eta _x`$ lattice and the zeroes of the $`\eta _y`$ lattice need not coincide . This method yields for the ratio of the second transition ($`H_2`$) to the upper critical field ($`H_{c2}`$)
$$\frac{H_2}{H_{c_2}^{ab}}=\frac{\beta _A\beta _m}{\beta _A\sqrt{1/\kappa _2}\beta _m}$$
(3)
where $`\beta _A=1.1596`$, $`\beta _m=(1\beta _2+\beta _3)S_1|\beta _2S_2|`$, $`S_1=\overline{|\eta _1|^2|\eta _2|^2}/(\overline{|\eta _1|^2}\overline{|\eta _2|^2})`$, $`S_2=\overline{(\eta _1\eta _2^{})^2}/(\overline{|\eta _1|^2}\overline{|\eta _2|^2})`$, the over-bar denotes a spatial average, and $`\beta _m`$ must be minimized with respect to $`\theta `$ and the displacement between the zeros of the $`\eta _x`$ and the $`\eta _y`$ lattices. Diagrams of the predicted vortex structures are shown in Fig. 1 and the phase diagram for the weak-coupling limit is shown in Fig. 2. Two vortex lattice configurations are found to be stable as a function of $`\nu `$. For $`0.23<\nu <1`$ (Phase 1), the displacement between the the zeroes of the $`\eta _x`$ and $`\eta _y`$ lattices is one half of a hexagonal lattice basis vector and $`\theta `$ varies with $`\nu `$ (and in general with field and temperature) so that the vortex lattice is not aligned with the ionic lattice. For $`1<\nu <0.23`$ (phase 2), the $`\eta _x`$ and $`\eta _y`$ lattices coincide and a vortex lattice basis vector lies along the $`\widehat{y}`$ direction. Kita has argued that for phase 1, at some field below $`H_2`$ there will be a first order transition from phase 1 to phase 2. For the field along $`\widehat{x}\pm \widehat{y}`$ the phase diagram is given by replacing $`\nu `$ with $`\nu `$. Recent experiments of Mao et al. show some support for the theory presented here.
## 4 Magnetic field along the $`c`$-axis
The free energy in Eq. 2 in the weak coupling limit has been used to determine the vortex lattice structure for the field along the $`c`$-axis for fields near the upper critical field . The main conclusion of this analysis is that vortex lattice is square for $`|\nu |>0.0114`$ and that the vortex lattice is oriented along the crystal lattice for $`\nu <0.0114`$ and rotated 45 degrees from the crystal lattice for $`\nu >0.0114`$. The analysis near $`H_{c1}`$ has also been carried out with the result that the vortex lattice at $`H_{c1}`$ will be hexagonal and with increasing field the lattice will continuously distort until a square vortex lattice is formed . A square vortex lattice has been observed by Small Angle Neutron scattering (SANS) and the orientation implies $`\nu >0.0114`$ . One notable feature of the SANS measurements is the observation of Bragg peaks beyond the first Bragg peak. The analysis of Eq. 2 has been extended to fields below $`H_{c2}`$ and the lowest Bragg peaks of the field distribution have been calculated and compared to experiment. One comparison is shown in Fig. 3 (note that in the SANS measurements there is flux pinning which is not included in the calculations). The large size of the higher order Bragg peaks relative to that expected for a single complex order parameter component (Abrikosov) theory and the agreement between the theory and the experiment gives support to the $`E_u`$ theory. However, note that it is possible that an Abrikosov theory with sufficiently large non-local corrections may also account for the observed results.
## 5 Acknowledgments
We acknowledge support from NSF DMR9527035 (DFA), the Swiss Nationalfonds, the U.K. E.P.S.R.C., and CREST of Japan Science and Technology Corporation. The neutron scattering was carried out at the Institut Laue-Langevin, Grenoble. We wish to thank T.M. Rice and M. Sigrist for informative discussions. |
warning/0001/cond-mat0001279.html | ar5iv | text | # Many–body Correlation Effects in the Ultrafast Nonlinear Optical Response of Confined Fermi Seas
## List of acronyms
| FES | Fermi Edge Singularity |
| --- | --- |
| FS | Fermi sea |
| e–h | Electron–hole |
| e–e | Electron–electron |
| h–h | Hole–hole |
| e–p | Electron–phonon |
| h–p | Hole–phonon |
| HFA | Hartree–Fock approximation |
| SBE’s | Semiconductor Bloch Equations |
| SP | Surface Plasmon |
| FWM | Four Wave Mixing |
| PP | Pump–probe |
| RPA | Random Phase Approximation |
| QW | Quantum Well |
| CCE | Coupled Cluster Expansion |
| fs | Femtosecond |
| ps | Picosecond |
| MDQW | Modulation–doped Quantum Well |
| TI–FWM | Time–integrated Four Wave Mixing |
## 1 Introduction and Scope of this Review
One of the most exciting thrusts in Physics focusses on understanding the properties of systems possessing many degrees of freedom. The goal is to relate such properties to the interactions among the elementary excitations. Such many–body problems have been a driving force in many different subfields of Physics and Chemistry: condensed matter physics, nuclear physics, quantum chemistry, high energy physics, $`\mathrm{}`$ In condensed matter physics, the study of many–body phenomena has led to many exciting discoveries: superconductivity, quantum Hall effect, Kondo effect, Mott transition, quantum phase transitions, Fermi Edge and X–ray Edge Singularities, $`\mathrm{}`$. In most cases, such fascinating phenomena have been observed in equilibrium systems at low temperatures. In comparison, the role of many–body effects in the dynamics of non–equilibrium many–electron systems presents a less explored frontier.
Ultrafast nonlinear optical spectroscopy provides unique and powerful tools for studying the dynamics of many–body correlations. In such experiments, an intense optical pulse, the pump, with duration shorter than the scattering time or the period of the elementary excitations photoexcites the system from its ground state. The subsequent dynamics is then monitored by using a second pulse, the probe. By measuring the amplitude and the phase of the optical field emerging from the photoexcited sample, one can gain valuable insight into the interplay between many–body and quantum confinement effects during very short time scales. A good understanding of the role of such fundamental physics can lead to new ideas regarding the next generation of photonic and opto–electronic devices. These operate under intense excitation conditions in the technologically important Tbit/sec regime and are often based on semiconductor heterostructures and novel low–dimensional structures. Such considerations have spurred intense efforts in recent years aimed at understanding how to control the optical dynamics during sub–picosecond time scales (1 Tbit/sec = (1 ps)<sup>-1</sup>). In the highly non–equilibrium femtosecond regime, the dominant effects on the nonlinear optical dynamics come from the Pauli exclusion principle (Phase Space Filling effects) and the Coulomb interactions among the optically–excited and Fermi sea carriers. The phonons also play an important role, especially during picosecond time scales.
Following their photoexcitation, the electrons and holes undergo a number of different relaxation stages as they scatter among themselves via the Coulomb interaction and with the phonons via the electron–phonon (e–p) interaction. Below we briefly discuss each of these relaxation stages. Initially (stage I, the coherent regime) the photoexcited electrons and holes have a well–defined phase relationship: since the photon momentum is negligible, they are restricted to have opposite momenta. Such a coherence is described by an optical polarization. In the very beginning, the photoexcited carriers interact via the bare Coulomb interaction. After time intervals of the order of the inverse plasmon frequency corresponding to the carrier density, the interactions become screened. The problem can then be approached in terms of interacting Fermi liquid quasi—particles . Even though in typical metals such screening occurs within an extremely short time interval $``$ 1 fs, in semiconductors the inverse plasma frequency is of the order of tens of femtoseconds and therefore the above regime is accessible experimentally by using $``$ 5–10 fs optical pulses. If one tunes the pump frequency above the onset of linear absorption, real carriers are photoexcited, in which case the duration of stage I is determined by the dephasing time, $`T_2`$. In metals, the latter is typically of the order of a few femtoseconds, while in semiconductors it can be as long as a few picoseconds. The dephasing processes are suppressed if one tunes the pump frequency below the onset of linear absorption. Under such below–resonance excitation conditions, the coherent effects become dominant since only virtual carriers are photoexcited and thus dissipative scattering processes cannot occur. The time–dependence of the optical polarization is then determined by the pulse duration, which leads to extremely fast nonlinearities.
Immediately after the e–h phase coherence is lost, incoherent populations of electrons and holes are formed. These are initially described by non–thermal density distributions. Subsequently, these evolve into hot thermal distributions characterized by an electron temperature that can far exceed that of the lattice. Such a time evolution (stage II) as the carriers equilibrate among themselves is mainly induced by the e–e, e–h, h–h, and e–p interactions. The duration of this stage II is determined by the population relaxation time, $`T_1`$, which can be of the order of picoseconds. Subsequently, the e–p interaction leads to the equilibration of the electronic and lattice temperatures on a time scale of tens of picoseconds, followed by the return to thermodynamic equilibrium after hundreds of picoseconds via the recombination of the photoexcited carriers or via the transport of heat to the surroundings (stage III).
A large body of experimental and theoretical work has focussed on different aspects of the above relaxation. Such developments in undoped semiconductors have been extensively described in a number of excellent review articles and books. Among the earlier reviews of interaction effects in highly–excited undoped semiconductors, we note those of Haug and Schmitt–Rink , Stahl and Baslev , Zimmermann , Schmitt–Rink et. al. , and Haug and Koch . More recent reviews focussing on the femtosecond coherent dynamics in undoped semiconductors include those by Mukamel , Binder and Koch , Shah , Haug and Jauho , Axt and Mukamel , and Chemla . We also note the collections of papers in Refs. . More recently, there has been a lot of interest in using ultrafast nonlinear spectroscopy to study the dynamics of many–body effects in metals. The charge carrier dynamics and the coherent ultrafast phenomena in bulk metals, metal interfaces, and metal surfaces were recently reviewed by Petek and Ogawa .
The purpose of this review article is to provide a comprehensive picture of recent experimental and theoretical advances in understanding femtosecond many–body phenomena in Fermi sea systems. As has been extensively discussed in the review articles and books listed above, even though a description of the ultrafast dynamics based on few–level models has been very successful in atomic and molecular systems, it faces serious shortcomings when applied to semiconductors and metals. The development of microscopic non–equilibrium many–body theories has been found to be necessary for describing the femtosecond dynamics and has made this field an exciting area of research in condensed matter physics. In view of the large body of work, especially in undoped semiconductors, we will concentrate here on specific aspects of the many–body physics that have not been extensively reviewed elsewhere. In particular, we will not touch on the impressive advances in undoped semiconductors and will attempt to provide a unified picture of the ultrafast dynamics in different confined Fermi sea systems. We will mainly focus on coherent effects due to the Coulomb interaction and will only briefly discuss the electron–phonon (e–p) interaction effects. The main body of this article discusses the recent breakthroughs in understanding the sub–picosecond dynamics in modulation–doped quantum wells (MDQW’s) and metal nanoparticles, with particular emphasis on non–Markovian memory effects induced by carrier–carrier many–body correlations. Such dynamical effects are beyond the scope of the dephasing and relaxation time approximations. One should note here that a detailed quantitative description of the experimental results presents a formidable task in the subpicosecond regime, which requires sophisticated numerical simulations as well as different approximations. This can make it difficult to discuss in a simple and intuitive way the important many–body physics conveyed by the experiment. Here we will focus on the simplest possible models that can capture the most important physics associated with the non–equilibrium Coulomb correlations and attempt to provide, to the degree possible, an intuitive explanation of the often delicate experimental features.
Let us now briefly discuss how many–body correlations enter into the ultrafast nonlinear optical response. The most commonly used theoretical approach starts with the Heisenberg equations of motion of the one–particle density matrix, $`\rho _{ij}(t)`$, where the indices $`i,j=1,2`$ label the conduction and valence band . The off–diagonal density matrix element, $`\rho _{12}(t)=P(t)`$, is the optical polarization, while the diagonal matrix elements, $`\rho _{ii}(t)`$, correspond to the electron and hole occupation numbers (density distributions). Such equations of motion must be solved to obtain the polarization, which in turn acts as the source term in the Maxwell equations that determine the signal measured in the experiment . The difficulty in treating the many–body interactions comes from the coupling of the one–particle density matrix to many–particle correlation functions (higher density matrices). As is typical for many–body systems, this leads to an infinite hierarchy of coupled equations which cannot be solved exactly . One must therefore devise a controlled truncation scheme. One possibility is to adopt the time–dependent Hartree–Fock approximation (HFA), which relies on the factorization of the many–particle correlation functions and only includes the two–particle interaction effects. Such an approach developed into the very successful Semiconductor Bloch Equations (SBE’s) formalism (see e.g. Refs. ). However, many recent experimental observations, reviewed e.g. in Ref. , could not be explained within such a mean field approach and were attributed to four–particle and higher order Coulomb correlations. In order to develop a many–body theory that can capture such effects, it was noted by several different groups that, in undoped semiconductors, the above infinite density matrix hierarchy truncates if one adopts an expansion in terms of the optical fields . This is due to the fact that, in the undoped semiconductor ground state, the conduction band is empty while the valence band is full. Such a method for treating the Coulomb correlations nonperturbatively is often referred to as the “dynamic controlled truncation scheme”. An alternative many–body approach has been based on diagrammatic expansions for Keldysh Green’s functions, whose multiple–time dependence is typically eliminated by using the Kadanoff–Baym ansatz .
Let us now discuss how non–Markovian memory effects enter into the density matrix equations of motion. Due to the many degrees of freedom that affect the dynamics, one usually distinguishes between a subsystem interacting directly with the optical fields (e.g. the photoexcited e–h pair) and a reservoir/bath consisting of all the other degrees of freedom (e.g. the phonons or the Fermi sea excitations). The polarization equation of motion can be cast in the form
$$\frac{}{t}P(t)=\frac{}{t}P(t)|_{\mathrm{coh}}+\frac{}{t}P(t)|_{\mathrm{scatt}}$$
(1)
where the first term describes the Hartree–Fock interactions among the coherent e–h pairs, while the second term describes the many–body scattering processes among the carriers as well as between the carriers and the bath excitations. The latter term can be expressed in the general form
$$\frac{}{t}P(t)|_{\mathrm{scatt}}=_{\mathrm{}}^t𝑑t^{}\mathrm{\Gamma }(t,t^{})P(t^{})$$
(2)
where $`\mathrm{\Gamma }(t,t^{})`$ is a memory kernel (similar equations can also be written for the carrier distribution functions). The importance of the many–body correlations of interest here is that they determine the dependence of $`\mathrm{\Gamma }(t,t^{})`$ on the two times $`t`$ and $`t^{}`$ and therefore govern the memory effects. In systems where only times $`t^{}t`$ contribute to the above integral, one can apply the Markovian approximation, in which case
$$\frac{}{t}P(t)|_{\mathrm{scatt}}\mathrm{\Gamma }(t)P(t).$$
(3)
In many cases, one can further approximate $`\mathrm{\Gamma }(t)`$ by a time–independent constant, the dephasing energy width $`1/T_2`$. For example, such an approximation allows one to include, in a semi–phenomenological way, the effect of the scattering processes in the SBE’s. However, when such a Markovian approximation breaks down, the many–body correlations can affect the qualitative ultrafast dynamics by inducing strong memory effects and non–exponential polarization decay. Such effects in undoped semiconductors are known to arise e.g. from exciton–exciton correlations or from e–p and carrier–carrier interactions in the quantum kinetic regime and have been extensively reviewed e.g. in Refs. . Here we focuss on analogous effects whose origin lies in the non–perturbative Coulomb correlations between the photoexcited carriers and the cold electron Fermi sea in MDQW’s and small metal nanoparticles.
This review is organized as follows. First we present an overview of the experimental results on the ultrafast nonlinear optical dynamics in MDQW’s (Section 2) and metal nanoparticles (Section 3). In Section 4 we outline the formalism necessary for describing the coherent ultrafast nonlinear optical response of Fermi sea systems. In Section 5, the ultrafast nonlinear dynamics of the Fermi Edge Singularity (FES) in the coherent regime is discussed. In Section 6 the size-dependent dynamical screening of the Coulomb interaction in metal nanoparticles is discussed and the results for the quasiparticle scattering rates are reviewed. In Section 7 these results are used to describe the ultrafast surface plasmon (SP) dynamics in small noble-metal nanoparticles. Section 8 concludes the review.
## 2 Nonlinear Optical Dynamics in Modulation–doped Quantum Wells
By separating the ionized impurity donors from the electrons, modulation doping can produce a high mobility two–dimensional electron gas. In contrast to the case of a photoexcited electron gas, the temperature of such a system can be lowered, which enables the observation of fascinating many–body phenomena such as the Quantum Hall Effect and the Fermi Edge Singularity .
Recently there has been a lot of interest in studying the ultrafast dynamics in MDQW’s. The role of the interactions between the photoexcited and the Fermi sea (FS) carriers during the thermalization stage II has been reviewed e.g. in Refs. and is only briefly discussed here. In the section 3.2 we will review corresponding measurements in metal films and nanoparticles and compare with MDQW’s. Much less is known about the role of many–body effects in the coherent regime (stage I). Even though some aspects of the effects of the e–e scattering have been described within the dephasing time approximation, the role of the e–h correlations between the photoexcited holes and the FS electrons has only recently been studied. The main difficulty in treating such e–h correlation effects comes from the non–perturbative nature of the Fermi edge singularity (FES), which dominates the absorption spectrum close to the absorption onset, and which cannot be described within the dephasing time approximation. Important here is that the e–h interactions between a heavy photoexcited hole and the FS-electrons lead to the scattering of a macroscopic number of low–energy FS-pairs, which readjusts the entire FS into a new orthogonal scattering state in the course of the optical excitation (Anderson orthogonality catastrophe ). Before we discuss the ultrafast dynamics due to such effects, we briefly summarize in the next section the linear optical properties of MDQW’s.
### 2.1 Fermi Edge Singularity in Linear Absorption
Close to the onset of linear absorption and for low–temperatures, the optical properties of MDQW’s are dominated by the FES. The latter is a many–body resonance that has been observed in doped semiconductors (see e.g ) as well as in metals , where it is referred to as the X–ray Edge Singularity. Despite the screening of the Coulomb interaction, the strength of this resonance is comparable to that of the undoped QW excitons. The non–Lorentzian lineshape of the FES can be viewed as originating from the decay of an excitonic bound state caused by its interactions with the gapless FS excitations . In the case of holes localized due to e.g. the disorder or for core holes in metals, the lineshape of the FES close to the onset of absorption can be approximately described by using the analytic power law expression
$$\alpha (\omega )𝒩\left(\frac{E_F}{\omega }\right)^\beta ,$$
(4)
where $`𝒩`$ is the density of states, $`E_F`$ is the Fermi energy, $`\omega `$ is the optical frequency measured from the Fermi edge, and $`\beta =2\delta /\pi \left(\delta /\pi \right)^2`$ is the FES exponent, where $`\delta \mathrm{tan}^1(\pi g)`$ is the s–wave phaseshift of the screened e–h potential $`V`$ evaluated at $`E_F`$, and $`g=V𝒩`$ is the dimensionless parameter characterizing the e–h scattering strength. Such a non–Lorentzian lineshape results from the competition between the Mahan singularity, due to the attractive interaction between the FS and the localized hole (vertex correction effect, ) and the Anderson orthogonality catastrophe due to the readjustment of the FS density profile during the optical transition (hole self–energy effect, ). In the case of finite hole mass $`m_h`$, the FES is broadened by an additional energy width of the order of the hole recoil energy $`m_e/m_hE_F`$, where $`m_e`$ is the electron mass .
In a first approximation, the excitonic effects in a MDQW can be described by extending the mean field approach (HFA) to include the effects of the screening and the Pauli blocking due to the FS . However, such a static treatment of the FS leads to a spurious bound state with respect to the Fermi energy, $`E_F`$, referred to in the following as the HFA bound state . This discrete state, with binding energy $`E_M`$, would appear at the energy $`E_FE_M`$. Obviously, for $`E_M<E_F`$ (as in typical MDQW’s), such a state cannot exist since it overlaps with the FS continuum, with which it interacts via the e–h potential . The “unbinding” of this HFA bound state occurs via its interactions with the FS excitations, which are not taken into account in the HFA . Note that, in two–dimensional systems, a static FS allows for bound states even for arbitrarily small attractive interactions and therefore such an unbinding cannot arise from static screening or from Pauli blocking effects. This spurious bound state could be artificially merged with the continuum by introducing a dephasing time comparable to its binding energy. Such an approximation, however, neglects competely the dynamical correlations between the photoexcited e–h pair and the FS excitations. A microscopic description of the unbinding of the HFA bound state presents a nontrivial problem due to the non–perturbative nature of the e–h correlations between the photoexcited hole and the FS excitations . Even in linear absorption, within Green function techniques, this problem is rather involved because vertex correction diagrams with arbitrarily many crossed e–h interaction lines are as divergent as the ladder diagrams and should be treated on equal footing . To perform such a task, one must sum up at least the parquet diagrams and address the three–body correlations between the photoexcited hole and a FS excitation . Therefore, alternative methods were developed for the case of linear absorption, based on Fermi’s golden rule with many–electron eigenstates expressed in terms of Slater determinants . Such approaches become exact in the limit $`m_e/m_h1`$ and describe quite accurately the FES lineshapes observed in typical MDQW’s . Another approach to the FES problem is based on the coupled cluster expansion (CCE) . This general many–body technique provided exact results in the limit $`m_e/m_h1`$ and was used to treat the hole recoil correlations in one dimension . In the latter case, an exact solution was obtained for $`m_e=m_h`$. The CCE has also been used to describe the e–e correlation effects . More importantly, however, this method is well–suited for describing correlations in non–equilibrium systems, where it retains the advantages of diagrammatic expansions without resorting to the Kadanoff-Baym ansatz or to the Markovian approximation .
### 2.2 Ultrafast Dynamics
In this section we summarize some recent studies of the dephasing and relaxation processes in MDQW’s. The earlier work focussed on the role of the interactions between the FS and the photoexcited carriers on the thermalization during relaxation stage II. In most cases, a non–thermal distribution of real electron and hole carriers was photoexcited well above the Fermi surface and the subsequent incoherent population dynamics was monitored with ultrafast pump–probe spectroscopy. Note that, in such experiments, the measured differential transmission is
$$DST(\omega ,\tau )=\frac{\mathrm{\Delta }T_s(\omega ,\tau )}{T_0(\omega )}=\frac{T_s(_p)T_s(_p=0)}{T_s(_p=0)},$$
(5)
where $`T_s(_p)`$ is the transmission coefficient in the probe direction in the presence of the pump field $`_p`$ and $`\tau `$ denotes the time delay between the pump and probe optical pulses. For samples with sufficiently small thickness $`d`$ such that $`\mathrm{\Delta }\alpha d`$ is small, the differential transmission reproduces the change, $`\mathrm{\Delta }\alpha (\omega ,\tau )`$, in the probe absorption coefficient $`\alpha (\omega ,\tau )`$ which is induced by the pump photoexcitation: $`DST(\omega ,\tau )\mathrm{\Delta }\alpha (\omega ,\tau )d`$.
For a simple intuitive interpretation of the thermalization experiments in MDQW’s, it is often assumed that, for free carriers photoexcited well above the onset of absorption and for time scales much longer than the dephasing time, the differential absorption spectrum maps the carrier distribution functions at the probe photon energy at time $`\tau `$. In particular, the following approximation is often used:
$$\mathrm{\Delta }\alpha (\omega ,\tau )=(1f_ef_h)\alpha (\omega )$$
(6)
where $`\alpha (\omega )`$ is the linear absorption coefficient and $`f_e`$ and $`f_h`$ are the electron and hole distribution functions at the corresponding energies. Within such an approximation, the time evolution of the differential absorption is determined by the carrier distribution functions at time $`\tau `$. There have been many theoretical calculations of the distribution function time evolution, mostly based on numerical solutions of the semiclassical Boltzmann equations (see e.g. Refs. ). A detailed review of such an approach to thermalization can be found in Ref. . More recently, the dynamics in the quantum kinetic regime and the shortcomings of the Boltzmann equations have been addressed, as reviewed e.g. in Ref. .
Knox et. al. studied the effect of the inelastic e–e scattering on the thermalization by photoexciting a non–thermal carrier distribution well above the Fermi surface and then monitoring its time evolution to a thermal distribution at the bottom of the band. They observed that the e–e interactions between the photoexcited and FS electrons significantly enhance the thermalization rates as compared to undoped semiconductors. In particular, the thermalization time in MDQW’s was found to be of the order of $``$10fs. The differential transmission lineshape mimicked that of a Boltzmann distribution peaked at the bottom of the conduction band, which indicates a non–degenerate carrier distribution. This is expected at room temperatures (as in the experiment of Ref. ), where the thermal energy exceeds the Fermi energy, and at high intensities, where the photoexcited carrier density exceeds the FS density. Knox et. al. also observed an instantaneous redshift due to the bandgap renormalization by the photoexcited carriers . However, as we discuss later, the situation changes drastically at low temperatures, where the FS electrons have a sharp Fermi–Dirac distribution. Wang et al. performed low temperature pump–probe measurements and found that, at low excitation densities, the differential transmission lineshape corresponded to a redshift of the FES due to the incoherent bandgap renormalization. Importantly, they showed that such a redshift slowly builds up with time as the hot electron temperature decreases due to the cooling to the lattice temperature. This observation points out that the bandgap renormalization depends not only on the carrier density but also on the carrier distribution. At high photoexcited densities, the heating of the electron gas becomes stronger and Wang et. al. observed a bleaching of the FES due to the smearing of the sharp Fermi surface.
Both the above experiments were interpreted in terms of e–e interactions between the photoexcited and the FS carriers. Tomita et. al. performed ultrafast luminescence measurements at low temperatures to investigate the role of the e–h interactions on the thermalization in n–type MDQW’s. In particular, they measured the time evolution of the luminescence intensity at the bandgap energy, where the optically–induced perturbation of the electron distribution due to, e.g., heating is minimal. Tomita et. al. found that the number of holes that thermalized to the top of the valence band as a function of time could not accounted for if one assumes a thermalized hole distribution. This result indicates that a substantial fraction of the photoexcited holes are non–thermal for time intervals as long as $``$ 800fs. Furthermore, Monte–Carlo simulations indicated that the e–h scattering is the dominant hole thermalization mechanism. Similar measurements in n–doped bulk GaAs were performed by Chebira et. al. (at low temperatures) and Zhou et. al. (at room temperature) . Woerner et. al. studied the hole relaxation in p–doped semiconductors by photoexciting holes from the heavy–hole to the split–off valence band and then monitoring the time evolution of the inter–valence–band absorption spectrum by using femtosecond pulses in the mid–infrared region. Due to the absence of conduction electrons, the dynamics is then solely determined by the h–h and h–p scattering. They found that the inter–valence–band scattering via the emission of e.g. optical phonons occurs within a short time interval $`<`$100fs, and that subsequently the phototexcited holes slowly thermalize via h–h scattering with the hole FS within a time interval $``$ 700fs.
An important effect of the quantum confinement in QW’s is the formation of discrete conduction and valence subbands. Recently, there has been much interest in exploring the role of the intersubband excitations in the ultrafast dynamics, partly motivated by device applications such as the quantum cascade laser . Lutgen et. al. performed femtosecond pump–probe measurements and showed that the intersubband nonlinear absorption is determined by both the intersubband and the intraband electron relaxation. They measured the effects of intersubband pump excitation on the nonlinear absorption due to probe–induced interband transitions from the valence to the conduction band. They observed a population relaxation that proceeds in two stages. First, the electrons photoexcited by the pump from the Fermi sea in the lowest conduction subband relax back via intersubband scattering within a time interval $``$ 1.3 ps. The latter is consistent with the intersubband relaxation times due to longitudinal optical phonon emission . This process is followed by intraband e–e relaxation similar to the experiments discussed above. Lutgen et. al. also measured the time evolution of the nonlinear absorption spectrum due to optical transitions between the conduction subbands. They observed an intial decay of the differential transmission within a time interval $``$ 2–3 ps, which they attributed to the population relaxation due to the intersubband scattering. This was followed by a slowly decaying differential transmission signal, whose time evolution depended strongly on the photoexcitation frequency (within the linear absorption linewidth). The latter regime was attributed to the intraband population relaxation within the lowest subband as the carriers cool to the lattice temperature.
We now come to the effect of a cold FS on the e–e scattering time. One would expect that the latter should become short as the FS density increases. This is indeed the case for low FS densities, high temperatures, or for large excess energies from the Fermi surface. Hawrylak et. al. obtained the energy–dependent e–e scattering times by calculating the equilibrium self–energies including the effects of the electron–plasmon and e–p interactions. They predicted step–like decreases in the scattering time for electron energy above the Fermi surface that exceeds the onset of plasmon and optical phonon emission. They also concluded that short–range e–e correlations (not included in the RPA) can significantly affect the scattering times at high energies. Such fast e–e scattering of photoexcited electrons well above the Fermi surface can explain for example the absence of spectral hole burning in the pump–probe spectra of MDQW’s . However, the situation changes drastically for electron energies close to the Fermi surface and temperatures smaller than the Fermi energy. Under such conditions, the e–e scattering is in fact suppressed in the presence of a degenerate FS. Indeed, a photoexcited electron lowers its energy by interacting with a FS electron while the latter scatters above the Fermi surface. For electrons close to the Fermi surface, the drastic reduction of the phase space available for scattering due to the Pauli blocking by the FS electrons as well as the screening leads to the strong suppression of the e–e scattering . This is due to the sharp Fermi–Dirac distribution of the FS electrons at low temperatures. In fact, at zero temperatures, the e–e scattering time $`\tau _e(E)`$ becomes infinite right at the Fermi surface, $`E=E_F`$ . According to Fermi liquid theory ,
$$\tau _e(E)\frac{E_F^2}{(EE_F)^2+(\pi k_BT)^2},$$
(7)
where $`k_BT`$ is the thermal energy. Note that the above result applies both to metals and MDQW’s due to the similar values of $`r_s`$. As the electron temperature increases, the smearing of the Fermi–Dirac distribution allows for the scattering of a photoexcited electron to states below $`E_F`$ and therefore $`\tau _e(E)`$ decreases as compared to the zero temperature limit. Even though the e–e scattering rate initially increases with electron density for low densities, it decreases at higher densities when the FS temperature is low and the Fermi edge is sharp . Therefore, one should expect delayed dephasing and thermalization processes as the electrons approach the Fermi surface (note that the above phase space restriction does no apply to the e–h scattering of the photoexcited hole with the FS electrons).
The effect of a cold FS on the dephasing time of the interband optical polarization was investigated by Kim et al. using transient FWM spectroscopy . In such experiments, two optical pulses separated by a time delay $`\tau `$ propagate along two different directions, $`𝐤_\mathrm{𝟏}`$ and $`𝐤_\mathrm{𝟐}`$, and interfere within the sample, thus generating a nonlinear polarization. In FWM spectroscopy, one measures the signal that emerges along the direction $`2𝐤_\mathrm{𝟐}𝐤_\mathrm{𝟏}`$. In the case of optically–thin samples, the time–resolved FWM signal corresponds to the square of the amplitude of the nonlinear polarization along the above direction . The time–integrated FWM signal (TI–FWM) is given by the integral over all times of the time–resolved FWM signal and is a function of the time delay $`\tau `$. To a first approximation, the FWM signal can be interpreted similar to atomic and molecular systems based on a two–level system . Within such a model, in the case of homogeneous broadening, the time–resolved FWM signal is emitted immediately after the second pulse arrives and decays with a decay time $`T_2/2`$, where $`T_2`$ is the dephasing time. The TI–FWM signal also decays with a decay time $`T_2/2`$ and vanishes for negative time delays where pulse 1 comes after pulse 2. In the case of inhomogeneous broadening, the time–resolved FWM signal is delayed by a time interval $`\tau `$ after the second pulse arrives and corresponds to a photon echo . The TI–FWM signal decays with a decay time $`T_2/4`$, which is a factor of 2 smaller than in the case of homogeneous broadening, and also vanishes for negative time delays. In undoped semiconductors, significant deviations from the above simple picture have been observed and have been attributed to Hartree–Fock exciton–exciton interactions and to exciton–exciton correlations . For example, such effects lead to a strong TI–FWM signal for negative time delays as well as a delayed time–resolved FWM signal. For a detailed review of the exciton–exciton interaction effects in undoped semiconductors, see e.g. Ref. .
In MDQW’s, Kim et. al. measured long dephasing times of a few picoseconds for energies within the frequency range of the FES (i.e. close to $`E_F`$). These decreased to less than 100fs as the excitation frequency exceeded the FES peak by an energy comparable to the Fermi energy. The above result is consistent with the Fermi liquid energy dependence of the e–e scattering time, Eq. 7, and points out that, unlike for the e–h correlations, in MDQW’s the e–e scattering is suppressed within the frequency range of the FES. This experiment also suggests that the hole dephasing times are of the order of a few picoseconds. In the time–resolved FWM signal, Kim et. al. observed photon–echo–like behavior, which indicates that, despite the strong FES peak in the linear absorption, the continuum of interband e–h states leads to behavior similar to that of an inhomogeneously–broadened system. Finally, in MDQW’s, the pair–pair interactions, analogous to the exciton–exciton interactions in undoped semiconductors, are screened out and the negative time delay TI–FWM signal characteristic of exciton–exciton interactions was absent for high Fermi sea densities . The above experimental results were interpreted by Hawrylak et. al. within the two–level system approximation (which neglects all FES correlation effects) by introducing energy–dependent dephasing times determined by the e–e and e–p scattering times.
In order to provide further insight into the dephasing processes in MDQW’s, Bar–Ad et. al. performed low temperature FWM measurements in the presence of a magnetic field parallel to the QW growth axis. For zero magnetic field, they observed a dephasing time $``$ 1ps. For magnetic fields such that only the lowest Landau level is occupied, they observed a suppression of the e–e scattering. In particular, the dephasing times were then found to be longer, $``$1.5 ps, and, unlike in the zero field case of Ref. , they did not depend significantly on the frequency of the optical excitation. Bar Ad et. al. attributed the above observations to the suppression of the phase space available for scattering induced by the magnetic confinement. Indeed, the application of a magnetic field changes the continuous density of states into a series of discrete highly degenerate Landau level peaks, broadened by the disorder. Finally, Bar–Ad et. al. observed a FES in the FWM spectrum at high magnetic fields, which they attributed to the optical transitions to the occupied lowest Landau level.
The dephasing of the intersubband optical polarization in MDQW’s was recently investigated by Kaindl et. al. using transient FWM in the midinfrared. They observed an exponential decay of the TI–FWM signal that was fairly insensitive to the temperature. By modelling their data with a HFA calculation, Kaindl at. al. deduced the values of the dephasing times to be of the order of hundreds of femtoseconds and concluded that the main dephasing mechanism is the e–e scattering. By increasing the Fermi sea density, they observed faster dephasing. In another experiment, Bonvalet et. al. used an extremely short, $``$ 12fs, pulse to excite a wavepacket consisting of states in the two lowest conduction subbands of a QW. By measuring the radiated coherent electromagnetic field, they observed a quantum beat at room temperature oscillating at the inverse intersubband excitation energy. From the decay of the signal they deduced an intersubband polarization dephasing time of 110fs in the case of a MDQW, which was shorter than the dephasing time of 180fs observed in an undoped QW. Such polarization dephasing is determined by the destructive interference effects as well as by the different scattering processes.
The above measurements of the polarization dephasing were performed under resonant photoexcitation conditions. In contrast to undoped semiconductors, the relaxation of such real pump–induced e–h pairs due to their e–e interactions with the FS will obscure any coherent effects. Fortunately, the dissipative processes can be suppressed by tuning the pump frequency below the onset of absorption. In this case, only virtual carriers are excited by the pump field and thus the coherent effects will dominate. For well–below–resonance pump excitation, the e–h pair phase is primarily determined by the optical field and the time dependence of the interband polarization follows that of the pump pulse (adiabatic following ). In this case, the pump creates a truly coherent e–h pair many–body state that lasts for the duration of the pulse. Such pump–probe experiments in atomic systems showed a resonance blueshift ( optical Stark effect) accompanied by a resonance bleaching due to the Phase Space Filling by the pump–induced e–h pairs . In the case of excitons, it was shown that the Coulomb interaction can significantly alter the bleaching observed in atomic systems. In particular, at low pump intensities or for pulse duration longer than the dephasing time, the interactions lead to a pure exciton blueshift without significant bleaching. For high intensities and pulse duration shorter than the dephasing time, the exciton bleaching is very strong. The exciton ac–Stark effect has been reviewed e.g. in Refs. . Such interaction effects in undoped semiconductors suggest that the many–body e–h correlations between the photoexcited holes and the FS electrons will significantly affect the coherent nonlinear optical response of the FES. Brener et. al. studied this issue by performing pump–probe measurements for pump detunings $``$ 50meV below the Fermi surface, much larger than the e–h Coulomb energy $`E_M`$ (which is much smaller than the Fermi energy $``$15meV). They observed a qualitatively different bleaching of the FES as compared to excitons in undoped QW’s, which they attributed to the distinct nature of the two resonances. Schäfer studied this issue by performing quantum kinetic calculations using Keldysh Green’s functions. Such an approach treats the HFA coherent effects and includes the scattering contributions within the second order Born, Markovian, and static screening approximations. The numerical results for Fermi energy equal to the exciton binding energy showed a strikingly different bleaching between the FES and the exciton, which was attributed to the e–e scattering.
As we discuss at length in section 5, the e–h correlations between the photoexcited hole and the FS electrons lead to strong dephasing and affect qualitatively the coherent ultrafast dynamics of the FES. An important difference from undoped semiconductors is that, due to its gapless excitation spectrum, a FS responds unadiabatically to time–dependent perturbations. In contrast, because of its finite Coulomb binding energy, an exciton can be polarized by the pump optical field without being ionized. A non–equilibrium treatment beyond the HFA is necessary in order to take into account the unadiabatic time–dependent change in the e–h pair–FS interactions and the e–h scattering processes induced by the ultrafast pump excitation in the coherent regime. As has already been noted in the case of photoexcitation of an electron gas within the continuum of states of an undoped semiconductor , the loss of coherence due to the many–body correlations cannot be fully described within the dephasing time approximation. In section 5 we show that memory effects due to the e–h correlations lead to a time evolution of the pump–probe signal characterized by the inverse Coulomb energy $`E_M`$ rather than by the dephasing time (as in the case of exciton bound states).
Concluding this section, let us briefly discuss a very recent FWM experiment by Fromer et. al. , which demonstrated for the first time non–Markovian memory effects in the Quantum Hall Effect regime. Even though for magnetic fields between 5.5 and 6.5 T the TI–FWM profile was found to be a single exponential with an unusually long decay time, for magnetic fields that exceed 7.5 T such a profile was more complicated and characterized by a change of slope indicating memory effects in the polarization dynamics. Such effects were also seen in the frequency domain: the FWM spectrum profile changed from a Lorentzian lineshape to an asymmetric one corresponding to a frequency dependent width, $`\mathrm{\Gamma }(\omega )`$. To interpret such behavior, we note that, at high magnetic fields such that the cyclotron energy, $`\mathrm{}\omega _c`$, is large compared to other characteristic energies of the system, the relaxation is dominated by intra-Landau-level processes. Such scattering by collective excitations involves the matrix elements of the dynamically screened interaction, $`U_{ij}^<(t,t^{})`$, which in the lowest Landau level have the form:
$`U_{ij}^<(t,t^{})`$ $`=`$ $`{\displaystyle \frac{d𝐪}{(2\pi )^2}e^{q^2l^2/2}v_q^2\overline{\chi }_q^<(t,t^{})c_{ij}(q)},`$ (8)
where $`\overline{\chi }_q^<(t,t^{})=\overline{\rho }_𝐪(t^{})\overline{\rho }_𝐪(t)`$ is the density–density correlation function projected onto the lowest Landau level , and $`\overline{\rho }_𝐪(t)`$ is the corresponding density operator. Here, $`v_q`$ is the unscreened Coulomb interaction, $`l=(\mathrm{}/eB)^{1/2}`$ is the magnetic length, and the coefficients $`c_{ij}(q)`$ with $`i,je,h`$ model the asymmetry in the e-e and e-h interaction matrix elements, which originates from the difference between the electron and hole lowest Landau level wavefunctions. Because of the breakdown of perturbation theory due to the Landau level degeneracy in 2D systems, it is incorrect to evaluate $`\chi _q^<(t,t^{})`$ within the standard RPA . Instead, one should account for the true excitations of the interacting two–dimensional electron liquid. Several models can be found in the literature, and we base our discussion on the magnetoroton model, which is the one best suited for the filling factors $`\nu `$ in the experiment of Ref. . The most salient features are, however, general and model independent. The magnetoroton dephasing mechanism is somewhat similar to that of acoustic phonon scattering. Under the experimental conditions of Ref. , to a very good approximation, the intra-Landau-level collective excitations are not affected by the small density of photogenerated carriers, so one can use the equilibrium density correlation function. The equations for the density matrix elements then read
$`{\displaystyle \frac{\rho _{ij}}{t}}|_{scatt}=i{\displaystyle \underset{k}{}}{\displaystyle _{\mathrm{}}^t}dt^{}G_i^r(tt^{})G_j^a(t^{}t)`$
$`\times ([U_{ik}^<(tt^{})U_{kj}^<(tt^{})]\rho _{ik}^<(t^{})\rho _{kj}^>(t^{})(<>)),`$ (9)
where $`G_i^{r/a}(t)`$ is the retarded/advanced Green function, $`\rho _{ij}^<=\rho _{ij}`$, and $`\rho _{ij}^>=\delta _{ij}\rho _{ij}`$. If all $`U_{ij}`$ are equal, i.e., $`c_{ij}(q)=1`$, then the polarization scattering term vanishes. This corresponds to identical electron and hole wavefunctions in the lowest Landau level. In practice, there is always an asymmetry between electrons and holes, due to, e.g., differing band offsets, lateral confinement, and disorder. Using the results of Ref. , Eq. (8) takes the form
$`U^<(t)={\displaystyle \frac{in}{2\pi }}{\displaystyle }`$ $`{\displaystyle \frac{d𝐪}{(2\pi )^2}}e^{q^2l^2/2}v_q^2c_{ij}(q)`$ (10)
$`\times \overline{s}_q[(N_q+1)e^{i\omega _qt}+N_qe^{i\omega _qt}],`$
where $`N_q`$ is the Bose distribution function for magnetorotons of energy $`\omega _q`$, and $`\overline{s}_q`$ is the static stucture factor of the 2D electron liquid in the lowest Landau level. By comparing Eqs. (2.2) and (10), we see that the $`\omega `$ dependence of $`\mathrm{\Gamma }(\omega )`$ is determined by the Fourier transform of $`U^<(t)`$, which is governed by the $`q`$ dependence of $`\overline{s}_q`$. In the lowest Landau level, we have $`\overline{s}_q=(2\nu ^11)\stackrel{~}{s}_q`$, where $`\stackrel{~}{s}_q(ql)^4`$ for $`ql1`$, $`\mathrm{exp}(q^2l^2/2)`$ for $`ql1`$, and $`\stackrel{~}{s}_q`$ displays a peak for $`ql1`$ that leads to the magnetoroton excitations. The corresponding resonance in $`\mathrm{\Gamma }(\omega )`$ near the magnetoroton energy leads to non–Markovian behavior with a characteristic response time of approximately the inverse of this energy. The latter is estimated from the gap $`\mathrm{\Delta }`$ at the magnetoroton dispersion minimum, $`\mathrm{\Delta }0.1(e^2/ϵl)`$ for the range of $`\nu `$ in this experiment , which for $`B=10`$ T is $`1.5`$ meV.
This concludes our overview of the main ultrafast dynamical features observed in MDQW’s. In the next section we discuss the ultrafast dynamics in the case of metal nanoparticles and compare to MDQW’s.
## 3 Nonlinear Optical Dynamics in Metal Nanoparticles
The properties of small metal particles in the intermediate regime between bulk–like and molecular behavior have been the subject of great interest lately. This was motivated in part by the need to understand how the properties of matter evolve at the transition from atoms and molecules to bulk solids. An additional motivation comes from the technological trend towards electronic and optoelectronic devices based on smaller and smaller solid state structures . Metal clusters are also being used in a variety of applications, ranging from catalysis to biological and medical applications ( see e.g. Refs. ). The electronic and thermodynamic properties of metal clusters have been extensively reviewed, e.g. in Refs. . The linear optical properties of metal nanoparticles have also been reviewed in detail by Kreibig and Vollmer .
It has been known for a long time that surface collective excitations play an important role in the absorption of light by metal nanoparticles. In large particles with sizes comparable to the wave–length of light $`\lambda `$ (but smaller than the bulk mean free path), the lineshape of the surface plasmon (SP) resonance is determined by the electromagnetic effects . However, as the size of the nanoparticle becomes smaller than the mean free path of electrons in the bulk metal, quantum confinement becomes important. In small nanoparticles with radii $`R\lambda `$, the absorption spectrum is governed by quantum confinement effects. For example, the momentum non–conservation due to the confining potential leads to the Landau damping of the SP and to a resonance linewidth inversely proportional to the nanoparticle size . A review of the extensive theoretical and experimental studies of this effect may be found in Ref. . More recently, it was established that the static nonlinear optical properties of small nanoparticles are also affected by the confinement. In particular, a size–dependent enhancement of the third–order nonlinear optical susceptibilities for monochromatic photoexcitation, caused by the elastic surface scattering of single–particle excitations, was predicted by Flytzanis and collaborators and observed experimentally by Yang et. al. . Dielectric confinement also enhances the optical nonlinearities close to the SP frequency . Before we proceed with the discussion of the ultrafast dynamics of such confined Fermi seas, let us briefly summarize in the next section the main features of the SP resonance in the linear absorption spectrum.
### 3.1 Surface Plasmon Resonance in Linear absorption
In this section we summarize the basic facts regarding the linear absorption by small metal particles embedded in a medium with dielectric constant $`ϵ_m`$. We will focus primarily on noble metal particles containing several hundreds of atoms; in this case, the confinement affects the extended electronic states even though the bulk lattice structure has been established. When the particle radii are small, $`R\lambda `$, so that only dipole surface modes can be optically excited and non–local effects can be neglected, the optical properties of this system are determined by the dielectric function
$$ϵ_{\mathrm{col}}(\omega )=ϵ_m+3pϵ_m\frac{ϵ(\omega )ϵ_m}{ϵ(\omega )+2ϵ_m},$$
(11)
where $`ϵ(\omega )=ϵ^{}(\omega )+iϵ^{\prime \prime }(\omega )`$ is the dielectric function of a metal particle and $`p1`$ is the volume fraction occupied by nanoparticles in the colloid. Since the $`d`$–electrons play an important role in the optical properties of noble metals, the dielectric function $`ϵ(\omega )`$ includes also the interband contribution $`ϵ_d(\omega )`$ due to transitions from the d–band to the s–p conduction band. For $`p1`$, the absorption coefficient of such a system is proportional to that of a single particle and is given by
$$\alpha (\omega )=9pϵ_m^{3/2}\frac{\omega }{c}\text{Im}\frac{1}{ϵ_s(\omega )},$$
(12)
where
$$ϵ_s(\omega )=ϵ_d(\omega )\omega _p^2/\omega (\omega +i\gamma _s)+2ϵ_m$$
(13)
plays the role of an effective dielectric function of a particle in the medium. Its zero, $`ϵ_s^{}(\omega _s)=0`$, determines the frequency of the SP, $`\omega _s`$. In Eq. (13), $`\omega _p`$ is the bulk plasmon frequency of the conduction electrons, and the width $`\gamma _s`$ characterizes the SP damping. The semiclassical result Eqs. (12) and (13) applies to nanoparticles with radii $`Rq_{_{TF}}^1`$, where $`q_{_{TF}}`$ is the Thomas–Fermi screening wave–vector ($`q_{_{TF}}^11`$ Å in noble metals). In this case, the electron density deviates from its classical shape only within a surface layer occupying a small fraction of the total volume . Quantum mechanical corrections, arising from the discrete energy spectrum, lead to a width $`\gamma _sv__F/R`$, where $`v__F=k__F/m`$ is the Fermi velocity . Even though $`\gamma _s/\omega _s(q_{_{TF}}R)^11`$, this damping mechanism dominates for sizes $`R<10`$ nm. Additional contributions to such a width come from the e–p, e–e, and electron–impurity interactions as well as from the disorder. On the other hand, in small clusters containing several dozens of atoms, this semiclassical approximation breaks down and density functional or ab initio methods should be used . In particular, discrete electronic levels whose width does not exceed their spacing are expected for particle sizes smaller than a few angstroms. In the latter regime, the e–e interactions are similar to those in atoms and molecules. One therefore expects that the role of the many–body correlations increases as we approach the crossover from quasi–continuous to discrete nanoparticle energy levels.
It should be noted that, in contrast to the scattering with surface collective excitations, the e–e scattering is not too sensitive to the nanoparticle size as long as the condition $`q_{_{TF}}R1`$ holds . Indeed, for such sizes, the static screening is essentially bulk–like. At the same time, the energy dependence of the bulk e–e scattering rate (Eq. 7) $`\gamma _e(EE_F)^2`$ comes from the phase–space restriction due to the energy and momentum conservation, and involves the exchange of typical momenta $`qq_{_{TF}}`$. If the size–induced momentum uncertainty $`\delta qR^1`$ is much smaller than $`q_{_{TF}}`$, the e–e scattering rate in a nanoparticle is not significantly affected by the confinement . Below we will see however that this is not the case for the electron–SP interactions.
### 3.2 Ultrafast Dynamics
Even though the electronic, thermodynamic, and optical properties of metal nanoparticles have been extensively studied, the role of confinement in the electron dynamics is much less understood. Examples of outstanding issues include the role of e–e interactions in the process of cluster fragmentation, the role of surface lattice modes in providing additional channels for intra-molecular energy relaxation, the influence of the electron and nuclear motion on the superparamagnetic properties of clusters, and the effect of confinement on the nonlinear optical properties and transient response under ultrafast excitation . These and other time–dependent phenomena can be studied with femtosecond nonlinear optical spectroscopy, which in these structures provides time resolution shorter than the relaxation times. Similar experimental studies in bulk metals and metal interfaces and surfaces have been reviewed e.g. in Ref. . Here we will focus on the more recent work in metal nanoparticles.
Extensive experimental studies of the electron relaxation in noble–metal nanoparticles have recently been performed using ultrafast pump–probe spectroscopy. In contrast to the situation in semiconductors, the dephasing processes in metals are very fast. In metal nanoparticles, the dephasing time can be deduced, e.g., from the linear absorption SP width to be of the order of a few fs. Such measurements were performed by Klar et. al. , who distinguished between homogeneous and inhomogeneous broadening by measuring directly the SP lineshape of a single nanoparticle using a near–field antenna effect. The decay time of the SP resonance has also been studied with second and third harmonic generation measurements .
Within several femtoseconds after the photoexcitation, the optically–excited electrons and holes form non–equilibrium populations. The pump–induced intraband transitions create a broad non–thermal electron distribution that extends from the Fermi energy up to the pump photon energy, while the interband transitions create additional electron and d-band hole populations when the pump frequency exceeds the d–band to conduction band transition threshold. Within a few fs, the high energy electrons scatter to lower energies due to e–e interactions with the FS electrons. The latter interactions are screened within time intervals of the order of the inverse plasma frequency, typically $``$ 1fs. Note that, similar to MDQW’s , at high electron excess energies the e–e scattering time is of the order of a fs, much shorter than the e–p scattering times (of the order of a ps). However, as such electrons scatter down to the Fermi surface, the e–e scattering times become much longer, of the order of several hundreds of fs. As discussed in the previous section, this is a consequence of the Pauli blocking of the phase space available for scattering and the screening of the e–e interaction, which leads to energy–dependent scattering times given by Eq. 7 . Therefore, one expects that thermalization should slow down as the electrons approach the Fermi surface. In the case of metal films, a non–thermal population consisting of a hot Fermi–Dirac distribution together with a tail of high energy electrons right above the Fermi surface (and a tail of low energy holes below the Fermi surface) was indeed observed in time–resolved photoemission experiments during time intervals of the order of hundreds of fs . Pump–probe experiments also showed a clear signature of delayed thermalization (see e.g. ). In particular, the rise time of the differential transmission and reflectance signal was observed to be of the order of hundreds of fs, much longer than the pulse duration.
A similar slow rise time of the differential transmission was also observed in the case of copper , gold , and silver nanoparticles as well as in gold nanoshells . It has been reported that, in small nanoparticles with diameters $``$5 nm, such a delayed thermalization regime (stage II) lasts longer than in metal films . At the same time, the e–p interaction for such sizes was found to be weaker as compared to the bulk . A possible explanation of the reduced e–p interaction is that, for the small sizes, the characteristic phonon energy, given by the Debye frequency, becomes smaller than the spacing between the nanoparticle energy levels close to the Fermi surface, which is of the order of $`v_F/R`$. As a result, the scattering of an electron via bulk–like phonons is suppressed . As the electrons approach the Fermi surface, the e–e scattering times become comparable to the e–p scattering times and therefore the e–p interaction also contributes to the thermalization of the electron gas. The suppression of the latter could lead to delayed internal thermalization of the electron gas.
The main spectral feature in the differential absorption of the copper and gold nanoparticles (where the onset of interband transitions is very close to SP frequency) was a transient asymmetric broadening of the SP resonance ( see e.g. Refs. ). Perner et. al observed a build–up of such a pump–induced broadening during time intervals of the order of 1ps. Bigot et. al. observed a similar buildup and also pointed out that the differential absorption lineshape could not be understood without including the effects of the energy–dependent e–e scattering on the interband dielectric function. They attributed the SP broadening mainly to the smearing of the Fermi–Dirac electron distribution close to the Fermi surface, due to the heating of the electron gas, which affects the interband dielectric function. Inouye et.al. arrived at similar conclusions and argued that, in the case of resonant interband excitations, the contribution of the intraband e–e scattering to the damping of the SP, which is described by the width $`\gamma _s`$ in the intraband dielectric function in Eq. (13) , plays a minor role.
The time–delay of the SP broadening can be understood as follows. For quasi–equilibrium conditions, the nonlinear absorption spectrum can be described using the linear absorption results but with dielectric constants determined by the time–dependent carrier populations. As discussed above, within the time resolution of the experiment, the electron distribution deviates from the equilibrium Fermi–Dirac distribution in the vicinity of the Fermi surface, due to the high electron temperature and the non–thermal population . As can be seen from the linear absorption expressions, such a smearing of the FS distribution leads to the broadening of the SP resonance. Therefore, the initial rise time of the differential transmission signal is determined by the non–equilibrium electrons ouside the equilibrium Fermi–Dirac distribution. Initially, such electrons occupy the broad photoexcited distribution and therefore their number is relatively small, determined by the pump intensity. However, their number increases with time as the e–e interaction leads to the scattering of more electrons out of the FS. This is the origin of the time–dependent increase in the SP broadening.
Since the electron heat capacity is much smaller than that of the lattice, an electron temperature much higher than that of the lattice can be reached during subpicosecond time scales. Subsequently, the electron and phonon baths equilibrate through the e–p interactions over time intervals of a few picoseconds. During this incoherent stage, the hot electrons can be characterized by a thermalized Fermi–Dirac distribution with time–dependent temperature $`T(t)`$, while the phonons can be characterized by a Bose–Einstein distribution with time–dependent temperature $`T_l(t)`$. During the cooling of the hot electron gas, the SP width decreases . As soon as the electrons have equilibrated among themselves, one can study the subsequent time evolution by using a set of coupled differential equations for the temperatures $`T(t)`$ and $`T_l(t)`$, referred to as the two–temperature model :
$`C_e(T){\displaystyle \frac{T}{t}}=`$ $`\left(\kappa _eT\right)G(TT_l)+P(𝐫,t),`$
$`C_l{\displaystyle \frac{T_l}{t}}=`$ $`G(TT_l),`$ (14)
where $`C_e(T)=\gamma T`$ and $`C_l`$ are the electron and lattice heat capacities, respectively, $`G`$ is the energy transfer coefficient between the electrons and the lattice, which is proportional to the e–p coupling constant , $`\kappa _e`$ is the electronic thermal conductivity, and the source term $`P(𝐫,t)`$ describes the local energy density per unit time absorbed by the electron system from the pump optical pulse. The first term on the rhs of Eq. (3.2) describes the thermal diffusion from the nanoparticle to the surrounding matrix, which occurs on a time scale of several tens of ps. The second term on the rhs of Eq. (3.2) determines the tranfer of heat from the electron gas to the lattice as the former cools down during a time interval of several ps. The above model assumes electron and lattice temperatures larger than the Debye temperature and that the interactions can maintain a quasi–equilibrium for both the electron and the phonon populations at all times. The cooling of the electron gas to the lattice manifests itself via an exponential decay of the differential transmission. According to Eq. (3.2), the decay time of the differential tranmission signal is inversely proportional to the strength of the e–p interaction . Stella et. al. performed transient reflectivity measurements in metallic tin nanoparticles with radii ranging from 2 nm to 6nm and saw a decay of the signal as a function of pump–probe delay that became faster by decreasing the nanoparticle radius. They deduced from their data a contribution to the energy relaxation time during the thermalization of the hot electron Fermi sea with the lattice that was inversely proportional to the nanoparticle radius. In a subsequent paper, Nisoli et. al. performed femtosecond pump–probe measurements in solid and liquid gallium nanoparticles with radii ranging from 5 nm to 9 nm and observed energy relaxation time constants that varied from 1.6 ps to 600 fs with decreasing nanoparticle size. They observed similar electron relaxation dynamics in the solid and liquid nanoparticle phases, which indicates that in this system, the scattering of electrons with bulk phonons plays only a minor role in the relaxation because of the reduction in the available phase space due to the quantum confinement. Instead, the hot electrons transfer their excess energy to the lattice through the generation of surface vibrational waves (capillary waves), which leads to an electron–surface phonon interaction $`G`$ inversely proportional to the nanoparticle radius . An additional contribution to this effect may come from the fact that the heat diffusion to the matrix becomes faster with decreasing nanoparticle size due to a thermal buildup that occurs for larger particles. A similar effect was observed by Halte et. al. in the case of very small silver nanoparticles . The role of surface effects on the dynamics was also discussed in Ref. .
In the gold and copper nanoparticles, condidered above, the SP frequency is very close to the onset of interband transitions. In contrast, in silver nanoparticles, the SP and interband transitions are well separated in energy. Halte et. al. compared the electron dynamics between silver thin films and silver nanoparticles embedded in glass in the same spectral range, and studied its dependence on the photoexcitation intensity. In thin films, they found that the RPA dielectric function corresponding to the instantaneous temperature determined by the two–temperature model reproduced their results quite well. In the case of the silver nanoparticles, the main feature in the differential absorption spectrum was an apparent SP redshift, which was attributed to the changes in the real part of the interband dielectric function due to the thermal broadening of the electron distribution close to the Fermi level. It was also shown that the finite e–e lifetime of the carriers photoexcited via interband transitions was essential for obtaining such a redshift. Finally, the strong non–parabolicity of the conduction band at the energies corresponding to the optical frequencies was shown to be important for interpreting the data. Averitt et. al. performed ultrafast pump–probe measurements in gold nanoshell structures, which consist of a dielectric core surrounded by a thin metallic shell of nanometer dimensions. In such structures, the SP frequency can be tuned by changing the ratio of the core diameter to the shell thickness, and the experiment was performed at frequencies far off the onset of the interband transitions. The delayed pump–induced broadening of the SP resonance suggests that a nonthermal electron population is present in these structures during the initial $``$100 fs following the pump photoexcitation. The decay of the pump–probe signal was found to be somewhat slower than in the bulk, which was attributed to a weaker e–p interaction due to the reduction of phase space available for scattering with bulk phonons induced by the confinement. Finally, they observed a transient blueshift of the SP resonance, attributed to the off–resonant interband transitions.
Recent experimental results in the case of small noble metal particles indicate that many–body correlation effects play an important role during the cooling of the electron gas to the lattice (relaxation stage III). Despite the similarities to the bulk–like behavior, observed, e.g., in metal films, certain aspects of the optical dynamics in nanoparticles are significantly different . For example, the experimental studies of copper nanoparticles by Bigot et. al. revealed that, for sizes smaller than $``$ 5nm, the decay times of the pump–probe signal depend strongly on the probe frequency in the immediate vicinity of the SP resonance. In particular, the relaxation is considerably slower at the SP resonance, and becomes faster right above and right below the SP frequency . This and other observations suggest that collective surface excitations play an important role in the electron dynamics in small metal particles. This important issue will be discussed at length in section 7.
## 4 Coherent Ultrafast Response of the Fermi Edge Singularity: Formalism
### 4.1 Basic equations
Before we proceed with our discussion of the time evolution of the FES in the coherent regime, we outline the main points of a formalism recently developed to account for the non–Markovian dynamics due to the dynamical FS response (e–h correlations). More details may be found in Appendix A.
As discussed in section 2.1, even in linear absorption, the HFA and the dephasing time approximation have serious shortcomings when used to describe the FES. In the case of a photoexcited electron gas, the correlation effects have been treated within the second Born approximation by using Keldysh Green functions (see e.g. Ref. ). However, at low temperatures and within the frequency range of the FES, the effects of the e–h interaction must be consistently accounted for to arbitrary order. For example, as discussed in section 2.1, even in linear absorption the vertex correction diagrams with arbitrarily many crossed e–h interaction lines are as divergent as the HFA ladder diagrams and should therefore be treated on equal footing. With Green functions, this requires summing up at least the parquet diagrams , a formidable task especially in the non–equilibrium femtosecond regime. Alternative methods are therefore desirable. In undoped semiconductors, the “dynamic controlled truncation scheme” has been used to treat the exciton–exciton correlations (see section 1). However, such a hierarchy of density matrix equations no longer truncates if the ground state of the semiconductor includes a FS. Furthermore, as we demonstrate in section 5.1, an expansion in terms of the optical fields breaks down for frequencies within the FES range. Finally, in view of the significant complexity of the problem, it is highly desirable to use a method that also allows for a physically intuitive interpretation of the results. The purpose of this section is to outline the main points of such a method, which will be used in section 5 to describe the coherent nonlinear response of the FES.
We first consider the case of one conduction subband and later extend the formalism to include the Coulomb coupling to a second subband. We also consider spinless electrons for simplicity. In the rotating frame , the Hamiltonian describing this system is
$$H_{\mathrm{tot}}(t)=H+H_p(t)+H_s(t).$$
(15)
The first term is the “bare” Hamiltonian,
$$H=\underset{𝐪}{}\epsilon _𝐪^vb_𝐪^{}b_𝐪+\underset{𝐪}{}(\epsilon _𝐪^c+\mathrm{\Omega })a_𝐪^{}a_𝐪+V_{ee}+V_{hh}+V_{eh},$$
(16)
where $`a_𝐪^{}`$ is the creation operator of a conduction electron with energy $`\epsilon _𝐪^c`$ and mass $`m_e`$, $`b_𝐪^{}`$ is the creation operator of a valence hole with energy $`\epsilon _𝐪^v`$ and mass $`m_h`$, $`V_{ee},V_{eh}`$, and $`V_{hh}`$ describe the e–e, e–h, and h–h interactions, respectively, and $`\mathrm{\Omega }=E_g+E_F(1+m_e/m_h)\omega _p`$ is the detuning of the central frequency of the optical fields, $`\omega _p`$, from the Fermi level, $`E_g`$ being the bandgap (we set $`\mathrm{}=1`$ everywhere). The second and third terms describe the coupling of the pump optical field, $`_p(t)e^{i𝐤_p𝐫i\omega _pt}`$, and the probe optical field, $`_s(t)e^{i𝐤_s𝐫i\omega _p(t\tau )}`$, respectively:
$`H_p(t)`$ $`=\mu _p(t)[e^{i𝐤_p𝐫}U^{}+\mathrm{H}.\mathrm{c}.],`$
$`H_s(t)=`$ $`\mu _s(t)[e^{i𝐤_s𝐫+i\omega _p\tau }U^{}+\mathrm{H}.\mathrm{c}.],`$ (17)
where the pump amplitude $`_p(t)`$ is centered at time $`t=0`$ and the probe amplitude $`_s(t)`$ is centered at the time delay $`t=\tau `$, $`\mu `$ is the interband transition matrix element, and
$$U^{}=\underset{𝐪}{}a_𝐪^{}b_𝐪^{}$$
(18)
is the optical transition operator. The conventions for the time delay $`\tau `$ are clarified in Appendix B.
In many experiments, the amplitude of the probe field is much smaller than that of the pump, $`|_p(t)||_s(t)|`$. In that case, as was shown in Ref. (see Appendix A), the experimentally–measurable nonlinear optical polarization can be obtained in terms of the linear response functions of a “pump–dressed” semiconductor to a probe field (note that, within $`\chi ^{(3)}`$, this is true even for comparable pulse amplitudes). This “dressed” system is described by a time–dependent effective Hamiltonian $`\stackrel{~}{H}(t)`$, which is obtained by performing a time–dependent Schrieffer–Wolff/Van Vleck canonical transformation on the Hamiltonian $`H+H_p(t)`$ . As we shall see later, in all the cases of interest, the effective Hamiltonian $`\stackrel{~}{H}(t)`$ has the same operator form as the bare Hamiltonian $`H`$, with the important difference that the band dispersions (effective masses) and interaction potentials are time–dependent through $`_p(t)`$. Thus, the calculation of the nonlinear absorption spectrum reduces to that of the linear absorption spectrum of the “pump–dressed” semiconductor with uncoupled “effective bands” — a great simplification that allows us to use straightforward generalizations of well established theoretical tools in order to treat the correlations. In fact, such an approach mimics nicely the spirit of the pump–probe experiments and allows for a physically intuitive interpetation of the results. It is important to note here that this “pump–dressed semiconductor” approach is not restricted to monochromatic pulses and is valid for any pulse duration .
The pump–probe nonlinear polarization has the following form (see Appendix A):
$$P_{𝐤_s}(t)=i\mu ^2e^{i𝐤_s𝐫i\omega _p(t\tau )}_{\mathrm{}}^t𝑑t^{}_s(t^{})\mathrm{\Phi }_0(t)|\stackrel{~}{U}(t)𝒦(t,t^{})\stackrel{~}{U}^{}(t^{})|\mathrm{\Phi }_0(t^{}).$$
(19)
Here, $`|\mathrm{\Phi }_0(t)`$ is the state evolved with $`\stackrel{~}{H}(t)`$ from the semiconductor ground state $`|0`$ of $`H`$, $`\stackrel{~}{U}^{}(t)`$ is the effective optical transition operator describing the probability amplitude for the photoexcitation of an e–h pair by the probe field in the presence of the pump excitation, and $`𝒦(t,t^{})`$ is the time–evolution operator satisfying the Schrödinger equation
$$i\frac{}{t}𝒦(t,t^{})=\stackrel{~}{H}(t)𝒦(t,t^{}).$$
(20)
The above equation describes the time evolution of a probe–photoexcited e–h pair in the presence of the pump excitation. The effective Hamiltonian and effective transition operator are given by (see Appendix A)
$$\stackrel{~}{H}(t)=H_0+\frac{\mu }{2}(_p(t)[\widehat{𝒫}(t),U^{}]+\mathrm{H}.\mathrm{c}.),$$
(21)
and
$$\stackrel{~}{U}^{}(t)=U^{}+\frac{1}{2}[\widehat{𝒫}(t),[U^{},\widehat{𝒫}^{}(t)]]+\frac{1}{2}[\widehat{𝒫}^{}(t),[U^{},\widehat{𝒫}(t)]],$$
(22)
where the operator $`\widehat{𝒫}^{}(t)`$, which generates the canonical transformation, satisfies the equation
$$i\frac{\widehat{𝒫}^{}(t)}{t}=[H,\widehat{𝒫}^{}(t)]+\mu _p(t)U^{},$$
(23)
with the initial condition $`\widehat{𝒫}^{}=0`$ before the pump arrives. Eqs. (21) and (22) include all the pump–induced corrections to $`\stackrel{~}{U}^{}(t)`$ and $`\stackrel{~}{H}(t)`$ up to the second order in the pump optical field and are valid when $`\left(\mu E_p/\mathrm{\Omega }\right)^2<1`$ (for off–resonant excitation) or $`\left(\mu E_pt_p\right)^2<1`$ (for resonant excitation), where $`t_p`$ is the pump duration.
It should be emphasized that, although Eq. (21) gives the effective Hamiltonian up to the second order in the pump field, $`_p_p^{}`$, the polarization expression Eq. (19) describes the effects of $`\stackrel{~}{H}`$ in all orders. For example, as we shall see in section 4.3, the pump–induced term in $`\stackrel{~}{H}`$ contains self–energy corrections to the electron/hole energies, which describe (among other effects) the resonance blueshift due to the ac–Stark effect. Although the magnitude of these self–energy corrections, calculated using Eq. (21), is quadratic in $`_p`$, the correct position of the resonance can only be obtained by evaluating the PP polarization (19) nonperturbatively (beyond $`\chi ^{(3)}`$), i.e., without resorting to the expansion of the time–evolution operator $`𝒦(t,t^{})`$ in the pump field. Importantly, as we demonstrate in 5.1, the same is true when calculating the effects of the self–energy corrections on the e–h correlations. As we shall see in section 5, such a nonperturbative (in the pump field) treatment of the nonlinear response of the FES is crucial for the adequate description of the PP spectrum at negative time delays. In Sections 4.2 and 4.4, we will describe the corresponding procedure, which accounts for the FS dynamical response. In contrast, the third–order polarization, $`\chi ^{(3)}`$, can be simply obtained from Eq. (19) by expanding $`𝒦(t,t^{})`$ to the first order in the pump–induced term in $`\stackrel{~}{H}`$ \[second term in Eq. (21)\]. We did not include in Eq. (19) the biexcitonic contribution (coming from the excitation of two e–h pairs by the pump and the probe pulses) since it vanishes for the negative time delays ($`\tau <0`$) considered below where the coherent effects dominate .
In addition to including important contributions beyond $`\chi ^{(3)}`$ via the solution of Eq. (20) as discussed above, the advantage of Eq. (19), as compared to the equations of motion for the polarization, comes from its similarity to the linear polarization that determines the linear absorption spectrum . This can be seen by setting $`_p(t)=0`$ in the Eqs. (21) and (22), in which case the effective time–evolution and optical transition operators transform into their “bare” counterparts: $`𝒦(t,t^{})e^{iH(tt^{})}`$ and $`\stackrel{~}{U}^{}(t)U^{}`$. Moreover, like $`U^{}`$, the effective transition operator $`\stackrel{~}{U}^{}(t)`$ creates a single e–h pair, while, as we shall see in section 4.3, the effective Hamiltonian $`\stackrel{~}{H}(t)`$ can be cast in a form similar to $`H`$. This allows one to interpret the Fourier transform of Eq. (19) as the linear absorption spectrum of a “pump–dressed” semiconductor with two uncoupled but time–dependent effective bands. This mapping simplifies significantly the calculation of the FES ultrafast nonlinear optical response by allowing a straightforward generalization of the CCE. It also allows one to interpret the various dynamical features in the nonlinear absorption spectra, originating from the correlation effects, within the familiar framework developed for linear spectroscopy.
### 4.2 Overview of the Coupled Cluster Expansion
In this section, we show how the time–dependent CCE can be used to study the effects of the e–h correlations (dynamical FS response) on the time evolution of the e–h pair photoexcited by the probe. Our goal is to evaluate the many–body state $`|\mathrm{\Psi }(t)=𝒦(t,t^{})\stackrel{~}{U}^{}(t^{})|\mathrm{\Phi }_0(t^{})`$ that enters in Eq. (19). This state satisfies the Schrödinger equation
$$i\frac{}{t}|\mathrm{\Psi }(t)=\stackrel{~}{H}(t)|\mathrm{\Psi }(t).$$
(24)
As already mentioned, $`\stackrel{~}{H}(t)`$ has the same form as the bare Hamiltonian $`H`$. This allows us to obtain $`|\mathrm{\Psi }(t)`$ through a straightforward generalization of the linear absorption calculations . After eliminating the valence hole degrees of freedom , $`|\mathrm{\Psi }(t)`$ is expressed in the CCE form
$$|\mathrm{\Psi }(t)=e^{S(t)}|\mathrm{\Phi }(t),$$
(25)
where the time–dependent operator $`S(t)`$ creates FS–pairs and is given by
$`S(t)=`$ $`{\displaystyle \underset{p>k_F,k<k_F}{}}s(𝐩,𝐤,t)a_𝐩^{}a_𝐤`$ (26)
$`+{\displaystyle \underset{p_1,p_2>k_F,k_1,k_2<k_F}{}}s_2(𝐩_\mathrm{𝟏},𝐩_\mathrm{𝟐},𝐤_\mathrm{𝟏},𝐤_\mathrm{𝟐},t)a_{𝐩_\mathrm{𝟏}}^{}a_{𝐩_\mathrm{𝟐}}^{}a_{𝐤_\mathrm{𝟐}}a_{𝐤_\mathrm{𝟏}}+\mathrm{},`$
while the state $`|\mathrm{\Phi }(t)`$, discussed in Section 4.4, describes the time evolution of the probe–induced e–h pair. In Eq. (26), the amplitude $`s(𝐩,𝐤,t)`$ describes the e–h correlations which, in particular, are responsible for the unbinding of the HFA bound state; the two–pair amplitude $`s_2`$ describes the e–e interaction effects at the RPA level and beyond. ¿From a physical point of view, the operator $`e^{S(t)}`$ describes the readjustment of the FS density profile during the optical transition in response to the FS interactions with the photoexcited e–h pair.
Substituting Eq. (25) into the Schrödinger equation Eq. (24), multiplying by the operator $`e^{S(t)}`$ from the lhs, and using the fact that $`[S(t),S(t^{})]=[\dot{S}(t),S(t^{})]=0`$, one obtains
$$i\frac{}{t}|\mathrm{\Phi }(t)+i\dot{S}(t)|\mathrm{\Phi }(t)=e^{S(t)}\stackrel{~}{H}(t)e^{S(t)}|\mathrm{\Phi }(t),$$
(27)
where the transformed Hamiltonian on the rhs can be expressed in terms of the commutator series (Baker–Campbell–Hausdorff expansion)
$$e^{S(t)}\stackrel{~}{H}(t)e^{S(t)}=\stackrel{~}{H}(t)+[\stackrel{~}{H}(t),S(t)]+\frac{1}{2}[[H,S(t)],S(t)]+\mathrm{}$$
(28)
An important advantage of the CCE is that, due to the FS momentum restrictions in Eq. (26), the above series terminates after the first few terms (three for quartic interactions) and a closed–form expression of the transformed Hamiltonian (28) can be obtained in terms of $`S(t)`$. By requiring that all FS–pair creation processes are eliminated from the above equation, one obtains the CCE equations for $`S(t)`$. Before proceeding with such a calculation however, one needs to derive explicit expressions for $`\stackrel{~}{H}(t)`$ and $`\stackrel{~}{U}(t)`$.
### 4.3 Effective Hamiltonian and Transition Matrix Elements
#### 4.3.1 Discussion for Doped Semiconductors
In the general case, the effective Hamiltonian $`\stackrel{~}{H}(t)`$ is given by Eq. (21) and the effective optical transition operator $`\stackrel{~}{U}(t)`$ by Eq. (22). For calculations, it is useful to re–express such equations in second–quantized form. In the case of undoped semiconductors, it was shown in Ref. that the effective Hamiltonian $`\stackrel{~}{H}(t)`$ describes the same interactions as the bare Hamiltonian $`H`$, however among quasiparticles with time–dependent properties determined by the pump polarization. Similarly, the effective interband transition matrix element in $`\stackrel{~}{U}(t)`$ was shown to have an imaginary part that describes the dephasing due to the exciton–exciton interactions. The purpose of this section is to present similar results in the case of a doped semiconductor .
We start with Eqs. (21) and (22), which express $`\stackrel{~}{H}(t)`$ and $`\stackrel{~}{U}(t)`$ in terms of the canonical transformation operator $`\widehat{𝒫}^{}(t)`$. The latter is given by Eq. (23), which includes the effects of the Coulomb interactions on the pump photoexcitation. It is important to realize that the effect of the e–h interactions on the pump and the probe photoexcitations is very different. For an adequate description of the FES, the e–h interactions should be taken into account non–perturbatively for the probe–photoexcited pair. Indeed, the nonlinear absorption spectrum at a given frequency $`\omega `$ close to the FES resonance is determined by the time–evolution for long times (of the order of the dephasing time $`T_2`$) of an e–h pair photoexcited by the probe at the energy $`\omega `$ . Since the characteristic “e–h interaction time” $`E_M^1`$ (inverse HFA bound state energy) that determines the non–exponential polarization decay of the FES is much shorter that the dephasing time, the long–time asymptotics of the response function (to the probe) depends non–perturbatively on the e–h interactions. In contrast, a short pump optical pulse excites a wavepacket of continuum e–h pair states (unlike in the discrete exciton case) with energy width $`t_p^1`$, which thus evolves during timescales comparable to the pulse duration $`t_p`$. Also, the corrections to the effective Hamiltonian are determined by the time evolution of the pump–induced carriers only up to times $`t_p`$ \[see Eq. (21)\]. Therefore, if the “e–h interaction time” $`E_M^1`$ is larger than $`t_p`$, $`t_pE_M<1`$, (i.e., if the pump pulse frequency width exceeds $`E_M`$), the e–h interactions can be treated perturbatively when describing the time–evolution of the pump–photoexcited pairs. This can also be shown explicitly for the third–order nonlinear polarization. In the general expression for $`\chi ^{(3)}`$, all contributions that depend on the pump are integrated over the width of the pump pulse; therefore, any resonant enhancement of $`\chi ^{(3)}`$ that depends on the pump frequency will be broadened out for sufficiently short pulses with frequency width that exceeds $`E_M`$. In other words, when deriving the pump–renormalized parameters, one can treat Coulomb interactions perturbatively if the above condition is fulfilled. In fact, the above situation is somewhat similar to the calculation of the linear absorption spectrum close to the indirect transition threshold, where perturbation theory can be used . Thus the above consideration applies even for long pulse durations provided that the detuning $`\mathrm{\Omega }`$ exceeds $`E_M`$. However, in order to obtain the full absorption spectrum, the time–evolution of the probe–photoexcited pair with such effective Hamiltonian (with perturbatively calculated time–dependent parameters) should be treated non–perturbatively. As can be seen from the above discussion, an important advantage of this formalism is that it seprates naturally between the perturbative and the non–perturbative interaction effects.
#### 4.3.2 Second Quantization Expressions
We now proceed with the derivation of the effective Hamiltonian. To lowest order in the interactions, $`\widehat{𝒫}^{}(t)`$ can be presented as
$`\widehat{𝒫}^{}(t)=`$ $`{\displaystyle \underset{𝐪}{}}𝒫_{eh}(𝐪,t)a_𝐪^{}b_𝐪^{}+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐩,𝐩^{},𝐤}{}}𝒫_{eh}^e(\mathrm{𝐩𝐩}^{};𝐤;t)b_{𝐤𝐩𝐩^{}}^{}a_𝐩^{}a_𝐩^{}^{}a_𝐤`$ (29)
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{𝐩,𝐩^{},𝐤}{}}𝒫_{eh}^h(\mathrm{𝐩𝐩}^{};𝐤;t)a_{𝐩+𝐩^{}𝐤}^{}b_𝐩^{}b_𝐩^{}^{}b_𝐤,`$
where $`𝒫_{eh}`$ is the probability amplitude for excitation of an e–h pair with zero momentum satisfying
$$i\frac{}{t}𝒫_{eh}(𝐪,t)=\left[\mathrm{\Omega }+\epsilon _𝐪^{(c)}+\epsilon _𝐪^{(v)}\right]𝒫_{eh}(𝐪,t)+\mu _p(t)\underset{𝐪^{}}{}v(𝐪𝐪^{})𝒫_{eh}(𝐪^{},t).$$
(30)
In Eq. (29),
$`𝒫_{eh}^e(\mathrm{𝐩𝐩}^{};𝐤;t)`$ $`=i{\displaystyle _{\mathrm{}}^t}𝑑t^{}e^{i(tt^{})(\mathrm{\Omega }+\epsilon _𝐩^c+\epsilon _𝐩^{}^c\epsilon _𝐤^c+\epsilon _{𝐩+𝐩^{}𝐤}^v)}`$
$`\times `$ $`[v(𝐩𝐤)[𝒫_{eh}(𝐩^{},t^{})𝒫_{eh}(𝐩+𝐩^{}𝐤,t^{})](𝐩𝐩^{})]`$ (31)
describes the scattering of the photoexcited e–h pair with an electron, and
$`𝒫_{eh}^h(\mathrm{𝐩𝐩}^{};𝐤;t)`$ $`=i{\displaystyle _{\mathrm{}}^t}𝑑t^{}e^{i(tt^{})\left(\mathrm{\Omega }+\epsilon _{𝐩+𝐩^{}𝐤}^c+\epsilon _𝐩^{}^v+\epsilon _𝐩^v\epsilon _𝐤^v\right)}`$
$`\times `$ $`[v(𝐩𝐤)[𝒫_{eh}(𝐩^{},t^{})𝒫_{eh}(𝐩+𝐩^{}𝐤,t^{})](𝐩𝐩^{})]`$ (32)
describes the scattering of the photoexcited e–h pair with a hole. The above expressions describe in the lowest order in the screened interaction $`\upsilon (𝐩𝐤)`$ the coherent pump–induced processes, the effects of the Hartree–Fock pair–pair and pair–FS interactions, and the dynamical FS response to the pump photoexcitation.
By substituting Eq. (29) into Eq. (22), we obtain the following expression for the effective optical transition operator:
$`\stackrel{~}{U}^{}(t)|\mathrm{\Phi }_0(t)=`$ $`{\displaystyle \underset{p>k_F}{}}M_𝐩(t)a_𝐩^{}b_𝐩^{}|0`$ (33)
$`+{\displaystyle \frac{1}{4}}{\displaystyle \underset{p,p^{}>k_F,k<k_F}{}}M_{\mathrm{𝐩𝐩}^{}𝐤}(t)a_𝐩^{}a_𝐩^{}^{}b_{𝐤𝐩𝐩^{}}^{}a_𝐤|0,`$
where the effective matrix element $`M_𝐩(t)`$ includes corrections due to phase space filling and Hartree–Fock interactions, and $`M_{\mathrm{𝐩𝐩}^{}𝐤}(t)`$ is the probability amplitude for indirect optical transitions induced by the pump optical field, which contribute to the pump–probe polarization in the second order in the interactions. The explicit expressions for $`M_𝐩(t)`$ and $`M_{\mathrm{𝐩𝐩}^{}𝐤}(t)`$ are given in Appendix C.
We turn now to the effective Hamiltonian $`\stackrel{~}{H}(t)`$. After substituting Eq. (29) into Eq. (21) we obtain that
$`\stackrel{~}{H}(t)={\displaystyle \underset{𝐪}{}}\epsilon _𝐪^v(t)b_𝐪^{}b_𝐪+{\displaystyle \underset{𝐪}{}}\epsilon _𝐪^c(t)a_𝐪^{}a_𝐪+V_{eh}(t)+V_{ee}(t),`$ (34)
where
$`\epsilon _𝐪^c(t)=\epsilon _𝐪^c\mu _p(t)\mathrm{Re}\left[𝒫_{eh}(𝐪,t){\displaystyle \underset{𝐪^{}}{}}𝒫_{eh}^e(\mathrm{𝐪𝐪}^{};𝐪;t)\right]`$ (35)
is the effective conduction electron energy;
$`\epsilon _𝐪^v(t)=E_g+\epsilon _𝐪^v\mu _p(t)\mathrm{Re}\left[𝒫_{eh}(𝐪,t){\displaystyle \underset{𝐪^{}}{}}𝒫_{eh}^h(\mathrm{𝐪𝐪}^{};𝐪;t)\right],`$ (36)
is the effective valence hole energy;
$`V_{eh}(t)={\displaystyle \underset{\mathrm{𝐤𝐤}^{}𝐪}{}}\upsilon _{eh}(𝐪;\mathrm{𝐤𝐤}^{};t)a_{𝐤+𝐪}^{}a_𝐤b_{𝐤^{}𝐪}^{}b_𝐤^{},`$ (37)
is the effective e–h interaction; and
$`V_{ee}(t)={\displaystyle \frac{1}{2}}{\displaystyle \underset{\mathrm{𝐤𝐤}^{}𝐪}{}}\upsilon _{ee}(𝐪;\mathrm{𝐤𝐤}^{};t)a_{𝐤+𝐪}^{}a_{𝐤^{}𝐪}^{}a_𝐤^{}a_𝐤,`$ (38)
is the effective e–e interaction. The explicit expressions for $`\upsilon _{eh}(𝐪;\mathrm{𝐤𝐤}^{};t)`$ and $`\upsilon _{ee}(𝐪;\mathrm{𝐤𝐤}^{};t)`$ are given in Appendix C. As can be seen, $`\stackrel{~}{H}(t)`$ has the same operator form as the bare Hamiltonian $`H`$. However, both the effective band dispersions and the effective interaction potentials are now dependent on time. Note here that the above pump–induced renormalizations only last for the pulse duration $`t_p`$. As discussed above, they are therefore perturbative in the screened interactions for $`t_pE_M<1`$ or for $`\mathrm{\Omega }>E_M`$.
Let us first discuss the effect of the pump–induced self–energy corrections to the conduction and valence band energies, given by the last terms in Eqs. (35) and (36). The dispersion of the effective band is shown in Fig. 1. As can be seen, the pump pulse leads to a bandgap increase as well as a change in the momentum dependence (band dispersion) that last as long as the pump pulse. The magnitude of the bandgap increase is of the order of $`(\mu _p)^2/\mathrm{\Omega }`$ (for off–resonant excitation) and $`(\mu _p)^2t_p`$ (for resonant excitation) and leads to, e.g., the ac–Stark blueshift. As we shall see, for pulse duration shorter than the dephasing time, it also leads to bleaching and gain right below the onset of absorption, analogous to the case of excitons or two–level systems. It should be emphasized that these are coherent effects that should not be confused with the incoherent bandgap redshift due to the e–e interactions among real photoexcited carriers . In particular, the above bandgap renormalization is induced by the transverse EM–field of the laser, as compared to the usual bandgap renormalization due to a longitudinal EM–field, i.e., Coulomb screening. The change in the band dispersion, whose relative magnitude is of the order of $`(\mu _p/\mathrm{\Omega })^2`$ (for off–resonant excitation) or $`(\mu _pt_p)^2`$ (for resonant excitation), can be viewed as an increase in the effective density of states and, to the first approximation, in the effective mass. This is important in doped semiconductors because, as we shall see later, it leads to an optically–induced time–dependent enhancement of the e–h interactions and scattering processes with the FS electrons.
The effective Hamiltonian $`\stackrel{~}{H}(t)`$ also includes pump–induced corrections in the effective interaction potentials, determined by the pair–pair and pair–FS interactions during the pump photoexcitation. By expanding Eqs. (128) and (129) for carrier energies close to the Fermi surface using Eqs. (4.3.2) and (4.3.2), one can show that these corrections vanish at the Fermi surface; for the typical FS excitation energies $`\mathrm{\Delta }\epsilon E_M`$ that contribute to the FES, their order of magnitude is $`(\mu _p\mathrm{\Delta }\epsilon /\mathrm{\Omega }^2)^2`$ (for off–resonant excitation) or $`(\mu _p\mathrm{\Delta }\epsilon t_p^2)^2`$ (for resonant excitation). Thus the corrections to the interaction potentials are suppressed for below–resonant excitation by a factor of $`(E_M/\mathrm{\Omega })^2`$, or for short pulses by a factor of $`(E_Mt_p)^2`$, as compared to the self–energy corrections. Such a suppression is due to the Pauli blocking effect and the screening, which leads to the vanishing of the pump–induced corrections to the interaction potentials at the Fermi surface. Similarly, the pump–induced indirect optical transition matrix elements $`M_{\mathrm{𝐩𝐩}^{}𝐤}(t)`$ are suppressed by the same factor as compared to the direct transition matrix element $`M_𝐩(t)`$ \[first term in Eq. (33)\], while they contribute to the pump–probe polarization only in the second order in the screened interactions. Therefore, in the doped case, the screened Coulomb interaction leads to subdominant parameter renormalizations to the effective Hamiltonian $`\stackrel{~}{H}(t)`$ and transition operator $`\stackrel{~}{U}^{}(t)`$ for sufficiently short pump pulses or for off–resonant excitation. In contrast, in undoped semiconductors, the exciton–exciton interactions lead to non–perturbative contributions discussed in Ref. .
### 4.4 Description of the Electron–Hole Pair Dynamics
In this section, we present the final formulae for the nonlinear pump–probe polarization of the interacting system by applying the CCE to the effective Hamiltonian $`\stackrel{~}{H}(t)`$ in order to treat the dynamical FS response. The CCE equation Eq. (27) contains the operator $`S(t)`$ described by a hierarchy of coupled equations for the amplitudes $`s,s_2,\mathrm{}`$, defined by Eq. (26). As discussed at length above, in the coherent regime of interest here, the e–e scattering effects are suppressed, and the nonlinear absorption spectrum is dominated by the e–h interactions. This allows us to use the dephasing time approximation for treating the probe–induced e–e scattering processes , in which case the above hierarchy terminates after $`s`$. Importantly, the e–h correlations (dynamical FS response) are still treated non-perturbatively, since they are determined by $`s`$ . Then all FS–pair creation processes can be eliminated explicitely from the rhs of Eq. (27), leading to the following nonlinear differential equation for the one–FS–pair scattering amplitude , $`s(𝐩,𝐤,t)`$:
$`i{\displaystyle \frac{s(𝐩,𝐤,t)}{t}}`$ $`\left[\epsilon _𝐩^c(t)\epsilon _𝐤^c(t)\right]s(𝐩,𝐤,t)=`$ (39)
$`V\left[1+{\displaystyle \underset{p^{}>k_F}{}}s(𝐩^{},𝐤,t)\right]\left[1{\displaystyle \underset{k^{}<k_F}{}}s(𝐩,𝐤^{},t)\right].`$
The e–h scattering processes described by the above equation are sketched in Fig. 2a. Here $`V`$ is the s–wave component of the screened interaction , approximated for simplicity by its value at the Fermi energy . This neglects plasmon effects, which are however small within the frequency range of the FES . Although this approximation is standard for the linear absorption case, its justification for the transient spectra requires more attention. Indeed, the characteristic time for screening buildup is of the order $`\mathrm{\Omega }_p^1`$, where $`\mathrm{\Omega }_p`$ is the typical plasma frequency corresponding to the FS . This time is however shorter than the typical pump duration $``$ 100fs and dephasing time (which is of the order of ps near the Fermi energy ), so the Coulomb interactions can be considered screened also for the nonlinear absorption case. In Eq. (39), we neglected the hole recoil energy contribution to the excitation energy since $`\epsilon _{k_F}^v(t)\epsilon _0^v(t)<E_M`$ due to a sufficiently heavy hole mass . Note that, by increasing the hole effective mass, the pump–induced hole self–energy, Eq. (36), reduces the hole recoil energy and thus the corresponding broadening . In real samples, the relaxation of the momentum conservation condition due to the disorder will also suppress the hole recoil broadening effects. In fact, the disorder can even lead to localized hole states , which corresponds to the infinite hole mass limit.
From Eq. (27) we then easily obtain the following expression for $`|\mathrm{\Phi }(t)`$:
$$|\mathrm{\Phi }(t)=\underset{p>k_F}{}\mathrm{\Phi }_𝐩(t,t^{})a_𝐩^{}b^{}|0,$$
(40)
where $`b^{}`$ is the creation operator of the zero–momentum hole state and the e–h pair wavefunction $`\mathrm{\Phi }_𝐩(t;t^{})`$ satisfies the “Wannier–like” equation of motion,
$`i{\displaystyle \frac{\mathrm{\Phi }_𝐩(t,t^{})}{t}}=[\mathrm{\Omega }+\epsilon _𝐩^c(t)+\epsilon _𝐩^v(t)`$ $`ϵ_A(t)i\mathrm{\Gamma }_𝐩]\mathrm{\Phi }_𝐩(t,t^{})`$
$`\stackrel{~}{V}(𝐩,t){\displaystyle \underset{p^{}>k_F}{}}\mathrm{\Phi }_𝐩^{}(t,t^{}),`$ (41)
where
$$\stackrel{~}{V}(𝐩,t)=V\left[1\underset{k^{}<k_F}{}s(𝐩,𝐤^{},t)\right],$$
(42)
is the effective e–h potential whose time– and momentum–dependence is determined by the response of the FS electrons to their interactions with the probe–induced e–h pair (vertex corrections) \[sketched schematically in Fig. 2(b), responsible for the unbinding of the HFA bound state\],
$$ϵ_A(t)=V\underset{k^{}<k_F}{}\left[1+\underset{p^{}>k_F}{}s(𝐩^{},𝐤^{},t)\right]$$
(43)
is the self–energy due the to the sudden appearance of the photoexcited hole, which leads to non–exponential polarization decay \[described by $`\mathrm{Im}ϵ_A(t)`$\] due to the Anderson orthogonality catastrophe and a dynamical resonance redshift \[described by $`\mathrm{Re}ϵ_A(t)`$\], and the dephasing width $`\mathrm{\Gamma }_𝐩`$ describes all additional dephasing processes (due to e–e interactions, hole recoil, and phonons). Eq. (4.4) should be solved with the initial condition $`\mathrm{\Phi }_𝐩(t^{},t^{})=M_𝐩(t^{})`$, where the effective matrix element $`M_𝐩(t^{})`$ is defined by Eq. (33).
It is worth stressing here the analogy between Eqs. (40) and (4.4) and the corresponding problem in undoped semiconductors. Indeed, Eq. (40) is the direct analog of an exciton state, whereas again Eq. (4.4) is very similar to a Wannier equation. However, Eqs. (39) and (4.4) include the effects of the interactions between the probe–photoexcited e–h pair and the FS–excitations, and the wavefunction $`\mathrm{\Phi }_𝐩(t,t^{})`$ describes the propagation of the photoexcited pair “dressed” by the FS excitations. Such a “dressing” is due to the dynamical FS response, which leads to the dynamical screening of the effective e–h interaction Eq. (42). The time–dependence of the latter is determined by the FS scattering amplitude $`s(𝐩,𝐤,t)`$ and is affected by the pump excitation as described by Eq. (39). One can easily verify that by setting $`s(𝐩,𝐤,t)=0`$ in Eq. (4.4), we recover the results of the Hartree–Fock (ladder diagram, static FS ) approximation. If one neglects the nonlinear (quadratic) term in Eq. (39), one recovers the three–body (Fadeev) equations . Note that the coupled equations for $`\mathrm{\Phi }_𝐩(t,t^{})`$ and $`s(𝐩,𝐤,t)`$, obtained by neglecting the multipair excitations in Eq. (26), can be extended to include the hole recoil–induced corrections .
Using Eqs. (25), (26), and (40), we now can express the pump–probe polarization Eq. (19) in terms of the e–h wavefunction $`\mathrm{\Phi }`$ and the effective transition matrix element $`M_𝐩`$. Assuming, for simplicity, a delta–function probe pulse centered at time delay $`\tau `$, $`_s(t)=_s\delta (t\tau )`$, we obtain a simple final expression:
$`P_{𝐤_s}(t)=i\theta (t\tau )\mu ^2_se^{i𝐤_s𝐫i\omega _p(t\tau )}{\displaystyle \underset{p>k_F}{}}M_𝐩(t)\mathrm{\Phi }_𝐩(t,\tau ).`$ (44)
Eq. (44) expresses the pump–probe polarization in terms of two physically distinct contributions. First is the effective transition matrix element $`M_𝐩(t)`$, which includes the effects of pair–pair and pair–FS interactions and Phase space filling effects due to the pump–induced carriers present during the probe photoexcitation. Second is the wavefunction $`\mathrm{\Phi }_𝐩(t)`$ of the e–h pair photoexcited by the probe, whose time dependence, determined by Eqs. (4.4) and (39), describes the formation of the absorption resonance. Despite the formal similarities, there are two important differences between the doped and the undoped cases. First, in the doped case, the time evolution of the e–h wavefunction $`\mathrm{\Phi }_𝐩(t)`$ is strongly affected by the interplay between the e–h correlations and the pump–induced transient changes in the bandgap and band dispersion relations. As we shall see later, this can be viewed as an excitation–induced dephasing. Second, unlike in the undoped case , the pump–induced corrections in the effective matrix element $`M_𝐩(t)`$ are perturbative in the screened interactions if the pump detuning or the pump frequency width exceed the Coulomb energy $`E_M`$.
### 4.5 The Case of Two Coupled Subbands
In this section, the above results are extended to the multisubband case in order to investigate the ultrafast PP dynamics of the FES–exciton hybrid formed in asymmetric QW’s with partially occupied subbands . In such one–sided MDQW structures, interband optical transitions from the valence band to several conduction subbands are allowed due to the finite overlap between the hole and electron envelope wave–functions. When the exciton Fano resonance from a higher empty conduction subband is nearly resonant (within a few meV) to the Fermi level, the FES coming from the lowest occupied subband is enhanced by over two orders of magnitude . Such many–body effects on the linear spectrum have been described by using the simple two–subband Hamiltonian
$$H=\underset{i𝐤}{}ϵ_{i𝐤}^ca_{i𝐤}^{}a_{i𝐤}+\underset{𝐤}{}(ϵ_𝐤^v+E_g)b_𝐤^{}b_𝐤\underset{ij}{}\underset{\mathrm{𝐩𝐤𝐪}}{}v_{ij}(𝐪)a_{i𝐩+𝐪}^{}a_{j𝐩}b_{𝐤𝐪}^{}b_𝐤,$$
(45)
where $`a_{i𝐤}^{}`$ and $`ϵ_{i𝐤}^c`$ are the creation operator and the energy of a conduction electron in the $`i`$th subband, $`b_𝐤^{}`$ and $`ϵ_𝐤^v`$ are those of a valence hole ($`E_g`$ is the bandgap), and $`v_{ij}(𝐪)`$ is the screened e–h interaction matrix with diagonal (off–diagonal) elements describing the intrasubband (intersubband) scattering. Due to the screening, the interaction potential is short–ranged and can be replaced by its s–wave component ; close to the Fermi surface, $`v_{ij}(𝐪)v_{ij}`$ . Here we consider the case where only the first subband is occupied, but the Fermi level is close to the exciton level (with binding energy $`E_B`$) below the bottom of the second subband. For large values of the FES–exciton splitting $`\mathrm{\Delta }E_FE_B`$, where $`\mathrm{\Delta }`$ is the subband separation, the linear absorption spectrum consists of two well separated peaks, the lower corresponding to the FES from subband 1, and the higher corresponding to the Fano resonance from the exciton of subband 2 broadened by its coupling to the continuum of states in subband 1. With decreasing $`\mathrm{\Delta }E_FE_B`$, the FES and the exciton become hybridized due to the intersubband scattering arising from the Coulomb interaction. This results in the transfer of oscillator strength from the exciton to the FES and a strong enhancement of the absorption peak near the Fermi level due to the resonant scattering of the photoexcited electron by the exciton level . Since the PP signal is linear in the probe field, the essential physics can be captured by assuming a $`\delta `$–function probe pulse, $`_\tau (t)=_\tau e^{i\omega _p\tau }\delta (t\tau )`$, and a Gaussian pump pulse. The PP polarization has the form ($`t>\tau `$)
$$P(t)=i_se^{i\omega _pt}\underset{ij}{}\mu _i\mu _j0|\stackrel{~}{U}_i(t)𝒦(t,\tau )\stackrel{~}{U}_j^{}(\tau )|0,$$
(46)
where $`𝒦(t,\tau )`$ is the time–evolution operator for the effective Hamiltonian,
$$\stackrel{~}{H}(t)=\underset{ij𝐤}{}ϵ_{ij𝐤}^c(t)a_{i𝐤}^{}a_{j𝐤}+\underset{𝐤}{}ϵ_𝐤^v(t)b_𝐤^{}b_𝐤+V_{eh}(t)+V_{ee}(t),$$
(47)
where $`V_{eh}`$ and $`V_{ee}`$ are the effective e-h and e-e interactions and $`\stackrel{~}{U}_i^{}(t)`$ is the effective transition operator given below. Here
$`ϵ_{ij𝐤}^c(t)=\delta _{ij}ϵ_{i𝐤}^c+\mathrm{\Delta }ϵ_{ij𝐤}^c(t),ϵ_𝐤^v(t)=ϵ_𝐤^v+\mathrm{\Omega }+\mathrm{\Delta }ϵ_𝐤^v(t)`$ (48)
are the band dispersions with pump–induced self–energies (to lowest order in the interactions)
$`\mathrm{\Delta }ϵ_{ij𝐤}^c(t)=`$ $`_p(t)[\mu _ip_{j𝐤}^{}(t)+\mu _jp_{i𝐤}(t)]/2,`$
$`\mathrm{\Delta }ϵ_𝐤^v(t)=`$ $`_p(t)\mathrm{Re}{\displaystyle \underset{i}{}}\mu _ip_{i𝐤}(t),`$ (49)
with $`p_{i𝐤}(t)`$ satisfying
$$i\frac{p_{i𝐤}(t)}{t}=(ϵ_{i𝐤}^c+ϵ_𝐤^v+\mathrm{\Omega })p_{i𝐤}(t)\underset{j𝐪}{}v_{ij}p_{j𝐪}(t)+\mu _i(t).$$
(50)
Note that the pump induces additional intersubband scattering, described by $`\mathrm{\Delta }ϵ_{12𝐤}^c(t)`$. To lowest order in the interactions, the effective transition operator appearing in Eq. (46) is given by
$$\stackrel{~}{U}_i^{}(t)=\underset{j𝐤}{}\varphi _{ij𝐤}(t)a_{j𝐤}^{}b_𝐤^{},$$
(51)
with
$$\varphi _{ij𝐤}(t)=\delta _{ij}\left[1\frac{1}{2}\underset{l}{}|p_{l𝐤}(t)|^2\right]\frac{1}{2}p_{i𝐤}(t)p_{j𝐤}^{}(t).$$
(52)
In the one–subband case and in the coherent limit, Eq. (52) takes a familiar form, $`\varphi _𝐤(t)=1|p_𝐤(t)|^2`$ — the usual Pauli blocking factor ; in a multi–subband case, the latter is a matrix.
Similar to the one–subband case, the state $`𝒦(t,\tau )\stackrel{~}{U}_i^{}(\tau )|0`$, entering into (46), can be viewed as describing the propagation of the e–h pair (with wavefunction $`\mathrm{\Phi }_{ij}(𝐤,t)`$) excited by the probe pulse at time $`\tau `$, dressed by the scattering of the FS excitations (dynamical FS response). The latter leads to a dynamical broadening described by the amplitude $`s_{ij}(𝐩,𝐤,t)`$ that satisfies the differential equation
$`i{\displaystyle \frac{s_{ij}(𝐩,𝐤,t)}{t}}=`$ $`(ϵ_{i𝐩}^cϵ_{j𝐤}^c)s_{ij}(𝐩,𝐤,t)`$ (53)
$`+{\displaystyle \underset{l}{}}[\mathrm{\Delta }ϵ_{il𝐩}^c(t)s_{lj}(𝐩,𝐤,t)\mathrm{\Delta }ϵ_{lj𝐤}^c(t)s_{il}(𝐩,𝐤,t)]`$
$`{\displaystyle \underset{l}{}}\stackrel{~}{v}_{il}(𝐩,t)[\delta _{lj}+{\displaystyle \underset{q>k_F}{}}s_{lj}(𝐪,𝐤,t)],`$
with initial condition $`s_{ij}(𝐩,𝐤,\tau )=0`$, and p and k labeling respectively the ($`i`$th subband) FS electron and the ($`j`$th subband) FS hole. Since only the first subband is occupied, the only non–zero components of $`s_{ij}`$ are $`s_{11}(𝐩,𝐤,t)`$ and $`s_{21}(𝐩,𝐤,t)`$, which describe the intra and intersubband FS excitations respectively. The photoexcited e–h pair wavefunction $`\mathrm{\Phi }_{ij}(𝐤,t,\tau )`$ satisfies the Wannier–like equation
$`i{\displaystyle \frac{\mathrm{\Phi }_{ij}(𝐤,t)}{t}}=`$ $`{\displaystyle \underset{l}{}}[ϵ_{il𝐤}^c(t)+\delta _{lj}[ϵ_𝐤^v(t)+ϵ_A(t)]i\mathrm{\Gamma }]\mathrm{\Phi }_{lj}(𝐤,t)`$ (54)
$`{\displaystyle \underset{l,q>k_F}{}}\stackrel{~}{v}_{il}(𝐤,t)\mathrm{\Phi }_{lj}(𝐪,t)`$
with initial condition $`\mathrm{\Phi }_{ij}(𝐤,\tau )=\varphi _{ij𝐤}(\tau )`$, where
$$ϵ_A(t)=\underset{k^{}<k_F}{}[v_{11}+\underset{p^{}>k_F}{}s_{11}(𝐩^{},𝐤^{},t)v_{11}]$$
(55)
is the self–energy due to the readjustment of the FS to the photoexcitation of a hole and $`\mathrm{\Gamma }`$ is the inverse dephasing time due to all the processes not included in $`H`$. In Eqs. (53) and (54),
$$\stackrel{~}{v}_{ij}(𝐤,t)=v_{ij}\underset{l,k^{}<k_F}{}s_{il}(𝐤,𝐤^{},t)v_{lj}$$
(56)
is the effective e–h potential whose time–dependence is due to the dynamical FS response . Note that it is the interplay between this effective potential and the pump-induced self-energies that gives rise to the unadiabatic FS response to the pump field. In terms of $`\mathrm{\Phi }_{ij}(𝐤,t)`$, the polarization (46) takes the simple form ($`t>\tau `$)
$$P(t)=i_se^{i\omega _pt}\underset{ijl}{}\mu _i\mu _j\underset{k>k_F}{}\mathrm{\Phi }_{il}(𝐤,t)\varphi _{ij𝐤}^{}(t),$$
(57)
with $`\varphi _{ij}(𝐤,t)`$ given by (52). The nonlinear absorption spectrum is then proportional to $`\mathrm{Im}P(\omega )`$, where $`P(\omega )`$ is the Fourier transform of the rhs of (57).
## 5 Coherent Ultrafast Dynamics of the Fermi Edge Singularity
### 5.1 Discussion of the Physics: Monochromatic Excitation
In this section, the role of the e-h correlations in the nonlinear optical response of the FES is illustrated in a simple way for monochromatic excitation . In the latter case, the theory discussed in the previous section applies for pump detunings larger than the characteristic Coulomb energy, $`\mathrm{\Omega }>E_M`$. As discussed above, close to the Fermi edge, the linear absorption spectrum of the FES can be approximately described using the simple power law expression in Eq. (4). The monochromatic pump excitation leads to a resonance blueshift, originating from the shift in the effective band energies \[see Fig. 1\], and to a bleaching mainly due to the Pauli blocking (Phase Space Filling) which reduces the effective transition matrix element (analogous to the dressed atom picture ). More importantly, however, the pump–induced change in the band dispersion increases the density of states $`𝒩`$ close to the Fermi surface and thus also increases both the e–h scattering strength $`g=V𝒩`$ and the phaseshift $`\delta \mathrm{tan}^1(\pi g)`$ that determine the FES lineshape (see Eq. (4)). This, in turn, leads to an increase in the FES exponent $`\beta `$ that determines the resonance width. In contrast, in the case of a bound excitonic state of dimensionality $`D`$ and Bohr radius $`a_B`$, the resonance width remains unchanged, while the oscillator strength, $`a_B^D`$, increases by a factor $`(1+D\mathrm{\Delta }m)`$, where $`\mathrm{\Delta }m`$ is the pump–induced change in the effective mass . Such an optically–induced enhancement of the exciton strength competes with the bleaching due to the Pauli blocking and the exciton–exciton interactions. This results in an almost rigid exciton blueshift, consistent with experiment and previous theoretical results .
However, in the case of a FES resonance, the pump–induced change in the exponent $`\delta `$ leads to a stronger oscillator strength enhancement than for a bound exciton state. Obviously, such an enhancement cannot be described perturbatively, i.e., with an expansion in terms of the optical field, since, as can be seen from Eq. (4), the corresponding corrections to the absorption spectrum diverge logarithmically for frequencies $`\omega `$ at the Fermi edge. Eq. (4) shows that the effect of the pump on the FES can be thought of as an excitation–induced dephasing that affects the frequency dependence of the resonance; again, this is in contrast to the case of the exciton. In the time domain, this also implies a memory structure related to the response–time of the FS-excitations. Therefore, the qualitative differences between the nonlinear optical response of the FES and the exciton originate from the fact that an exciton is a discrete bound state, while the FES is a continuum many–body resonance. The FS responds unadiabatically to the pump–induced change in the density of states via an increase in the e–h scattering of low–energy pair excitations. Such scattering processes, which determine the response of the Fermi sea to the hole potential in the course of the optical excitation, are responsible for the unbinding and broadening of the HFA bound state . Therefore, the pump field changes the broadening and dephasing effects even for below–resonant photoexcitation. On the other hand, due to the finite Coulomb binding energy of the exciton, the pump optical field can polarize such a bound state and change its Bohr radius without ionizing it.
### 5.2 Results for Short–pulse Excitation
In this section, we discuss the nonlinear absorption of the FES in the case of short pulse excitation. The results presented here were obtained by solving numerically the differential equations (4.4) and (39), using the Runge–Kutta method, for Gaussian pulses with duration $`t_p=2.0E_F^1`$ . In order to suppress the incoherent effects due to the e–e scattering of real carrier populations with the FS, below–resonant pump excitation and negative time delays were considered. Under such excitation conditions, the coherent effects in which we are interested dominate, and the Coulomb–induced corrections to the effective parameters, discussed in section 4.3, are perturbative. The goal of this section is to study the role of the dynamical FS response (e–h correlations) on the pump–probe dynamics. For this reason, we compare the results of the theory outlined in section 4 to those of the HFA, obtained by setting $`s(𝐩,𝐤,t)=0`$ in Eq. (4.4). As mentioned above, in the latter case, the (spurious) HFA bound state does not interact with the FS pair excitations, even though it can merge with the continuum when one introduces a very short dephasing time.
In Fig. 3, the linear absorption lineshape (i.e., in the absence of pump) of the FES is compared to the HFA (without the dynamical FS response). The parameter values $`g=0.4`$ and $`\mathrm{\Gamma }=0.1E_F`$ were used, which were previously used to fit the experimental spectra in modulation doped quantum wells . For better visibility, we shifted the curves in order to compare their lineshapes. The linear absorption FES lineshape is consistent with that obtained in Ref. . On the other hand, the HFA spectrum is characterized by the coexistence of the bound state and a continuum contribution due to the fact that, in 2D, a bound state exists even for an arbitrary weak attractive potential. We note that if one limits oneself to linear absorption, it is possible to artificially shorten the dephasing time $`T_2=\mathrm{\Gamma }^1`$, mainly determined by the hole recoil effects, by taking $`\mathrm{\Gamma }E_M`$. Then the spurious discrete state and the continuum merge, and the discrepancy between the two linear absorption lineshapes decreases. This trick has been used for phenomenological fits of linear absorption experimental data. However, as we discuss below, in the nonlinear absorption case the differences in the transient spectra are significant so that the processes beyond HFA can be observed experimentally.
Let us turn to the time evolution of the pump–probe spectra. In Fig. 4 we show the nonlinear absorption spectra calculated by including the dynamical FS response \[Fig. 4(a)\] and within the HFA \[Fig. 4(b)\] at a short time delay $`\tau =t_p/2`$. The main features of the spectrum are a pump–induced resonance bleaching, blueshift, and gain right below the onset of absorption. For off–resonant pump, these transient effects vanish for positive time delays after the pump is gone, and persist for negative time delays shorter than the dephasing time $`T_2=\mathrm{\Gamma }^1`$. Similar features were also obtained for different values of the pump amplitude, duration, and detuning. They are mainly due to the broadening induced by the transient renormalization of the energy band dispersion \[Eqs. (35) and (36)\] when its duration $`t_p`$ is shorter than the dephasing time (analogous to excitons and two–level systems).
Let us now turn to the role of the e–h correlations. In Fig. 5 we compare the differential transmission spectrum calculated by including the dynamical FS response or within the HFA for long and short negative time delays. Note that, in PP spectroscopy, the experimentally measured differential transmission is given by Eq. (5) and in the weak signal regime, it reproduces the pump–induced changes in the probe absorption coefficient $`\alpha (\omega ,\tau )`$: $`DST(\omega ,\tau )\mathrm{\Delta }\alpha (\omega ,\tau )`$. Fig. 5(a) shows the results obtained for a long time delay, $`\tau =1.5T_2=15.0E_F^1`$, in which case frequency domain oscillations are observed. These oscillations are similar to those seen in undoped semiconductors and two–level systems ; however, their amplitude in the FES case is reduced. On the other hand, as shown in Fig. 5(b), for time delays comparable to the pulse duration, $`\tau =0.1T_2=t_p/2=1.0E_F^1`$, the main features are a blueshift and bleaching. In this case the e–h correlations lead to a substantially larger width and asymmetric lineshape of the differential transmission spectrum. This comes from the different response of the FES to the pump–induced dispersion renormalizations when the e–h correlations are accounted for. This is more clearly seen in Fig. 6, where the magnitude of the resonance decrease, evaluated at the peak frequency, is plotted as a function of $`\tau `$. Clearly, the bleaching of the FES peak is substantially stronger when the dynamical FS response is included than in the HFA case. Note that for $`|\tau |\mathrm{\Gamma }^1`$ the FES resonance is actually enhanced by the pump, as can be seen more clearly in Fig. 7. The time dependence of the resonance bleaching is strikingly different in the two cases. In the HFA case, the $`|DST(\omega ,\tau )|`$ evaluated at the instantaneous peak frequency decays over a time scale $`|\tau |\mathrm{\Gamma }^1`$, i.e. during the dephasing time. This is similar to results obtained for a two–level system with the same effective parameters. On the other hand, the decay of $`|DST(\omega ,\tau )|`$ at the peak frequency is much faster when we take into account the e–h correlations. Note that the above results were obtained for off–resonant excitation. Under resonant conditions, a spectral hole is produced. In Fig. 8 we compare the resonance blueshifts, evaluated at the peak frequency, as a function of $`\tau `$. Again, a larger blueshift is predicted when the dynamical FS response is included. This suggests that in the experiment of Ref. , where similar blueshifts were observed in two quantum well samples (one MDQW with a FES and one undoped sample with a 2D-exciton) the effective parameters were larger in the latter case, due to the absence of screening and exciton–exciton interaction effects.
In order to gain qualitative understanding of the role of the e–h correlations, let us for a moment neglect the momentum dependence of the pump–renormalization of the band dispersion and the phase space filling effects and consider the bleaching caused by a rigid semiconductor band shift $`\mathrm{\Delta }E_g(t)`$, obtained from the pump–induced self–energies, Eqs. (35) and (36), evaluated at the bottom of the band (note that $`\mathrm{\Delta }E_g(t)`$ lasts for the duration of pulse). Within this model, the pump excitation has no effect on the e–h scattering amplitude $`s(𝐩,𝐤,t)`$ \[see Eq. (39)\]. It is thus convenient to factorize the effects of the rigid band shift on the e–h wavefunction $`\mathrm{\Phi }_𝐩(t,t^{})`$:
$$\mathrm{\Phi }_𝐩(t,t^{})=e^{i_t^{}^t\mathrm{\Delta }E_g(t^{\prime \prime })𝑑t^{\prime \prime }}\stackrel{~}{\mathrm{\Phi }}_𝐩(t,t^{}).$$
(58)
This relation is general and defines $`\stackrel{~}{\mathrm{\Phi }}_𝐩(t,t^{})`$, which does not depend on $`\mathrm{\Delta }E_g(t)`$. In the special case of a rigid shift, $`\stackrel{~}{\mathrm{\Phi }}_𝐩(t,t^{})`$ coincides with $`\mathrm{\Phi }_𝐩^0(tt^{})`$ describing the propagation of the probe–photoexcited e–h pair in the absence of the pump pulse. By substituting into Eq. (58) the long–time asymptotic expression $`\stackrel{~}{\mathrm{\Phi }}_𝐩(t,t^{})=\mathrm{\Phi }_𝐩^0(tt^{})[i(tt^{})E_F]^{\beta 1}`$ that gives the linear absorption spectrum of the FES at $`\omega E_F`$, and substituting the resulting $`\mathrm{\Phi }_𝐩(t,t^{})`$ into Eq. (44), we obtain a simple analytic expression for the effect of a pump–induced rigid band shift on the nonlinear absorption spectrum:
$$\alpha (\omega )\mathrm{Re}_\tau ^{\mathrm{}}𝑑te^{i(\omega +i\mathrm{\Gamma })(t\tau )i_\tau ^t\mathrm{\Delta }E_g(t^{})𝑑t^{}}[i(t\tau )E_F]^{\beta 1}.$$
(59)
For $`\mathrm{\Delta }E_g(t)=0`$ one, of course, recovers the linear FES absorption in the vicinity of the Fermi edge . For $`\beta =0`$, Eq. (59) gives the absorption of the non–interacting continuum.
The physics of the FES can be seen from Eq. (59). For $`\beta =1`$, this gives a discrete Lorentzian peak corresponding to the HFA bound state. However, during the optical transition, the e–h pair interacts with the FS electrons, leading to the readjustment of the FS density profile via the scattering of FS pairs. This results in the broadening of the discrete HFA bound state, which is governed by the time evolution of the FS. Such time evolution is unadiabatic due to the low–energy FS pairs, which leads to the characteristic power–law time dependence of the broadening factor in Eq. (59). The interaction with the FS-pairs determines the exponent, $`0\beta 1`$, of the latter, which leads to a non–Lorentzian lineshape in the frequency domain and a non–exponential decay in the time domain. A detailed discussion of the above physics and the analogy to phonon sidebands and collision broadening may be found in Ref. . However, the CCE must be used in order to calculate the spectrum at all frequencies (and not just asymptotically close to $`E_F`$ as with Eq. (59)) and, most importantly, to describe the non–equilibrium FS and e–h pair response to the time–dependent increase in the effective mass/density of states, not included in Eq. (59).
The resonance bleaching obtained from Eq. (59) as a function of $`\tau `$, is shown in Fig. 6(b) for $`\beta =0.6`$, corresponding to the value of the parameters used in Fig. 6(a) ($`g=0.4`$), together with the HFA result ($`\beta =1`$). Comparing Figs. 6(a) and (b), one can see that the rigid band shift approximation qualitatively accounts for the dynamics, but that there are strong discrepancies (see vertical scales), whose origin is discussed below. Both the magnitude and the time–dependence of the bleaching depends critically on the value of $`\beta `$, which characterizes the interaction of the photoexcited e–h pair with the FS excitations. Because of such coupling, many polarization components are excited in the case of the continuum FES resonance, and it is their interference that governs the dynamics of the PP signal. Such interference is also responsible for the resonance enhancement and differential transmission oscillations at $`\tau <0`$ shown in Figs. 5 and 7. As $`\beta `$ increases the interference effects are suppressed because the energy width of the continuum states contributing to the FES narrows. In fact, this energy width is directly related to that of the linear absorption resonance. This is clearly seen in Fig. 9 where we show the effect of increasing $`g`$ on the dynamics of the bleaching. It becomes more bound–state–like as, with increasing $`g`$, the FES resonance becomes narrower. On the other hand, in the HFA case, the decay rate is $`T_2=\mathrm{\Gamma }^1`$, i. e. it is independent on $`g`$, when $`E_M`$ becomes smaller than $`\mathrm{\Gamma }`$, while for $`E_M\mathrm{\Gamma }`$, the contribution of the continuum states produce a faster decay.
Although the transient rigid band shift approximation, Eq. (59), explains some of the features of the dynamics of the bleaching, it strongly overestimates its magnitude. This is because Eq. (59) neglects the response of the many–body system to the pump–induced renormalization of the band’s dispersion. Such a transient change in the dispersion, which can be viewed as an increase in the density of states/effective mass for the duration of the pump, is important because it results in an enhancement of the e–h scattering. For example, in the case of monochromatic excitation, this leads to the change in the exponent $`\beta `$ of the broadening prefactor in the integrand of Eq. (59), as discussed Section 5.1. For the short–pulse case, it is not possible to describe analytically the effect of the pump on the e–h scattering processes, due to the non–equilibrium unadiabatic FS response. The latter can however be described in a simple way with the numerical solution of Eqs. (39) and (4.4), which is presented in Figs. 10 and 11 and discussed below.
In order to show the role of the pump–induced renormalization of the band dispersion in the presense of the dynamical FS response, we plot in Fig. 10 the function
$$F(\omega ,\tau )=\mathrm{Im}\underset{p>k_F}{}\stackrel{~}{\mathrm{\Phi }}_𝐩(\omega ,\tau )$$
(60)
where $`\stackrel{~}{\mathrm{\Phi }}_𝐩(\omega ,\tau )`$ is the Fourier transform of $`\stackrel{~}{\mathrm{\Phi }}_𝐩(t,\tau )`$ defined by Eq. (58). Note that, in the presence of the band dispersion renormalization, the wave–function $`\stackrel{~}{\mathrm{\Phi }}_𝐩`$ (which is independent of $`\mathrm{\Delta }E_g(t)`$) no longer coincides with $`\mathrm{\Phi }_𝐩^0`$ as in Eq. (59). As can be seen in Fig. 10(a), when the e–h correlations are taken into account, the pump–induced redistribution of oscillator strength between the states of the continuum that contribute to the resonance manifests itself as a dynamical redshift. This shift opposes the rigid band blueshift $`\mathrm{\Delta }E_g(t)`$ (when the latter is included). At the same time, the resonance strength is enhanced significantly. The latter effect originates from the interplay between the transient increase in the effective mass/density of states of the photoexcited e–h pair and the “dressing” of this pair with the FS excitations \[described by the effective potential $`\stackrel{~}{V}(𝐩,t)`$ in Eq. (4.4)\]. In contrast, such an oscillator strength enhancement is suppressed in the HFA (which neglects the e–h correlations), as seen in Fig. 10(b), in which case the main feature is the redshift of the resonance due to the pump–induced increase of the binding energy $`E_M`$ coming from the transient increase in the effective mass.
In Fig. 11 the effect of the renormalization of the band dispersion on the nonlinear absorption spectrum is shown. The optically–induced increase in the e–h interactions enhances significantly the strength of the FES and compensates part of the bleaching induced by the rigid band shift. A smaller enhancement is also seen in the HFA, where the pump–induced increase in the binding energy $`E_M`$ competes with the effects of the bandgap renormalization.
### 5.3 Ultrafast Dynamics of the FES–Exciton Hybrid
In this section we present the results for the evolution of the PP spectra of the FES–exciton hybrid . The spectra in Ref. were obtained by the numerical solution of the coupled equations (54) and (53), with the time–dependent band dispersions $`ϵ_{ij𝐤}^c(t)`$ and $`ϵ_𝐤^v(t)`$. In typical MDQW’s, the intersubband scattering strength characterized by the screened potential $`v_{12}`$ is much smaller than the intrasubband scattering strength characterized by $`v_{ii}`$ (a value $`v_{12}/v_{11}0.2`$ was deduced from the fit to the linear absorption spectrum in ). As discussed in the previous section, in the absence of coupling ($`v_{12}=0`$), the different nature of the exciton and FES leads to distinct dynamics under ultrafast excitation. In the presence of the coupling, one should expect new effects coming from the interplay of this difference and the intersubband scattering that hybridizes the two resonances. Indeed, it was demonstrated in Ref. that, at negative time delays, the PP spectrum undergoes a drastic transformation due to a transient light–induced redistribution of the oscillator strength between the FES and the exciton. We discuss this at length in this section and show that such a redistribution is a result of the dynamical FS response to the pump pulse. In fact, the ultrafast PP spectra of the FES–exciton hybrid can serve as an experimental test of the difference between the FES and exciton dynamics.
The calculations in Ref. were performed at zero temperature for below–resonant pump with detuning $`\mathrm{\Omega }E_F`$ and duration $`t_pE_F/\mathrm{}=2.0`$, and by adopting the typical values of parameters $`v_{12}/v_{11}=0.2`$, $`\mathrm{\Gamma }=0.1E_F`$, and $`v_{11}𝒩=0.3`$, $`𝒩`$ being the density of states, previously extracted from fits to the linear absorption spectra ($`E_F1520`$ meV in typical GaAs/GaAlAs QW’s ). Note, however, that similar results were also obtained for a broad range of parameter values. In Fig. 12(a) we plot the nonlinear absorption spectra at different negative time delays $`\tau <0`$. For better visibility, the curves are shifted vertically with decreasing $`|\tau |`$ (the highest curve represents the linear absorption spectrum). For the chosen value of the subband separation $`\mathrm{\Delta }`$, the FES and excitonic components of the hybrid are distinguishable in the linear absorption spectrum, with the FES peak carrying larger oscillator strength. It can be seen that, at short $`\tau <0`$, the oscillator strength is first transferred to the exciton and then, with further increase in $`|\tau |`$, back to the FES. At the same time, both peaks experience a blueshift, which is larger for the FES than for the exciton peak because the ac–Stark effect for the exciton is weaker due to the subband separation $`\mathrm{\Delta }`$ and the correlation effects.
The transient exchange of oscillator strength originates from the different nature of the FES and exciton components of the hybrid. At negative time delays, the time–evolution of the exciton is governed by its dephasing time, which is essentially determined by the homogeneous broadening $`\mathrm{\Gamma }`$ (in doped systems the exciton–exciton correlations do not play a significant role due to the screening). The pump pulse first leads to a bleaching of the exciton peak, which then recovers its strength at $`|\tau |\mathrm{}/\mathrm{\Gamma }`$. On the other hand, since the FES is a many–body continuum resonance, (i) the bleaching of the FES peak is stronger, and (ii) the polarization decay of the FES is determined not by $`\mathrm{\Gamma }`$, but by the scattering with the low–lying FS excitations. This leads to much faster dynamics, roughly determined by the inverse Coulomb energy $`E_M`$ (see discussion in the previous section). However, the time–evolution of the hybrid spectrum is not a simple superposition of the dynamics of its components. Indeed, the pump-induced self-energies lead to the flattening of the subbands or, to the first approximation, to a time–dependent increase in the effective mass (and hence the density of states), which in turn increases the e–h scattering . Important is, however, that, due to the subband separation and different nature of the resonances, such an increase is stronger for the FES. Therefore, the effect of the pump is to reduce the excitonic enhancement of the FES peak (coming from the resonant scattering of the photoexcited electron by the exciton level) as compared to the linear absorption case, resulting in the oscillator strength transfer from the FES back to exciton. In fact, such a transfer is strong even for smaller $`\mathrm{\Delta }`$ \[see Fig. 12(b)\]. It should be emphasized that the above feature cannot be captured within the HFA. Indeed, the latter approximates the FES by a bound state and thus neglects the difference between the FES and exciton dynamics originating from the unadiabatic response of the FS to the change in the e–h correlations. This is demonstrated in Fig. 12(c) where we show the spectra obtained without the FS dynamical response, i.e., by setting $`s_{ij}=0`$. Although in that case both peaks show blue shift and broadening, there is no significant transfer of oscillator strength
This concludes our discussion of the coherent nonlinear response of the FES. In the rest of this article, we will review the role of size–dependent correlation effects, due to quasiparticle scattering via surface collective modes, on the ultrafast dynamics of the SP resonance in small metal nanoparticles.
## 6 Quasiparticle Scattering with Surface Collective Modes in Metal Nanoparticles
### 6.1 Electron–Electron Interactions in Metal Nanoparticles
Before we proceed with the ultrafast dynamics, we discuss in this section the effect of the surface collective excitations on the e–e interactions in a spherical metal particle. In particular, we present a detailed derivation of the dynamically screened Coulomb potential by generalizing a method previously developed for calculations of local field corrections to the optical fields .
The potential $`U(\omega ;𝐫,𝐫^{})`$ at point $`𝐫`$ arising from an electron at point $`𝐫^{}`$ is determined by the equation
$`U(\omega ;𝐫,𝐫^{})=u(𝐫𝐫^{})+{\displaystyle 𝑑𝐫_1𝑑𝐫_2u(𝐫𝐫_1)\mathrm{\Pi }(\omega ;𝐫_1,𝐫_2)U(\omega ;𝐫_2,𝐫^{})},`$ (61)
where $`u(𝐫𝐫^{})=e^2|𝐫𝐫^{}|^1`$ is the unscreened Coulomb potential and $`\mathrm{\Pi }(\omega ;𝐫_1,𝐫_2)`$ is the polarization operator. There are three contributions to $`\mathrm{\Pi }`$, arising from the polarization of the conduction electrons, the $`d`$–electrons, and the medium surrounding the nanoparticles: $`\mathrm{\Pi }=\mathrm{\Pi }_c+\mathrm{\Pi }_d+\mathrm{\Pi }_m`$. It is useful to rewrite Eq. (61) in the “classical” form
$$(𝐄+4\pi 𝐏)=4\pi e^2\delta (𝐫𝐫^{}),$$
(62)
where $`𝐄(\omega ;𝐫,𝐫^{})=U(\omega ;𝐫,𝐫^{})`$ is the screened Coulomb field and $`𝐏=𝐏_c+𝐏_d+𝐏_m`$ is the electric polarization vector, related to the potential $`U`$ as
$$𝐏(\omega ;𝐫,𝐫^{})=e^2𝑑𝐫_1\mathrm{\Pi }(\omega ;𝐫,𝐫_1)U(\omega ;𝐫_1,𝐫^{}).$$
(63)
In the random phase approximation (RPA), the intraband polarization operator is given by
$`\mathrm{\Pi }_c(\omega ;𝐫,𝐫^{})={\displaystyle \underset{\alpha \alpha ^{}}{}}{\displaystyle \frac{f(E_\alpha ^c)f(E_\alpha ^{}^c)}{E_\alpha ^cE_\alpha ^{}^c+\omega +i0}}\psi _\alpha ^c(𝐫)\psi _\alpha ^{}^c(𝐫)\psi _\alpha ^c(𝐫^{})\psi _\alpha ^{}^c(𝐫^{}),`$ (64)
where $`E_\alpha ^c`$ and $`\psi _\alpha ^c`$ are the single–electron eigenenergies and eigenfunctions in the nanoparticle, and $`f(E)`$ is the Fermi–Dirac distribution (we set $`\mathrm{}=1`$). Since we are interested in frequencies much larger than the single–particle level spacing, $`\mathrm{\Pi }_c(\omega )`$ can be expanded in terms of $`1/\omega `$. For the real part, $`\mathrm{\Pi }_c^{}(\omega )`$, we obtain in the leading order
$$\mathrm{\Pi }_c^{}(\omega ;𝐫,𝐫_1)=\frac{1}{m\omega ^2}[n_c(𝐫)\delta (𝐫𝐫_1)],$$
(65)
where $`n_c(𝐫)`$ is the conduction electron density. In the following we assume, for simplicity, a step density profile, $`n_c(𝐫)=\overline{n}_c\theta (Rr)`$, where $`\overline{n}_c`$ is the average density. The leading contribution to the imaginary part, $`\mathrm{\Pi }_c^{\prime \prime }(\omega )`$, is proportional to $`\omega ^3`$, so that $`\mathrm{\Pi }_c^{\prime \prime }(\omega )\mathrm{\Pi }_c^{}(\omega )`$.
By using Eqs. (65) and (63), we obtain a familiar expression for $`𝐏_c`$ at high frequencies,
$$𝐏_c(\omega ;𝐫,𝐫^{})=\frac{e^2n_c(𝐫)}{m\omega ^2}U(\omega ;𝐫,𝐫^{})=\theta (Rr)\chi _c(\omega )𝐄(\omega ;𝐫,𝐫^{}),$$
(66)
where $`\chi _c(\omega )=e^2\overline{n}_c/m\omega ^2`$ is the conduction electron susceptibility. Note that, for a step density profile, $`𝐏_c`$ vanishes outside the particle. The $`d`$–band and dielectric medium contributions to $`𝐏`$ are also given by similar relations,
$`𝐏_d(\omega ;𝐫,𝐫^{})=\theta (Rr)\chi _d(\omega )𝐄(\omega ;𝐫,𝐫^{}),`$ (67)
$`𝐏_m(\omega ;𝐫,𝐫^{})=\theta (rR)\chi _m𝐄(\omega ;𝐫,𝐫^{}),`$ (68)
where $`\chi _i=(ϵ_i1)/4\pi `$, $`i=d,m`$ are the corresponding susceptibilities and the step functions account for the boundary conditions . Using Eqs. (66)–(68), one can write a closed–form equation for $`U(\omega ;𝐫,𝐫^{})`$. Using Eq. (63), the second term of Eq. (61) can be presented as $`e^2𝑑𝐫_1u(𝐫𝐫_1)𝐏(\omega ;𝐫_1,𝐫^{}).`$ Substituting the above expressions for $`𝐏`$, we then obtain after integration by parts
$`ϵ(\omega )U(\omega ;𝐫,𝐫^{})=`$ $`{\displaystyle \frac{e^2}{|𝐫𝐫^{}|}}+{\displaystyle }d𝐫_1_1{\displaystyle \frac{1}{|𝐫𝐫_1|}}_1[\theta (Rr)\chi (\omega )`$ (69)
$`+\theta (rR)\chi _m]U(\omega ;𝐫_1,𝐫^{})`$
$`+i{\displaystyle 𝑑𝐫_1𝑑𝐫_2\frac{e^2}{|𝐫𝐫_1|}\mathrm{\Pi }_c^{\prime \prime }(\omega ;𝐫_1,𝐫_2)U(\omega ;𝐫_2,𝐫^{})},`$
with
$$ϵ(\omega )1+4\pi \chi (\omega )=ϵ_d(\omega )\omega _p^2/\omega ^2,$$
(70)
$`\omega _p^2=4\pi e^2\overline{n}_c/m`$ being the plasmon frequency in the conduction band. The last term in the rhs of Eq. (69), proportional to $`\mathrm{\Pi }_c^{\prime \prime }(\omega )`$, can be regarded as a small correction. To solve Eq. (69), we first eliminate the angular dependence by expanding $`U(\omega ;𝐫,𝐫^{})`$ in spherical harmonics, $`Y_{LM}(\widehat{𝐫})`$, with coefficients $`U_{LM}(\omega ;r,r^{})`$. Using the corresponding expansion of $`|𝐫𝐫^{}|^1`$ with coefficients $`Q_{LM}(r,r^{})=\frac{4\pi }{2L+1}r^{L1}r^L`$ (for $`r>r^{}`$), we get the following equation for $`U_{LM}(\omega ;r,r^{})`$:
$`ϵ(\omega )U_{LM}(\omega ;r,r^{})=`$ $`Q_{LM}(r,r^{})+4\pi \left[\chi (\omega )\chi _m\right]{\displaystyle \frac{L+1}{2L+1}}\left({\displaystyle \frac{r}{R}}\right)^LU_{LM}(\omega ;R,r^{})`$ (71)
$`+ie^2{\displaystyle \underset{L^{}M^{}}{}}{\displaystyle 𝑑r_1𝑑r_2r_1^2r_2^2Q_{LM}(r,r_1)\mathrm{\Pi }_{LM,L^{}M^{}}^{\prime \prime }(\omega ;r_1,r_2)}`$
$`\times U_{L^{}M^{}}(\omega ;r_2,r^{}),`$
where
$`\mathrm{\Pi }_{LM,L^{}M^{}}^{\prime \prime }(\omega ;r_1,r_2)={\displaystyle 𝑑\widehat{𝐫}_1𝑑\widehat{𝐫}_2Y_{LM}^{}(\widehat{𝐫}_1)\mathrm{\Pi }_c^{\prime \prime }(\omega ;𝐫_1,𝐫_2)Y_{L^{}M^{}}(\widehat{𝐫}_2)},`$ (72)
are the coefficients of the multipole expansion of $`\mathrm{\Pi }_c^{\prime \prime }(\omega ;𝐫_1,𝐫_2)`$. For $`\mathrm{\Pi }_c^{\prime \prime }=0`$, the solution of Eq. (71) can be presented in the form
$`U_{LM}(\omega ;r,r^{})=a(\omega )e^2Q_{LM}(r,r^{})+b(\omega ){\displaystyle \frac{4\pi e^2}{2L+1}}{\displaystyle \frac{r^Lr^L}{R^{2L+1}}},`$ (73)
with frequency–dependent coefficients $`a`$ and $`b`$. Since $`\mathrm{\Pi }_c^{\prime \prime }(\omega )\mathrm{\Pi }_c^{}(\omega )`$ for the relevant frequencies, the solution of Eq. (71) in the presence of the last term can be written in the same form as Eq. (73), but with modified $`a(\omega )`$ and $`b(\omega )`$. Substituting Eq. (73) into Eq. (71), we obtain after lengthy algebra in the lowest order in $`\mathrm{\Pi }_c^{\prime \prime }`$
$$a(\omega )=ϵ^1(\omega ),b(\omega )=ϵ_L^1(\omega )ϵ^1(\omega ),$$
(74)
where
$$ϵ_L(\omega )=\frac{L}{2L+1}ϵ(\omega )+\frac{L+1}{2L+1}ϵ_m+iϵ_{cL}^{\prime \prime }(\omega ),$$
(75)
is the effective dielectric function, whose zero, $`ϵ_L^{}(\omega _L)=0`$, determines the frequency of the collective surface excitation with angular momentum $`L`$ ,
$$\omega _L^2=\frac{L\omega _p^2}{Lϵ_d^{}(\omega _L)+(L+1)ϵ_m}.$$
(76)
In Eq. (75), $`ϵ_{cL}^{\prime \prime }(\omega )`$ characterizes the damping of the $`L`$–pole collective mode by single–particle excitations, and is given by
$$ϵ_{cL}^{\prime \prime }(\omega )=\frac{4\pi ^2e^2}{(2L+1)R^{2L+1}}\underset{\alpha \alpha ^{}}{}|M_{\alpha \alpha ^{}}^{LM}|^2[f(E_\alpha ^c)f(E_\alpha ^{}^c)]\delta (E_\alpha ^cE_\alpha ^{}^c+\omega ),$$
(77)
where $`M_{\alpha \alpha ^{}}^{LM}`$ are the matrix elements of $`r^LY_{LM}(\widehat{𝐫})`$. Due to the momentum non–conservation in a nanoparticle, the matrix elements are finite, which leads to the size–dependent width of the $`L`$–pole mode :
$$\gamma _L=\frac{2L+1}{L}\frac{\omega ^3}{\omega _p^2}ϵ_{cL}^{\prime \prime }(\omega ).$$
(78)
For $`\omega \omega _L`$, one can show that the width, $`\gamma _Lv_F/R`$, is independent of $`\omega `$. Note that, in noble metal particles, there is an additional d–electron contribution to the imaginary part of $`ϵ_L(\omega )`$ at frequencies above the onset $`\mathrm{\Delta }`$ of the interband transitions.
Putting everything together, we arrive at the following expression for the dynamically–screened interaction potential in a nanoparticle:
$`U(\omega ;𝐫,𝐫^{})={\displaystyle \frac{u(𝐫𝐫^{})}{ϵ(\omega )}}+{\displaystyle \frac{e^2}{R}}{\displaystyle \underset{LM}{}}{\displaystyle \frac{4\pi }{2L+1}}{\displaystyle \frac{1}{\stackrel{~}{ϵ}_L(\omega )}}\left({\displaystyle \frac{rr^{}}{R^2}}\right)^LY_{LM}(\widehat{𝐫})Y_{LM}^{}(\widehat{𝐫}^{}),`$ (79)
with $`\stackrel{~}{ϵ}_L^1(\omega )=ϵ_L^1(\omega )ϵ^1(\omega )`$. Equation (79) represents a generalization of the plasmon pole approximation to spherical metal particles. The two contributions to the rhs originate from two types of dynamical screening. The first describes the usual bulk-like screening of the Coulomb potential by the electrons inside the particle. The second contribution is a new effective interaction induced by the surface: the potential of an electron inside the nanoparticle excites high–frequency surface collective modes, which in turn act as image charges that interact with the second electron. It should be emphasized that, unlike in the case of the optical fields, the surface–induced dynamical screening of the Coulomb potential is size–dependent.
Note that the excitation energies of the surface collective modes are lower than the bulk plasmon energy, also given by Eq. (76) but with $`ϵ_m=0`$. As discussed below, this opens up new channels of quasiparticle scattering.
### 6.2 Conduction Electron Scattering
In this section, we present the calculation of the rates of electron scattering in the conduction band accompanied by the emission of surface collective modes and discuss its possible experimental manifestations. In the first order in the surface–induced potential, given by the second term in the rhs of Eq. (79), the corresponding scattering rate can be obtained from the Matsubara self–energy
$`\mathrm{\Sigma }_\alpha ^c(i\omega )={\displaystyle \frac{1}{\beta }}{\displaystyle \underset{i\omega ^{}}{}}{\displaystyle \underset{LM}{}}{\displaystyle \underset{\alpha ^{}}{}}{\displaystyle \frac{4\pi e^2}{(2L+1)R^{2L+1}}}{\displaystyle \frac{|M_{\alpha \alpha ^{}}^{LM}|^2}{\stackrel{~}{ϵ}_L(i\omega ^{})}}G_\alpha ^{}^c(i\omega ^{}+i\omega ),`$ (80)
where $`G_\alpha ^c=(i\omega E_\alpha ^c)^1`$ is the non-interacting Green function of the conduction electron. Here the matrix elements $`M_{\alpha \alpha ^{}}^{LM}`$ are calculated with the one–electron wave functions $`\psi _\alpha ^c(𝐫)=R_{nl}(r)Y_{lm}(\widehat{𝐫})`$. Since $`|\alpha `$ and $`|\alpha ^{}`$ are the initial and final states of the scattered electron, the main contribution to the $`L`$th term of the angular momentum sum in Eq. (80) will come from electron states with energy difference $`E_\alpha E_\alpha ^{}\omega _L`$. Therefore, $`M_{\alpha \alpha ^{}}^{LM}`$ can be expanded in terms of the small parameter $`E_0/|E_\alpha ^cE_\alpha ^{}^c|E_0/\omega _L`$, where $`E_0=(2mR^2)^1`$ is the characteristic confinement energy. The leading term can be obtained by using the following procedure . We present $`M_{\alpha \alpha ^{}}^{LM}`$ as
$$M_{\alpha \alpha ^{}}^{LM}=c,\alpha |r^LY_{LM}(\widehat{𝐫})|c,\alpha ^{}=\frac{c,\alpha |[H,[H,r^LY_{LM}(\widehat{𝐫})]]|c,\alpha ^{}}{(E_\alpha ^cE_\alpha ^{}^c)^2},$$
(81)
where $`H=H_0+V(r)`$ is the Hamiltonian of an electron in a nanoparticle with confining potential $`V(r)=V_0\theta (rR)`$. Since $`[H,r^LY_{LM}(\widehat{𝐫})]=\frac{1}{m}[r^LY_{LM}(\widehat{𝐫})]`$, the numerator in Eq. (81) contains a term proportional to the gradient of the confining potential, which peaks sharply at the surface. The corresponding contribution to the matrix element describes the surface scattering of an electron making the $`L`$–pole transition between the states $`|c,\alpha `$ and $`|c,\alpha ^{}`$, and gives the dominant term of the expansion. Thus, in the leading order in $`|E_\alpha ^cE_\alpha ^{}^c|^1`$, we obtain
$`M_{\alpha \alpha ^{}}^{LM}=`$ $`{\displaystyle \frac{c,\alpha |[r^LY_{LM}(\widehat{𝐫})]V(r)|c,\alpha ^{}}{m(E_\alpha ^cE_\alpha ^{}^c)^2}}`$ (82)
$`={\displaystyle \frac{LR^{L+1}}{m(E_\alpha ^cE_\alpha ^{}^c)^2}}V_0R_{nl}(R)R_{n^{}l^{}}(R)\phi _{lm,l^{}m^{}}^{LM},`$
with $`\phi _{lm,l^{}m^{}}^{LM}=𝑑\widehat{𝐫}Y_{lm}^{}(\widehat{𝐫})Y_{LM}(\widehat{𝐫})Y_{l^{}m^{}}(\widehat{𝐫})`$. Note that, for $`L=1`$, Eq. (82) becomes exact. For electron energies close to the Fermi level, $`E_{nl}^cE_F`$, the radial quantum numbers are large, and the product $`V_0R_{nl}(R)R_{n^{}l^{}}(R)`$ can be evaluated by using semiclassical wave–functions. In the limit $`V_0\mathrm{}`$, this product is given by $`2\sqrt{E_{nl}^cE_{n^{}l^{}}^c}/R^3`$, where $`E_{nl}^c=\pi ^2(n+l/2)^2E_0`$ is the electron eigenenergy for large $`n`$. Substituting this expression into Eq. (82) and then into Eq. (80), we obtain
$`\mathrm{\Sigma }_\alpha ^c(i\omega )=`$ $`{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{i\omega ^{}}{}}{\displaystyle \underset{L}{}}{\displaystyle \underset{n^{}l^{}}{}}C_{ll^{}}^L{\displaystyle \frac{4\pi e^2}{(2L+1)R}}{\displaystyle \frac{E_{nl}^cE_{n^{}l^{}}^c}{(E_{nl}^cE_{n^{}l^{}}^c)^4}}`$ (83)
$`\times {\displaystyle \frac{(4LE_0)^2}{\stackrel{~}{ϵ}_L(i\omega ^{})}}G_\alpha ^{}^c(i\omega ^{}+i\omega ),`$
with
$$C_{ll^{}}^L=\underset{M,m^{}}{}|\phi _{lm,l^{}m^{}}^{LM}|^2=\frac{(2L+1)(2l^{}+1)}{8\pi }_1^1𝑑xP_l(x)P_L(x)P_l^{}(x),$$
(84)
where $`P_l(x)`$ are Legendre polynomials; we used properties of the spherical harmonics in the derivation of Eq. (84). For $`E_{nl}^cE_F`$, the typical angular momenta are large, $`lk__FR1`$, and one can use the large–$`l`$ asymptotics of $`P_l`$; for the low multipoles of interest, $`Ll`$, the integral in Eq. (84) can be approximated by $`\frac{2}{2l^{}+1}\delta _{ll^{}}`$. After performing the Matsubara summation, we obtain for the imaginary part of the self–energy that determines the electron scattering rate
$`\text{Im}\mathrm{\Sigma }_\alpha ^c(\omega )={\displaystyle \frac{16e^2}{R}}E_0^2{\displaystyle \underset{L}{}}{\displaystyle 𝑑Eg_l(E)\frac{L^2EE_\alpha ^c}{(E_\alpha ^cE)^4}\text{Im}\frac{N(E\omega )+f(E)}{\stackrel{~}{ϵ}_L(E\omega )}},`$ (85)
where $`N(E)`$ is the Bose distribution and $`g_l(E)`$ is the density of states of a conduction electron with angular momentum $`l`$,
$$g_l(E)=2\underset{n}{}\delta (E_{nl}^cE)\frac{R}{\pi }\sqrt{\frac{2m}{E}},$$
(86)
where we replaced the sum over $`n`$ by an integral (the factor of 2 accounts for spin).
Each term in the sum in the rhs of Eq. (85) represents a channel of electron scattering mediated by a collective surface mode with angular momentum $`L`$. For low $`L`$, the difference between the energies of modes with successive values of $`L`$ is larger than their widths, so the different channels are well separated. Note that since all $`\omega _L`$ are smaller than the frequency of the bulk plasmon, one can replace $`\stackrel{~}{ϵ}_L(\omega )`$ by $`ϵ_L(\omega )`$ in the integrand of Eq. (85) for frequencies $`\omega \omega _L`$.
Consider now the $`L=1`$ term in Eq. (85), which describes the SP–mediated scattering channel. The main contribution to the integral comes from the SP pole in $`ϵ_1^1(\omega )=3ϵ_s^1(\omega )`$, where $`ϵ_s(\omega )`$ is the same as in Eq. (13). To evaluate the integral in Eq. (85), we can in the first approximation replace $`\text{Im}ϵ_s^1(\omega )`$ by a Lorentzian,
$`\text{Im}{\displaystyle \frac{1}{ϵ_s(\omega )}}=`$ $`{\displaystyle \frac{\gamma _s\omega _p^2/\omega ^3+ϵ_d^{\prime \prime }(\omega )}{[ϵ^{}(\omega )+2ϵ_m]^2+[\gamma _s\omega _p^2/\omega ^3+ϵ_d^{\prime \prime }(\omega )]^2}}`$ (87)
$`{\displaystyle \frac{\omega _s^2}{ϵ_d^{}(\omega _s)+2ϵ_m}}{\displaystyle \frac{\omega _s\gamma }{(\omega ^2\omega _s^2)^2+\omega _s^2\gamma ^2}},`$
where $`\omega _s\omega _1=\omega _p/\sqrt{ϵ_d^{}(\omega _s)+2ϵ_m}`$ and $`\gamma =\gamma _s+\omega _sϵ_d^{\prime \prime }(\omega _s)`$ are the SP frequency and width, respectively. For typical widths $`\gamma \omega _s`$, the integral in Eq. (85) can be easily evaluated, yielding
$$\text{Im}\mathrm{\Sigma }_\alpha ^c(\omega )=\frac{24e^2\omega _sE_0^2}{ϵ_d^{}(\omega _s)+2ϵ_m}\frac{E_\alpha ^c\sqrt{2m(\omega \omega _s)}}{(\omega E_\alpha ^c\omega _s)^4}[1f(\omega \omega _s)].$$
(88)
Using the relation $`e^2k_F[ϵ_d^{}(\omega _s)+2ϵ_m]^1=3\pi \omega _s^2/8E_F`$, the SP–mediated scattering rate, $`\gamma _e^s(E_\alpha ^c)=\text{Im}\mathrm{\Sigma }_\alpha ^c(E_\alpha ^c)`$, takes the form
$$\gamma _e^s(E)=9\pi \frac{E_0^2}{\omega _s}\frac{E}{E_F}\left(\frac{E\omega _s}{E_F}\right)^{1/2}[1f(E\omega _s)].$$
(89)
Recalling that $`E_0=(2mR^2)^1`$, we see that the scattering rate of a conduction electron is size–dependent: $`\gamma _e^sR^4`$. At $`E=E_F+\omega _s`$, the scattering rate jumps to the value $`9\pi (1+\omega _s/E_F)E_0^2/\omega _s`$, and then increases with energy as $`E^{3/2}`$ (for $`\omega _sE_F`$). This should be contrasted with the usual (bulk) plasmon–mediated scattering, originating from the first term in Eq. (79), in which case the rate decreases as $`E^{1/2}`$ above the onset .
Note that the total electron scattering rate is the sum, $`\gamma _e+\gamma _e^s`$, of the SP–mediated ($`\gamma _e^s`$) and the bulk–like ($`\gamma _e`$) scattering rates. In order to be observable, the former should exceed the latter. The typical size at which $`\gamma _e^s`$ becomes important can be estimated by equating $`\gamma _e^s`$ and the Fermi liquid e–e scattering rate , $`\gamma _e(E)=\frac{\pi ^2q_{_{TF}}}{16k__F}\frac{(EE_F)^2}{E_F}`$. For energies $`EE_F+\omega _s`$, the two rates become comparable for
$$(k__FR)^212\frac{E_F}{\omega _s}\left(1+\frac{E_F}{\omega _s}\right)^{1/2}\left(\frac{k__F}{\pi q_{_{TF}}}\right)^{1/2}.$$
(90)
In the case of a copper nanoparticle with $`\omega _s2.2`$ eV, we obtain $`k__FR8`$, which corresponds to a radius of $`R3`$ nm. At the same time, in this energy range, the width $`\gamma _e^s`$ exceeds the mean level spacing $`\delta `$, so that the energy spectrum is still continuous. The strong size dependence of $`\gamma _e^s`$ indicates that, although $`\gamma _e^s`$ increases with energy slower than $`\gamma _e`$, the SP–mediated scattering should dominate for nanometer–sized particles. Note also that the size and energy dependences of the scattering in the different channels are similar: the rate of scattering via the $`L`$th channel is given by Eq. (89) with $`\omega _s`$ replaced by $`\omega _L`$, Eq. (76), and the numerical factor 9 replaced by $`3L(2L+1)`$.
Concluding this section, we have shown that the SP–mediated scattering is the dominant scattering mechanism of conduction electrons in nanometer-sized nanoparticles for energies larger than $`\omega _s`$ but smaller than $`\omega _p`$. The scattering rate in the $`L`$th channel, $`\gamma _e^L`$, increases with energy, in sharp contrast with the bulk–plasmon–mediated scattering rate. The total scattering rate as a function of energy represents a series of steps at $`E=\omega _L`$, on top of a smooth energy increase. We expect that this new effect should be observable experimentally by measuring e–e scattering rate in size–selected cluster beams in time–resolved two–photon photoemission spectrum .
### 6.3 $`d`$–Band Hole Scattering
We now turn to the interband processes in noble metal particles and consider the scattering of a $`d`$–hole into the conduction band. The corresponding surface–induced potential, given by the $`L`$th term in Eq. (79), has the form
$`U_{LM}(\omega ;𝐫,𝐫^{})={\displaystyle \frac{4\pi }{2L+1}}{\displaystyle \frac{1}{ϵ_L(\omega )}}{\displaystyle \frac{e^2}{R}}\left({\displaystyle \frac{rr^{}}{R^2}}\right)^LY_{LM}(\widehat{𝐫})Y_{LM}^{}(\widehat{𝐫}^{}),`$ (91)
where $`ϵ_L(\omega )`$ is given by Eq. (75) With this potential, the $`d`$–hole self–energy is given by
$`\mathrm{\Sigma }_\alpha ^d(i\omega )={\displaystyle \frac{e^2}{R^{2L+1}}}{\displaystyle \underset{\alpha ^{}}{}}|\stackrel{~}{M}_{\alpha \alpha ^{}}^{LM}|^2{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{i\omega ^{}}{}}{\displaystyle \frac{G_\alpha ^{}^c(i\omega ^{}+i\omega )}{ϵ_L(i\omega ^{})}},`$ (92)
where $`\stackrel{~}{M}_{\alpha \alpha ^{}}^{LM}=c,\alpha |r^LY_{LM}(\widehat{𝐫})|d,\alpha ^{}`$ is the interband transition matrix element \[compare with Eq. (81)\]. Since the final state energies in the conduction band are high (in the case of interest here, they are close to the Fermi level), the matrix element can be approximated by a bulk–like expression $`\stackrel{~}{M}_{\alpha \alpha ^{}}^{LM}=\delta _{\alpha \alpha ^{}}\stackrel{~}{M}_{\alpha \alpha }^{LM}`$, the corrections due to surface scattering being suppressed by a factor of $`(k__FR)^11`$.
The largest contribution to the self energy Eq. (92) comes from the dipole channel, $`L=1`$, mediated by the SP. In this case, after performing the frequency summation, we obtain for Im$`\mathrm{\Sigma }_\alpha ^d`$
$`\mathrm{Im}\mathrm{\Sigma }_\alpha ^d(\omega )={\displaystyle \frac{9e^2\mu ^2}{m^2(E_\alpha ^{cd})^2R^3}}\text{Im}{\displaystyle \frac{N(E_\alpha ^c\omega )+f(E_\alpha ^c)}{ϵ_s(E_\alpha ^c\omega )}},`$ (93)
where $`E_\alpha ^{cd}=E_\alpha ^cE_\alpha ^d`$ and $`\mu `$ is the interband dipole matrix element . We see that the scattering rate of a $`d`$-hole with energy $`E_\alpha ^d`$, $`\gamma _h^s(E_\alpha ^d)=\text{Im}\mathrm{\Sigma }_\alpha ^d(E_\alpha ^d)`$, has a strong $`R^3`$ dependence on the nanoparticle size, which is, however, different from that of the intraband scattering, Eq. (89).
The most important difference between the interband and the intraband SP–mediated scattering rates lies in their dependence on energy. Since the surface–induced potential, Eq. (91), only allows for vertical (dipole) interband single–particle excitations, the phase space available for the scattering of a $`d`$–hole with energy $`E_\alpha ^d`$ is restricted to a single final state in the conduction band, with energy $`E_\alpha ^c`$. As a result of this restriction, the $`d`$–hole scattering rate, $`\gamma _h^s(E_\alpha ^d)`$, exhibits a peak as the difference between the energies of final and initial states, $`E_\alpha ^{cd}=E_\alpha ^cE_\alpha ^d`$, approaches the SP frequency $`\omega _s`$ \[see Eq. (93)\]. In contrast, the energy dependence of $`\gamma _e^s`$ is smooth due the larger phase space available for scattering within the conduction band. This leads to the additional integral over final state energies in Eq. (85), which smears out the SP resonant enhancement of the intraband scattering.
As we show later, the fact that the scattering rate of a $`d`$–hole is dominated by the SP resonance affects strongly the nonlinear optical dynamics in small nanoparticles. This is the case, in particular, when the SP frequency, $`\omega _s`$, is close to the onset of interband transitions, $`\mathrm{\Delta }`$, as, e.g., in copper and gold nanoparticles . Indeed, if the optical pulse excites an e–h pair with excitation energy $`\omega `$ close to $`\mathrm{\Delta }`$, the $`d`$–hole can subsequently scatter into the conduction band by emitting a SP. According to Eq. (93), for $`\omega \omega _s`$, such a scattering process should be resonantly enhanced. In order to have an observable effect on the absorption spectrum, the scattering rate of the photoexcited $`d`$–hole should be comparable (or larger) than that of the photoexcited electron. Close to $`E_F`$, the electron scattering in the conduction band comes from a two–quasiparticle process; the corresponding rate in noble–metals is estimated as $`\gamma _e10^2`$ eV. If one assumes the bulk value for $`\mu `$ ($`2\mu ^2/m1`$ eV near the L-point ), then $`\gamma _h^s`$ exceeds $`\gamma _e`$ for $`R<2.5`$ nm. In fact, one would expect that, in nanoparticles, $`\mu `$ is larger than in the bulk due to the localization of the conduction electron wave–functions .
## 7 Surface Plasmon Nonlinear Optical dynamics
In this section, we discuss the effect of the SP–mediated interband scattering on the ultrafast optical dynamics in noble metal nanoparticles. We are interested in the situation when the hot electron distribution has already thermalized and the electron gas is cooling to the lattice (stage III). In this case the transient response of a nanoparticle can be described by the time–dependent absorption coefficient $`\alpha (\omega ,t)`$, given by Eq. (12) with time–dependent temperature. In noble–metal particles, the temperature dependence of $`\alpha `$ originates from two different sources. First is the phonon–induced correction to $`\gamma _s`$, which is proportional to the lattice temperature $`T_l(t)`$. As mentioned in section 3, for small nanoparticles this effect is relatively weak. Note that, as discussed e.g. in Ref. , the contribution to $`\gamma _s`$ coming from the e–e interaction , which depends on the electron temperature, plays a minor role when the interband transitions are resonantly excited. Second, near the onset of the interband transitions, $`\mathrm{\Delta }`$, the absorption coefficient depends on the electron temperature $`T(t)`$ via the interband dielectric function $`ϵ_d(\omega )`$ \[see Eqs. (12) and (13)\]. In fact, in copper or gold nanoparticles, $`\omega _s`$ can be tuned close to $`\mathrm{\Delta }`$, so the SP damping by interband e–h excitations leads to an additional broadening of the absorption peak . In this case, it is the temperature dependence of $`ϵ_d(\omega )`$ that dominates the pump–probe dynamics. Below we show that, near the SP resonance, both the temperature and frequency dependence of $`ϵ_d(\omega )=1+4\pi \chi _d(\omega )`$ are strongly affected by the SP–mediated interband scattering.
We start with the RPA expression for the interband susceptibility , $`\chi _d(i\omega )=\stackrel{~}{\chi }_d(i\omega )+\stackrel{~}{\chi }_d(i\omega )`$,
$$\stackrel{~}{\chi }_d(i\omega )=\underset{\alpha }{}\frac{e^2\mu ^2}{m^2(E_\alpha ^{cd})^2}\frac{1}{\beta }\underset{i\omega ^{}}{}G_\alpha ^d(i\omega ^{})G_\alpha ^c(i\omega ^{}+i\omega ),$$
(94)
where $`G_\alpha ^d(i\omega ^{})`$ is the Green function of a $`d`$–electron. With the $`d`$-band fully occupied, the only allowed SP–mediated interband scattering is that of the $`d`$–hole. We assume here, for simplicity, a dispersionless $`d`$–band with energy $`E^d`$. Substituting $`G_\alpha ^d(i\omega ^{})=[i\omega ^{}E^d+E_F\mathrm{\Sigma }_\alpha ^d(i\omega ^{})]^1`$, with $`\mathrm{\Sigma }_\alpha ^d(i\omega )`$ given by Eq. (92), and performing the frequency summation, we obtain
$$\stackrel{~}{\chi }_d(\omega )=\frac{e^2\mu ^2}{m^2}\frac{dE^cg(E^c)}{(E^{cd})^2}\frac{f(E^c)1}{\omega E^{cd}+i\gamma _h^s(\omega ,E^c)},$$
(95)
where $`g(E^c)`$ is the density of states of conduction electrons. The scattering rate of a $`d`$-hole, $`\gamma _h^s(\omega ,E^c)=\mathrm{Im}\mathrm{\Sigma }^d(E^c\omega )`$, is obtained from Eq. (93) with $`E_d=E^c\omega `$:
$$\gamma _h^s(\omega ,E^c)=\frac{9e^2\mu ^2}{m^2(E^{cd})^2R^3}f(E^c)\text{Im}\frac{1}{ϵ_s(\omega )},$$
(96)
where $`N(\omega )`$ is negligible for frequencies $`\omega \omega _sk_BT`$ . The rate $`\gamma _h^s(\omega ,E^c)`$ exhibits a sharp peak as a function of the frequency of the probe optical field. The reason for this is that the scattering rate of a $`d`$–hole with energy $`E`$ depends explicitly on the difference between the final and initial states, $`E^cE`$, as discussed above; therefore, for a $`d`$–hole with energy $`E=E^c\omega `$, the dependence on the final state energy, $`E^c`$, cancels out in $`ϵ_s(E^cE)`$ \[see Eq. (93)\]. In other words, the optically–excited $`d`$–hole scatters resonantly into the conduction band as $`\omega `$ approaches $`\omega _s`$. It is important to note that $`\gamma _h^s(\omega ,E^c)`$ is, in fact, proportional to the absorption coefficient $`\alpha (\omega )`$ \[see Eq. (12)\]. Therefore, $`\alpha `$ and $`\gamma _h^s`$ should be calculated self–consistently from Eqs. (12), (13), (95), and (96).
It should be emphasized that the effect of $`\gamma _h^s`$ on $`ϵ_d^{\prime \prime }(\omega )`$ increases with temperature. Indeed, it can be seen from Eq. (96) that the value of $`\gamma _h^s`$ is appreciable only if $`E^cE_F>k_BT`$. Since the main contribution to $`\stackrel{~}{\chi }_d^{\prime \prime }(\omega )`$ comes from energies $`E^c\omega \mathrm{\Delta }+E_F`$, the effect of $`d`$–hole scattering on the absorption becomes important only for elevated electron temperatures: $`k_BT>\omega _s\mathrm{\Delta }`$. As a result, near the SP resonance, the time evolution of the differential absorption, which is governed by the temperature dependence of $`\alpha `$, becomes strongly size–dependent.
The numerical results discussed below are taken from Refs. . In the experiment of Bigot et. al. , the pump–probe measurements were performed on $`R2.5`$ nm copper nanoparticles. The SP frequency, $`\omega _s2.22`$ eV, was slightly above the onset of the interband transitions, $`\mathrm{\Delta }2.18`$ eV. In order to describe the time–evolution of the differential absorption spectra, one first needs to determine the time–dependence of the electron temperature, $`T(t)`$, due to the relaxation of the electron gas to the lattice. For this, a simple two–temperature model is employed, defined by heat equations for $`T(t)`$ and the lattice temperature $`T_l(t)`$:
$`C(T){\displaystyle \frac{T}{t}}`$ $`=`$ $`G(TT_l),`$
$`C_l{\displaystyle \frac{T_l}{t}}`$ $`=`$ $`G(TT_l),`$ (97)
where $`C(T)=\gamma T`$ and $`C_l`$ are the electron and lattice heat capacities, respectively, and $`G`$ is the e–p coupling . The parameter values used in Refs. were $`G=3.5\times 10^{16}`$ Wm<sup>-3</sup>K<sup>-1</sup>, $`\gamma =70`$ Jm<sup>-3</sup>K<sup>-2</sup>, and $`C_l=3.5`$ Jm<sup>-3</sup>K<sup>-1</sup>, and the initial condition was taken as $`T_0=1000`$ K. The time–dependent absorption coefficient $`\alpha (\omega ,t)`$ was calculated self–consistently; the differential transmission is proportional to $`\alpha _r(\omega )\alpha (\omega ,t)`$, where $`\alpha _r(\omega )`$ was calculated at the room temperature.
In Fig. 13, the calculated nonlinear absorption spectra are shown for various nanoparticle sizes. Fig. 13(a) shows the spectra at several time delays for $`R=5.0`$ nm; for this size, the SP–mediated d–hole scattering has no effect. With decreasing nanoparticle size, the linear absorption spectra are not significantly altered, as can be seen in Figs. 13(b) and (c). However, the change in the nonlinear absorption spectra becomes pronounced at short time delays corresponding to higher temperatures \[see Figs. 13(b) and (c)\]. This effect is more clearly seen in the differential transmission spectra, shown in Fig. 14, which undergo a qualitative transformation with decreasing size.
Note that it is necessary to include the intraband e–e scattering in order to reproduce the differential transmission lineshape observed in the experiment . For optically excited electron energy close to $`E_F`$, this can be achieved by adding the e–e scattering rate $`\gamma _e(E^c)[1f(E^c)][(E^cE_F)^2+(\pi k_BT)^2]`$ to $`\gamma _h^s`$ in Eq. (95). The difference in $`\gamma _e(E^c)`$ for $`E^c`$ below and above $`E_F`$ leads to a lineshape similar to that expected from the combination of red–shift and broadening \[see Fig. 14(a)\].
In Figs. 14(b) and (c) the differential transmission spectra are shown with decreasing nanoparticle size. For $`R=2.5`$ nm, the apparent red–shift is reduced \[see Fig. 14(b)\]. This change can be explained as follows. Since here $`\omega _s\mathrm{\Delta }`$, the SP is damped by the interband excitations. This broadens the spectra for $`\omega >\omega _s`$, so that the absorption peak is asymmetric. The $`d`$–hole scattering with the SP enhances the damping; since the $`\omega `$–dependence of $`\gamma _h^s`$ follows that of $`\alpha `$, this effect is larger above the resonance. On the other hand, the efficiency of the scattering increases with temperature, as discussed above. Therefore, for short time delays, the relative increase in the absorption is larger for $`\omega >\omega _s`$. With decreasing size, the strength of this effect increases further, leading to an apparent blue–shift \[see Fig. 14(c)\]. Such a strong change in the absorption dynamics originates from the $`R^3`$ dependence of the $`d`$–hole scattering rate; reducing the size by the factor of two results in an enhancement of $`\gamma _h^s`$ by an order of magnitude.
In Fig. 15, the time evolution of the differential transmission are shown for several frequencies close to $`\omega _s`$. It can be seen that the relaxation is slowest at the SP resonance; this characterizes the robustness of the collective mode, which determines the peak position, versus the single–particle excitations, which determine the resonance width. For larger sizes, at which $`\gamma _h^s`$ is small, the change in the differential transmission decay rate with frequency is smoother above the resonance \[see Fig. 15(a)\]. This stems from the asymmetric lineshape of the absorption peak, mentioned above: the absorption is larger for $`\omega >\omega _s`$, so that its relative change with temperature is weaker. For smaller nanoparticle size, the decay rates become similar above and below $`\omega _s`$ \[see Fig. 15(b)\]. This change in the frequency dependence is related to the stronger SP damping for $`\omega >\omega _s`$ due to the $`d`$–hole scattering, as discussed above. Since this additional damping is reduced with decreasing temperature, the relaxation is faster above the resonance. This rather “nonlinear” relation between the time–evolution of the pump–probe signal and that of the temperature becomes even stronger for smaller sizes \[see Fig. 15(c)\]. In this case, the frequency dependence of the differential transmission decay below and above $`\omega _s`$ is reversed. Note that a frequency dependence consistent with our calculations presented in Fig. 15(b) was, in fact, observed in the experiment of Ref. , shown in Fig. 15(d).
## 8 Conclusions
In this article we presented an overview of ultrafast nonlinear optics experiments that probe the dynamics of modulation–doped quantum wells and metal nanoparticles. We also discussed some examples where the ultrafast dynamics of such confined Fermi seas cannot be described in terms of the usual dephasing and relaxation time approximations. These examples show the important role of dynamical and non–equilibrium many–body correlations during femtosecond time scales.
More specifically, we discussed a recent theory for the ultrafast nonlinear optical response of the FES. We focussed on coherent effects, which dominate the pump–probe spectra during negative time delays and off–resonant excitation conditions. We demonstrated that the dynamical FS response leads to qualitatively different coherent dynamics of the FES as compared to the Hartree–Fock approximation. In the latter case, the time evolution of the resonance bleaching is governed by the dephasing time. In contrast, in the former case, polarization interference and e–h correlation effects dominate. This results in an initial fast FES dynamics, with a response time determined by the characteristic inverse Coulomb energy $`E_M`$ as well as an apparent resonance enhancement during negative time delays, followed by a long–time decay determined by the dephasing time. Such dynamical features should be observable in ultrafast PP experiments. Using a simple model, we showed that the different dynamics of the FES and Hartree-Fock treatment can be attributed to the non–Lorentzian broadening of the HFA bound state due to its interactions with the gapless FS excitations, a process which is, of course, beyond the dephasing time approximation. In addition, we showed that the pump excitation directly affects the strength of the e–h scattering processes, which changes the frequency dependence of the resonance. The latter can be thought of as an excitation–induced dephasing effect that leads to a transient enhancement of the FES. These results indicate that ultrafast spectroscopy provides a powerful tool to study the role of correlations in the nonlinear response of a Fermi liquid during time scales shorter than the dephasing times. The dynamical features discussed above can also be used as an experimental signal to probe the crossover from FES to exciton bound state (exciton Mott transition) as a function of the FS density.
The above e–h correlation effects and the fundamental differences in the dynamics between an exciton state and a FES many–body resonance can be best observed experimentally in asymmetric modulation–doped quantum wells in the case where the lowest occupied subband is in close proximity to the second unoccupied subband and is coupled strongly to it via intersubband e–h scattering. In such a system there is a strong redistribution of the oscillator strength between the FES and exciton peaks which is caused by the different dynamics of the FES and exciton components of the hybrid as well as by their coupling due to the e–h correlations. This originates from the dynamical Fermi sea response and leads to a strong transient changes in the PP spectra.
We also discussed the role of size–dependent correlations in the electron relaxation in small metal particles. We identified a new mechanism of quasiparticle scattering, mediated by collective surface excitations, which originates from the surface–induced dynamical screening of the e–e interactions. The behavior of the corresponding scattering rates with varying energy and temperature differs substantially from that in the bulk metal. We showed that the conduction electron scattering rate increases with energy, in sharp contrast to the bulk behavior, which could be observed in two–photon photoemission measurements. We also found that in noble metal particles, the resonant energy dependence of the $`d`$–hole scattering rate affects strongly the differential absorption.
An important aspect of the SP–mediated scattering is its strong dependence on size. Our estimates show that it becomes comparable to the usual Fermi–liquid scattering in nanometer–sized particles. This size regime is, in fact, intermediate between “classical” particles with sizes larger than 10 nm, where the bulk–like behavior dominates, and very small clusters with only dozens of atoms, where the metallic properties are completely lost. Although the static properties of nanometer–sized particles are also size–dependent, the deviations from their bulk values do not change the qualitative features of the electron dynamics. In contrast, the size–dependent many–body effects discussed here do affect the dynamics in a significant way during time scales comparable to the relaxation times. The SP–mediated interband scattering reveals itself in the size–dependence of the transient pump–probe spectra. In particular, as the nanoparticle size decreases, the calculated time–resolved differential absorption lineshape shows a transition from an apparent redshift to a blueshift. This transition, absent in the RPA, comes from the correlations between collective surface and single–particle excitations. At the same time, near the SP resonance, these correlation leads to significant size–dependent changes in the frequency dependence of the relaxation time of the pump–probe signal. These results indicate the need for a systematic experimental studies of the size–dependence of the transient nonlinear optical response, as we approach the transition from boundary–constrained nanoparticles to molecular clusters. We expect that in the coming years ultrafast nonlinear optical spectroscopy will provide new insight into many–body effects in strongly correlated, and magnetic Fermi sea systems .
## Acknowledgements
This work was supported by the NSF grant ECS-9703453 and by HITACHI Ltd. Many of the calculations discussed in this review were performed by N. Primozich. It is a pleasure to acknowledge numerous fruitful discussions with D. S. Chemla and J.–Y. Bigot.
## Appendix A APPENDIX A
In this appendix we briefly outline the formalism of Ref. . The pump–probe signal is determined by the polarization,
$$P(t)=\mu e^{i\omega _pt}\mathrm{\Psi }(t)|U|\mathrm{\Psi }(t),$$
(98)
where the state $`|\mathrm{\Psi }(t)`$ satisfies the time-dependent Schrödinger equation,
$$\left[i\frac{}{t}H_{tot}(t)\right]|\mathrm{\Psi }(t)=\left[i\frac{}{t}HH_p(t)H_s(t)\right]|\mathrm{\Psi }(t)=0,$$
(99)
with the Hamiltonians $`H`$ and $`H_{p,s}(t)`$, given by Eqs. (16)–(4.1). The Hilbert space of the bare semiconductor, i.e. in the absence of optical fields, consists of disconnected subspaces $`\zeta \{\nu _{eh}\}`$ which are labeled by the number of (interband) e–h pairs, $`\nu _{eh}`$. The corresponding bare Hamiltonian, $`H`$, conserves the number of e–h pairs in each band separately and in the $`\zeta \{\nu _{eh}\}`$-basis has a block–diagonal form. The Hamiltonians $`H_{p,s}(t)`$ couple the different subspaces $`\zeta \{\nu _{eh}\}`$ by causing interband transitions.
In the description of PP experiments, we are interested only in the polarization component propagating along the probe direction $`𝐤_s`$. For a weak probe, the nonlinear signal then arises from the linear response of the system in the presence of the pump, described by the time–dependent Hamiltonian $`H+H_p(t)`$, to the probe–induced perturbation $`H_s(t)`$. Within $`\chi ^{(3)}`$, the above is true even for comparable pump and probe amplitudes. However, since, in contrast to $`H`$, the Hamiltonian $`H+H_p(t)`$ does not conserve the number of carriers in each band separately, the calculation of the linear response function is not practical. Therefore, we seek to replace $`H+H_p(t)`$ by an effective Hamiltonian $`\stackrel{~}{H}(t)`$ that does conserve the number of e–h pairs in each of its Hilbert subspaces (i.e., is “block–diagonal”). As derived in Ref. , this can be accomplished in any given order in the pump field $`_p(t)`$ and for any pulse duration by using a time–dependent Schrieffer–Wolff /Van Vleck canonical transformation . Here it is sufficient to “block–diagonalize” the Hamiltonian $`H+H_p(t)`$ up to the second order in $`_p(t)`$. The transformation that achieves this has the form $`e^{\widehat{T}_2}e^{\widehat{T}_1}[H+H_p(t)]e^{\widehat{T}_1}e^{\widehat{T}_2}`$, where the anti–Hermitian operators $`\widehat{T}_1(t)`$ and $`\widehat{T}_2(t)`$ create/annihilate one and two e–h pairs, respectively.
We proceed with the first step and eliminate the single-pair pump-induced transitions in the time–dependent Schrödinger equation of the pump/bare–semiconductor system,
$$\left[i\frac{}{t}HH_p(t)\right]|\mathrm{\Psi }(t)=0.$$
(100)
This is achieved by substituting $`|\mathrm{\Psi }(t)=e^{\widehat{T}_1(t)}|\chi (t)`$ and acting with the operator $`e^{\widehat{T}_1(t)}`$ on the lhs of Eq. (100),
$$e^{\widehat{T}_1(t)}\left[i\frac{}{t}H\right]e^{\widehat{T}_1(t)}|\chi (t)=e^{\widehat{T}_1(t)}\left[H_p(t)\right]e^{\widehat{T}_1(t)}|\chi (t).$$
(101)
The anti–Hermitian operator $`\widehat{T}_1(t)`$ has a decomposition
$$\widehat{T}_1(t)=\widehat{𝒫}(t)e^{i𝐤_p𝐫}\widehat{𝒫}^{}(t)e^{i𝐤_p𝐫},$$
(102)
where $`\widehat{𝒫}^{}(t)`$ and $`\widehat{𝒫}(t)`$ create and annihilate single e–h pairs, respectively. The effective Hamiltonian can be found from the condition that the terms describing single–pair interband transitions cancel each other in Eq. (101). In Ref. , it was shown that the multiple commutators of $`\widehat{𝒫}(t)`$ with its time derivatives can be eliminated from Eq. (101) to all orders. By expanding Eq. (101) and neglecting third or higher order terms in $`\widehat{𝒫}`$, we obtain the following equation:
$$i\frac{\widehat{𝒫}^{}(t)}{t}=[H,\widehat{𝒫}^{}(t)]+\mu _p(t)U^{},$$
(103)
with initial condition $`\widehat{𝒫}^{}(\mathrm{})=0`$. The formal solution is
$$\widehat{𝒫}^{}(t)=i\mu _{\mathrm{}}^t𝑑t^{}_p(t^{})e^{iH(tt^{})}U^{}e^{iH(tt^{})}.$$
(104)
Note that, since the Hamiltonian $`H`$ conserves the number of e–h pairs and the optical transition operator $`U^{}`$ creates a single e–h pair, $`\widehat{𝒫}^{}(t)`$ also creates a single e–h pair. Furthermore, since both $`H`$ and $`U^{}`$ conserve momentum, so does $`\widehat{𝒫}^{}(t)`$. Eq. (101) then takes the form
$`[i{\displaystyle \frac{}{t}}\stackrel{~}{H}(t)]|\chi (t)=+{\displaystyle \frac{\mu }{2}}[_p(t)e^{2i𝐤_p𝐫}[U^{},\widehat{𝒫}^{}(t)]+\mathrm{H}.\mathrm{c}.]|\chi (t),`$ (105)
where
$$\stackrel{~}{H}(t)=H+\frac{\mu }{2}(_p(t)[\widehat{𝒫}(t),U^{}]+\mathrm{H}.\mathrm{c}.)$$
(106)
is the sought time–dependent effective Hamiltonian that conserves the number of e–h pairs and $`\widehat{𝒫}^{}(t)`$ is given by Eq. (103). Note that, since $`\widehat{𝒫}^{}(t)`$ is linear in the pump field $`_p`$, the pump–induced term in $`\stackrel{~}{H}(t)`$ \[second term in Eq. (106)\] is quadratic ($`_p_p^{}`$). The rhs of Eq. (105) describes the pump–induced two–pair transitions. These can be eliminated as well by performing a second canonical transformation, $`|\chi (t)=e^{\widehat{T}_2(t)}|\mathrm{\Phi }(t)`$. Following the same procedure, we use the anti–Hermiticity of $`\widehat{T}_2(t)`$ to decompose it as
$$\widehat{T}_2(t)=\widehat{𝒫}_2(t)e^{2i𝐤_p𝐫}𝒫_2^{}(t)e^{2i𝐤_p𝐫},$$
(107)
where $`𝒫_2^{}(t)`$ and $`\widehat{𝒫}_2(t)`$ create and annihilate two e–h pairs, respectively. Substituting the above expression for $`|\chi (t)`$ into Eq. (105) and requiring that all two–pair transitions cancel out, we obtain the following equation for $`𝒫_2^{}(t)`$,
$$i\frac{𝒫_2^{}(t)}{t}=[\stackrel{~}{H}(t),𝒫_2^{}(t)]+\frac{\mu }{2}_p(t)[\widehat{𝒫}^{}(t),U^{}].$$
(108)
Note that $`\widehat{𝒫}_2^{}(t)`$ only affects the PP polarization via higher order ($`_p^4`$) corrections, which are neglected here. However, it does determine the four–wave–mixing (FWM) polarization (see below).
To obtain the condition of validity of this approach, it is useful to write down a formal solution (104) of Eq. (103) in the basis of the N–hole many–body eigenstates, $`|\alpha N`$, with energies $`E_{\alpha N}`$, of the Hamiltonian $`H`$. Here the index $`\alpha `$ labels all the other quantum numbers, so that $`N`$=0 corresponds to the semiconductor ground state $`|0`$, $`|\alpha 1`$ denotes the one–pair states (exciton eigenstates in the undoped case, with $`\alpha `$ labeling both bound and scattering states), $`|\alpha 2`$ denotes the two–pair (biexciton in the undoped case) eigenstates, etc. In this basis, the solution of Eq. (103) can be written as
$$\frac{\beta N+1|\widehat{𝒫}^{}(t)|\alpha N}{\beta N+1|U^{}|\alpha N}=i\mu _{\mathrm{}}^t𝑑t^{}_p(t^{})e^{i(tt^{})\left(\mathrm{\Omega }+\mathrm{\Delta }E_{\alpha \beta }^N\right)}e^{\mathrm{\Gamma }(tt^{})},$$
(109)
where we separated out the detuning $`\mathrm{\Omega }`$ and denoted $`\mathrm{\Delta }E_{\alpha \beta }^N=E_{\beta N+1}E_{\alpha N}`$. It can be seen that for resonant excitation (small $`\mathrm{\Omega }`$) the rhs of Eq. (109) is of the order of $`\mu _pt_p`$. Thus, for short pulses, this parameter justifies the expansion in terms of the optical fields. For off-resonant excitation, this expansion is valid for any pulse duration provided that $`\mu _p/\mathrm{\Omega }<1`$. Similar conditions can be obtained for the two–pair transition described by $`\widehat{𝒫}_2`$.
The nonlinear polarization Eq. (98) can now be expressed in terms of the linear response to the probe field:
$$P(t)=\mu e^{i\omega _pt}\mathrm{\Psi }(t)|U|\mathrm{\Psi }(t)=\mu e^{i\omega _pt}\mathrm{\Phi }(t)|U_T(t)|\mathrm{\Phi }(t),$$
(110)
where, in the first order in $`_s(t)`$, the state $`|\mathrm{\Phi }(t)`$ is given by
$$|\mathrm{\Phi }(t)=|\mathrm{\Phi }_0(t)i\mu _{\mathrm{}}^tdt^{}𝒦(t,t^{})[_s(t^{})e^{i𝐤_s𝐫+i\omega _p\tau }U_T^{}(t^{})+\mathrm{H}.\mathrm{c}.]|\mathrm{\Phi }_0(t^{}).$$
(111)
Here $`𝒦(t,t^{})`$ is the time-evolution operator satisfying
$$i\frac{}{t}𝒦(t,t^{})=\stackrel{~}{H}(t)𝒦(t,t^{}),$$
(112)
and $`U_T^{}(t)=e^{\widehat{T}_2(t)}e^{\widehat{T}_1(t)}U^{}e^{\widehat{T}_1(t)}e^{\widehat{T}_2(t)}`$ is the (transformed) optical transition operator. In Eq. (111), $`|\mathrm{\Phi }_0(t)=𝒦(t,\mathrm{})|0`$ is the time–evolved ground state $`|0`$; since $`\stackrel{~}{H}(t)`$ conserves the number of e–h pairs, $`|\mathrm{\Phi }_0(t)`$ contains no interband e–h pairs (in undoped semiconductors, it coincides with the ground state, $`|\mathrm{\Phi }_0(t)=|0`$). ¿From Eqs. (111) and (110), the polarization $`P(t)`$ takes the form
$`P(t)=i\mu ^2e^{i\omega _pt}`$ $`{\displaystyle _{\mathrm{}}^t}dt^{}[\mathrm{\Phi }_0(t)|U_T(t)𝒦(t,t^{})[_s(t^{})e^{i𝐤_s𝐫+i\omega _p\tau }U_T^{}(t^{})`$ (113)
$`+_s^{}(t^{})e^{i𝐤_s𝐫i\omega _p\tau }U_T(t^{})]|\mathrm{\Phi }_0(t^{})`$
$`\mathrm{\Phi }_0(t^{})|[_s(t^{})e^{i𝐤_s𝐫+i\omega _p\tau }U_T^{}(t^{})`$
$`+_s^{}(t^{})e^{i𝐤_s𝐫i\omega _p\tau }U_T(t^{})]𝒦(t^{},t)U_T(t)|\mathrm{\Phi }_0(t)].`$
The above expression for the total polarization contains contributions propagating in various directions. To obtain the polarization propagating in a specific direction, one has to expand the effective–transition operator $`U_T(t)`$ in terms of $`\widehat{T}_1`$ and $`\widehat{T}_2`$. Using Eqs. (102) and (107) and keeping only terms contributing to PP and FWM polarizations, we obtain
$`U_T^{}(t)=U_1^{}(t)+U_2^{}(t)e^{i𝐤_p𝐫}+U_{FWM}(t)e^{2i𝐤_p𝐫}+\mathrm{},`$ (114)
where (to lowest order in the pump field)
$`U_1^{}(t)`$ $`=`$ $`U^{}+{\displaystyle \frac{1}{2}}[\widehat{𝒫}(t),[U^{},\widehat{𝒫}^{}(t)]]+{\displaystyle \frac{1}{2}}[\widehat{𝒫}^{}(t),[U^{},\widehat{𝒫}(t)]],`$
$`U_2^{}(t)`$ $`=`$ $`[\widehat{𝒫}^{}(t),U^{}],`$
$`U_{FWM}^{}(t)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[[U,\widehat{𝒫}^{}(t)],\widehat{𝒫}^{}(t)][U,\widehat{𝒫}_2^{}(t)].`$ (115)
Here operators $`U_1^{}(t)\stackrel{~}{U}(t)`$ and $`U_{FWM}^{}(t)`$ create one e–h pair, while $`U_2^{}(t)`$ creates two e–h pairs (note that $`U_{FWM}`$ in Eq. (A) annihilates an e–h pair).
### Pump–probe polarization
In order to extract the PP polarization from Eq. (113), one should retain only terms that are proportional to $`e^{i𝐤_s𝐫}`$. Substituting Eqs. (114) into Eq. (113), we obtain $`P_{𝐤_s}(t)=e^{i𝐤_s𝐫i\omega _p(t\tau )}[\stackrel{~}{P}^{(1)}(t)+\stackrel{~}{P}^{(2)}(t)]`$, where
$`\stackrel{~}{P}^{(1)}(t)=i\mu ^2{\displaystyle _{\mathrm{}}^t}𝑑t^{}_s(t^{})\mathrm{\Phi }_0(t)|\stackrel{~}{U}(t)𝒦(t,t^{})\stackrel{~}{U}^{}(t^{})|\mathrm{\Phi }_0(t^{}),`$ (116)
and
$`\stackrel{~}{P}^{(2)}(t)=i\mu ^2{\displaystyle _{\mathrm{}}^t}𝑑t^{}_s(t^{})\mathrm{\Phi }_0(t)|U_2(t)𝒦(t,t^{})U_2^{}(t^{})|\mathrm{\Phi }_0(t^{}).`$ (117)
Note that the above formulae apply to both undoped and doped semiconductors.
Eqs. (116)–(117) express the nonlinear PP polarization in terms of the linear response to the probe field of a system described by the time–dependent effective Hamiltonian (106). Such a form for the nonlinear response allows one to distinguish between two physically distinct contributions to the optical nonlinearities. Assuming that a short probe pulse arrives at $`t=\tau `$, consider the first term, Eq. (116), which gives the single–pair (exciton for undoped case) contribution to the PP polarization. At negative time delays, $`\tau <0`$, the probe excites an e–h pair, described by state $`\stackrel{~}{U}^{}(\tau )|\mathrm{\Phi }_0(\tau )U^{}|0`$; since the probe arrives before the pump, the effective transition operator coincides with the “bare”one \[see Eqs. (A) and (104)\]. The first contribution to the optical nonlinearities comes from the effective Hamiltonian, $`\stackrel{~}{H}(t)`$, governing the propagation of that interacting e–h pair in the interval $`(\tau ,t)`$ via the time–evolution operator $`𝒦(t,\tau )`$. Note that since the pump pulse arrives at $`t=0`$, for $`|\tau |\mathrm{\Gamma }^1`$, the negative time–delay signal vanishes. At $`t>0`$, the e–h pair (exciton in the undoped case) “feels” the effect of the pump via mainly the transient bandgap shift, leading, e.g., to ac–Stark effect, and the change in the band dispersions (increase in effective mass/density of states), leading to enhanced e–h scattering (exciton binding energy in the undoped case). Note that $`\stackrel{~}{H}(t)`$ also contains a contribution coming from the interactions between probe– and pump–excited e–h pairs, which are however perturbative in the doped case for short pulses or off–resonant excitation and lead to subdominant corrections. Importantly, the response of the system to the optically–induced corrections in $`\stackrel{~}{H}(t)`$ takes into account all orders in the pump field, which is necessary for the adequate description, e.g., of the ac–Stark effect and the pump–induced changes in the e–h correlations. Indeed, although the pump–induced term in Eq. (106) is quadratic in $`_p`$, the time–evolution of the interacting e–h pair is described without expanding $`𝒦(t,\tau )`$ in the pump field. On the other hand, the third–order polarization ($`\chi ^{(3)}`$) can be obtained by expanding $`𝒦(t,\tau )`$ to the lowest order. The second contribution to the optical nonlinearities comes from the matrix element of the final state, $`\mathrm{\Phi }_0(t)|\stackrel{~}{U}(t)`$ in Eq. (116). The latter, given by Eq. (A), contains the lowest order (quadratic) pump–induced terms which describe the Pauli blocking, pair–pair, and pair–FS interaction effects (exciton–exciton interactions in the undoped case ). Note that the matrix element of the initial state contributes for positive time delays, i.e., if the probe arrives after the pump pulse. In this case, however, the pump–induced term in the effective Hamiltonian (106) vanishes (since it lasts only for the duration of the pulse) so that for positive $`\tau >t_p`$ the PP signal is determined by the matrix elements rather than by $`\stackrel{~}{H}(t)`$. If the probe arrives during the interaction of the system with the pump pulse ($`\tau t_p`$), both the effective Hamiltonian and the matrix elements contribute to the polarization. In this case, there is also a biexcitonic contribution \[given by Eq. (117)\] coming from a simultaneous excitation of two e–h pairs by the pump and the probe. However, such a biexciton state does not contribute to negative ($`\tau <0`$) time delays. As can be seen from the above discussion, our theory separates out a number of contributions that play a different role for different time delays and excitation conditions.
### FWM polarization
By extracting from Eq. (113) all the terms propagating in the the FWM direction, $`2𝐤_p𝐤_s`$, we obtain (for delta–function probe $`_s(t)=_s\delta (t\tau )`$)
$`P_{FWM}(t)=`$ $`ie^{i(2𝐤_p𝐤_s)𝐫i\omega _p(t+\tau )}\theta (t\tau )\mu ^2_s^{}`$ (118)
$`\times [\mathrm{\Phi }_0(t)|U𝒦(t,\tau )U_{FWM}^{}(\tau )|\mathrm{\Phi }_0(\tau )(t\tau )],`$
where the FWM transition operator $`U_{FWM}^{}(t)`$ is given by Eq. (A). It is convenient to express $`U_{FWM}^{}(t)`$ in terms of the “irreducible” two–pair operator $`W^{}(t)=\frac{1}{2}\widehat{𝒫}^2\widehat{𝒫}_2^{}`$, satisfying
$$i\frac{W^{}(t)}{t}=[\stackrel{~}{H}(t),W^{}(t)]+\mu _p(t)U^{}\widehat{𝒫}^{}(t).$$
(119)
In terms of $`W^{}(t)`$, the state $`U_{FWM}^{}(t)|\mathrm{\Phi }_0(t)`$ in Eq. (118) can be presented as a sum of two– and one–pair contributions:
$`U_{FWM}^{}(t)|\mathrm{\Phi }_0(t)=UW^{}(t)|\mathrm{\Phi }_0(t)\widehat{𝒫}^{}(t)U\widehat{𝒫}^{}(t)|\mathrm{\Phi }_0(t).`$ (120)
For undoped semiconductors, $`|\mathrm{\Phi }_0(t)`$ represents the ground state $`|0`$ of $`H`$. The operator $`W^{}(t)`$, being quadratic in the pump field ($`_p^2`$), describes the simultaneous excitation of two interacting e–h pairs by the pump pulse and includes the biexciton and exciton–exciton scattering effects. Introducing the amplidutes $`\chi _𝐤`$ and $`\mathrm{\Phi }_𝐤`$ as
$$U_{FWM}^{}(t)|0=\underset{𝐤}{}\chi _𝐤(t)a_𝐤^{}b_𝐤^{}|0,$$
(121)
and
$$\stackrel{~}{𝒦}(t,t^{})U_{FWM}^{}(t^{})|0=\underset{𝐤}{}\mathrm{\Phi }_𝐤(t,t^{})a_𝐤^{}b_𝐤^{}|0,$$
(122)
where $`\mathrm{\Phi }_𝐤(t,t^{})`$ satisfies
$$i_t\mathrm{\Phi }_𝐤(t,t^{})=\underset{𝐪}{}0|b_𝐤a_𝐤\stackrel{~}{H}(t)a_𝐪^{}b_𝐪^{}|0\mathrm{\Phi }_𝐪(t,t^{}),$$
(123)
with initial condition $`\mathrm{\Phi }_𝐤(t,t)=\chi _𝐤(t)`$, one finally obtains for the FWM polarization
$$P_{FWM}(t)=ie^{i(2𝐤_p𝐤_s)𝐫i\omega _p(t+\tau )}\theta (t\tau )\mu ^2_s^{}\underset{𝐤}{}[\mathrm{\Phi }_𝐤(t,\tau )\mathrm{\Phi }_𝐤(\tau ,t)].$$
(124)
Note that the third order polarization \[corresponding to $`\chi ^{(3)}`$\] is obtained by replacing $`\stackrel{~}{H}`$ with $`H`$.
## Appendix B APPENDIX B
In this appendix we clarify our convention for the time delay $`\tau `$ and relate it to the most commonly used conventions in PP and FWM. In the generic experimental configuration two laser pulses $`E_1(t)e^{i𝐤_1𝐫i\omega (tt_1)}`$, and $`E_2(t)e^{i𝐤_2𝐫i\omega (tt_2)}`$, respectively centered at time $`t=t_1`$ and $`t=t_2`$ are incident on a sample. Let us define $`\mathrm{\Delta }t`$ as,
$$\mathrm{\Delta }t=t_1t_2$$
(125)
and consider a FWM experiment where the signal is measured in the direction $`2𝐤_2𝐤_1`$. Then for a two-level-atom, the signal vanishes for $`\mathrm{\Delta }t<0`$, while for $`\mathrm{\Delta }t>0`$ its amplitude, which decays with time as $`e^{t/T_2}`$, is determined by the Pauli blocking. In a system with Coulomb interactions (such as a semiconductor) a FWM signal is observed both for $`\mathrm{\Delta }t<0`$ and $`\mathrm{\Delta }t>0`$. The $`\mathrm{\Delta }t<0`$ signal is entirely due to the Coulomb interaction.
In PP experiments, one usually chooses one pulse (the “pump”) to have an amplitude $`E_p`$ much larger than that of the other pulse (the “probe”), $`E_s`$. As discussed in the text, a weak probe measures the linear response of the system (bare or dressed by the pump). If we choose that $`E_p=E_2`$, the pump induces coherent and incoherent populations when it arrives in the sample before the probe i.e. for $`t_2<t_1`$. This is usually defined as “positive” time delay $`\tau =t_2t_1>0`$ in the PP literature. Note that $`\tau =\mathrm{\Delta }t`$, i.e., the “regular” sequence in PP experiments is the reverse of that of FWM experiments. For $`\tau <0`$, the origin of the PP signal is that the probe creates a linear polarization in the sample which lasts for time $`T_2=\mathrm{\Gamma }^1`$ and, consequently, is scattered by polarization excited by the pump field. The signal observed for $`\tau <0`$ is therefore due to coherent effects.
In FWM experiments, there is no restriction on the magnitude of the two incident fields $`E_2`$ and $`E_1`$, which are often chosen to have amplitudes of the same order. Note however that, at the $`\chi ^{(3)}`$ level, the FWM and pump–probe polarizations are linear in the $`E_1(t)`$ field and thus the above linear response calculation applies even for comparable pump and probe amplitudes.
## Appendix C APPENDIX C
In this Appendix we present the explicit expressions for the renormalized transition matrix elements in the presence of the pump excitation. The direct transition matrix element is given by
$`M_𝐩(t)`$ $`=1|𝒫_{eh}(𝐩,t)|^2+[𝒫_{\mathrm{eh}}^{}(𝐩,t){\displaystyle \underset{k^{}<k_F}{}}𝒫_{eh}^\mathrm{e}(\mathrm{𝐩𝐤}^{};𝐤^{};t)+\mathrm{H}.\mathrm{c}.]`$ (126)
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{p^{}>k_F}{}}𝒫_{eh}^{}(𝐩^{},t)[𝒫_{eh}^\mathrm{e}(𝐩^{}𝐩;𝐩^{};t)𝒫_{eh}^\mathrm{h}(𝐩^{}𝐩;𝐩^{};t)`$
$`𝒫_{eh}^\mathrm{h}(𝐩^{}𝐩;𝐩;t)+𝒫_{eh}^\mathrm{e}(𝐩^{}𝐩;𝐩;t)]`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{p^{}>k_F}{}}\left[𝒫_{eh}(𝐩,t)+𝒫_{eh}(𝐩^{},t)\right]`$
$`\times \left[𝒫_{eh}^\mathrm{e}(\mathrm{𝐩𝐩}^{};𝐩;t)+𝒫_{eh}^\mathrm{h}(\mathrm{𝐩𝐩}^{};𝐩;t)\right]^{}.`$
The first term on the rhs of the above equation describes the phase space filling contribution, while the rest of the terms are due to the mean field pair–pair and pair–FS interactions.
The pump–induced indirect transition matrix element is given by
$`M_{\mathrm{𝐩𝐩}^{}𝐤}(t)=`$ $`\left[𝒫_{eh}(𝐤,t)𝒫_{eh}(𝐩+𝐩^{}𝐤,t)\right]^{}𝒫_{eh}^\mathrm{e}(\mathrm{𝐩𝐩}^{};𝐤;t)`$ (127)
$`\left[𝒫_{eh}(𝐩,t)+𝒫_{eh}(𝐩^{},t)\right]^{}`$
$`\times \left[𝒫_{eh}^\mathrm{e}(\mathrm{𝐩𝐩}^{};𝐩;t)+𝒫_{eh}^\mathrm{h}(\mathrm{𝐩𝐩}^{};𝐩+𝐩^{}𝐤;t)\right]`$
$`+𝒫_{eh}(𝐩+𝐩^{}𝐤,t)[𝒫_{eh}^\mathrm{e}(𝐤,𝐩+𝐩^{}𝐤;𝐩^{};t)`$
$`𝒫_{eh}^\mathrm{e}(𝐤,𝐩+𝐩^{}𝐤;𝐩;t)]^{}`$
$`+𝒫_{eh}^{}(𝐤,t)[𝒫_{eh}^\mathrm{h}(𝐤,𝐩+𝐩^{}𝐤;𝐩^{};t)+𝒫_{eh}^\mathrm{e}(𝐩,𝐩^{};𝐤;t)`$
$`𝒫_{eh}^\mathrm{e}(𝐩,𝐩^{};𝐩+𝐩^{}𝐤;t)𝒫_{eh}^\mathrm{h}(𝐤,𝐩+𝐩^{}𝐤;𝐩;t)]`$
$`+𝒫_{eh}(𝐩,t)𝒫_{eh}^\mathrm{e}(𝐤,𝐩+𝐩^{}𝐤;𝐩^{};t)`$
$`𝒫_{eh}(𝐩^{},t)𝒫_{eh}^e(𝐤,𝐩+𝐩^{}𝐤;𝐩;t).`$
The effective e–h potential is given by
$`\upsilon _{eh}(𝐪;\mathrm{𝐤𝐤}^{};t)=\upsilon (𝐪){\displaystyle \frac{\mu }{2}}_p(t)`$ $`[𝒫_{eh}^e(𝐤+𝐪,𝐤^{};𝐤^{}+𝐪;t)`$ (128)
$`+𝒫_{eh}^h(𝐤,𝐤^{}+𝐪;𝐤^{};t)`$
$`+𝒫_{eh}^e(𝐤,𝐤^{}+𝐪;𝐤^{};t)`$
$`+𝒫_{eh}^h(𝐤+𝐪,𝐤^{};𝐤^{}+𝐪;t)],`$
and the effective e–e potential is given by
$`\upsilon _{ee}(𝐪;\mathrm{𝐤𝐤}^{};t)=\upsilon (𝐪)+{\displaystyle \frac{\mu }{4}}_p(t)`$ $`[𝒫_{eh}^e(𝐤+𝐪,𝐤^{}𝐪;𝐤^{};t)`$ (129)
$`𝒫_{eh}^e(𝐤+𝐪,𝐤^{}𝐪;𝐤;t)`$
$`+𝒫_{eh}^e(𝐤,𝐤^{};𝐤^{}𝐪;t)`$
$`𝒫_{eh}^e(𝐤,𝐤^{};𝐤+𝐪;t)].`$ |
warning/0001/hep-th0001146.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The present article is an extended account of the work reported in Ref. . We consider a classical system of electric charges which make a source of the electromagnetic field and move in the self field. However, we take into account that this source is immersed in the real vacuum, and the field that it generates excites the vacuum charges. The problem is figuring out the vacuum back-reaction on the motion of the source.
The electromagnetic field generated by a source is a solution of the expectation-value equations in the in-vacuum state. For this state to exist , the source is assumed asymptotically static in the past. In consequence of this assumption, the solution always contains a contribution of the static vacuum polarization whose principal effect is screening of the monopole moment of the source. The polarization occurs in the whole of space and increases infinitely in the vicinity of the source thereby causing the ultraviolet disaster. However, we show that this infinity does not affect the motion of the source. After the self-action of a pointlike charge is properly eliminated, the force exerted by the source on itself is finite. By choosing the spatial scale of the system exceeding the Compton size of the vacuum particle, we abstract ourselves from the static polarization altogether.
The effect that we are concerned about is the vacuum instability caused by a nonstationarity of the source . A nonstationary electromagnetic field is capable of creating in the vacuum real particles having the electric charge. When the frequency of the source exceeds the threshold of pair creation, it emits a flux of energy and charge carried by the created particles. The main question is how much energy can be extracted from a source by means of this mechanism? An attempt to answer this question without taking into account the back-reaction of the vacuum on the motion of the source leads to a contradiction with the energy conservation law . The radiation rate grows unboundedly with the energy of the source, and, at a sufficiently high energy, the source appears to give more than it has.
The problem of self-consistent motion is solved below for the simplest model of a pair-creating source. The model is a charged spherical shell expanding in the self field. This choice is made to avoid the complications connected with an emission of the electromagnetic waves. A high-frequency source will generally emit both the electromagnetic waves by the classical mechanism and charged particles by the quantum mechanism. The two radiations overlap in a nontrivial way: the energy of the vacuum of charged particles goes partially into the electromagnetic radiation and amplifies it . Accounting for the vacuum back-reaction is then necessary already for a removal of the infrared disaster, and the problem of restoration of the energy conservation law concerns both components of radiation . By choosing the source spherically symmetric, we exclude the emission of waves, and, thereby, put off the solution of this more complicated problem.
The solution in the case of the spherical shell is in the fact that there emerges a new kinematic bound on the velocity of the source. Raising the energy of the source results in an increase of its acceleration, which causes an intensification of the vacuum particle production, which entails a reinforcement of its back-reaction, which results in a deceleration of the source. As a result, however high the energy may be, the velocity of a charged body cannot approach the speed of light closer than a certain limit. Within a given type of coupling, this limit is universal. It does not depend on the parameters of the source, only on the coupling constant. The back-reaction effect is nonanalytic in the coupling constant and restores completely the conservation laws. Up to 50% of energy and charge can be extracted from the source by raising its initial energy.
The effect of vacuum instability is of significance for rapidly moving, or high-frequency sources. The high-frequency approximation is the only approximation made in the solution. The phenomenological, or axiomatic theory of the vacuum is used in which the expectation-value equations are specified by a set of operator form factors. In the high-frequency approximation, only the polarization operator is involved in the calculation of the induced charge<sup>1</sup><sup>1</sup>1For the calculation of the induced energy, one needs also the gravitational form factors .. Its relevant properties are postulated in Section 2.
In Section 3, a closed set of equations is obtained for the motion of the source in the self electromagnetic field. The technique needed for dealing with the nonlocal expectation-value equations is presented in Section 4. In Section 5, the static solution is considered, valid outside some future light cone. The solution for a moving source is obtained in Section 6, and its ultraviolet behaviour is studied in Sections 7 and 8. In Sections 9-11, the force of the vacuum back-reaction is calculated, and the equation of motion of the source is solved in the high-frequency approximation. The rate of emission of charge is calculated in Section 12.
## 2 Electrically charged source coupled to the vacuum charges
Our starting point is the action for the electromagnetic field generated by a source
$$S=S_{\text{cl}}+S_{\text{vac}}+\mathrm{SS}$$
(2.1)
in which the source is a set of particles with masses $`M_i`$ and charges $`e_i`$ :
$$\mathrm{SS}=\underset{i}{}𝑑s\left(\frac{M_i}{2}g_{\mu \nu }\left(x_i(s)\right)\frac{dx_i^\mu }{ds}\frac{dx_i^\nu }{ds}+e_iA_\mu \left(x_i(s)\right)\frac{dx_i^\mu }{ds}\right).$$
(2.2)
Here $`g_{\mu \nu }`$ is the flat metric
$$g_{\mu \nu }dx^\mu dx^\nu =dt^2+dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2),$$
(2.3)
$`A_\mu (x)`$ is the electromagnetic potential, and $`s`$ is the proper time of the $`i`$-th particle. The quantities
$$M=\underset{i}{}M_i,e=\underset{i}{}e_i$$
(2.4)
are the full mass and charge of the source ($`c=1`$). For definiteness, the charge $`e`$ will be considered positive.
In Eq. (2.1), $`S_{\text{cl}}`$ is the classical action of the electromagnetic field
$$S_{\text{cl}}=\frac{1}{16\pi }dxg^{1/2}F_{\mu \nu }F^{\mu \nu },F_{\mu \nu }=_\mu A_\nu _\nu A_\mu ,$$
(2.5)
and $`S_{\text{vac}}`$ is the effective action of the vacuum charges. The phenomenological, or axiomatic theory of the vacuum will be used in which $`S_{\text{vac}}`$ is taken as the general expansion over the basis of nonlocal invariants :
$$S_{\text{vac}}=\frac{1}{16\pi }𝑑xg^{1/2}F_{\mu \nu }f(\mathrm{})F^{\mu \nu }+O(F\times F\times F\mathrm{}).$$
(2.6)
The higher-order terms of this expansion are of the form
$$𝑑xg^{1/2}f(\mathrm{}_1,\mathrm{}_2,\mathrm{}_3,\mathrm{}\mathrm{}_{1+2},\mathrm{}_{1+3},\mathrm{})F_1\times F_2\times F_3\mathrm{}$$
(2.7)
where $`f`$’s are functions of the D’Alembert operators (the form factors). In (2.7), the operator $`\mathrm{}_n`$ acts on $`F_n`$ , and the operator $`\mathrm{}_{m+n}`$ acts on the product $`F_mF_n`$ . All form factors are assumed to admit spectral representations through the resolvents $`1/(\mu ^2\mathrm{})`$ .
The expectation-value equations of the electromagnetic field, associated<sup>2</sup><sup>2</sup>2For the expectation-value equations there is no direct least-action principle but they differ from the variational equations of the Feynman effective action only by the boundary conditions for the resolvents . with the action (2.1) are the following equations for the current $`j^\alpha `$ :
$$_\beta F^{\alpha \beta }=4\pi j^\alpha ,$$
(2.8)
$$j^\alpha +f(\mathrm{})j^\alpha +O(F\times F\mathrm{})=j_{\text{bare}}^\alpha $$
(2.9)
with the retarded resolvents for $`f(\mathrm{})`$ and the higher-order form factors. The current $`j_{\text{bare}}^\alpha `$ as given by the action (2.2) is of the form
$$j_{\text{bare}}^\alpha (x)=\underset{i}{}𝑑se_i\frac{1}{g^{1/2}(x)}\delta ^{(4)}\left(xx_i(s)\right)\frac{dx_i^\alpha (s)}{ds}.$$
(2.10)
The full set of exact form factors provides a complete phenomenological description of the vacuum of particles having a given type of charge (here the electric charge). On the other hand, the form factors are to be calculated from some quantum-field model, and, even within a given model, this calculation can never be complete let alone the fact that one never knows the ultimate model of the vacuum. The virtue of the axiomatic approach is in the fact that, for the physical questions of interest, the detailed form of the form factors is unimportant. Only some of their properties are important. These properties should be postulated, and the expectation-value problem solved with the form factors which are otherwise arbitrary. The axiomatic approach is model-independent, and at the same time it is a tool for testing various models and approximations therein. The present paper gives an example of this approach.
A possibility of truncating the series (2.6) depends on the expectation-value problem in question which, in its turn, is specified by the properties of the source<sup>3</sup><sup>3</sup>3We assume that there are no incoming electromagnetic waves.. Our concern in the present paper is a high-frequency source creating pairs. Let $`\nu `$ be its typical frequency. In the problem of particle creation by a nonstationary external or mean field, this field is considered high-frequency if the energy $`\mathrm{}\nu `$ dominates both the rest energy of the vacuum particle and its static (Coulomb) energy in this field . The high-frequency approximation is the condition of validity of the expansion (2.6) .
The linear expectation-value equations obtained by truncating Eq. (2.9) are solved by the ansatz
$$j^\alpha =j_{\text{bare}}^\alpha \gamma (\mathrm{})j_{\text{bare}}^\alpha $$
(2.11)
in which $`\gamma (\mathrm{})`$ is some retarded form factor. Its relevant properties are to be postulated. We assume that the function $`\gamma (\mathrm{})`$ is analytic in the complex plane of $`\mathrm{}`$ except at the real negative half-axis where it has a cut:
$$\frac{1}{2\pi \mathrm{i}}[\gamma (\mu ^2\mathrm{i0})\gamma (\mu ^2+\mathrm{i0})]=\mathrm{\Delta }(\mu ^2).$$
(2.12)
The properties of the spectral-mass function $`\mathrm{\Delta }(\mu ^2)`$ that need to be specified are (i) positivity
$$\mathrm{\Delta }(\mu ^2)0,$$
(2.13)
(ii) the presence of a lower bound in the spectrum
$$\mathrm{\Delta }(\mu ^2)\theta (\mu ^24m^2),m0,$$
(2.14)
and (iii) finiteness at large spectral mass
$$\mathrm{\Delta }(\mu ^2)|_{\mu ^2\mathrm{}}=\frac{\kappa ^2}{24\pi }0.$$
(2.15)
Eq. (2.14) introduces the mass of the vacuum particles $`m`$ , and Eq. (2.15) introduces the coupling constant $`\kappa ^2`$ . Finally, the function $`\gamma (\mathrm{})`$ must satisfy the normalization condition
$$\gamma (0)=0.$$
(2.16)
Redefining the spectral-mass function as
$$\mathrm{\Delta }(\mu ^2)=\frac{\kappa ^2}{24\pi }\mathrm{\Gamma }(\mu ^2),\mu ^24m^2,$$
(2.17)
we obtain from Eqs. (2.12)-(2.16)
$$\gamma (\mathrm{})=\frac{\kappa ^2}{24\pi }\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)\left(\frac{1}{\mu ^2\mathrm{}}\frac{1}{\mu ^2}\right).$$
(2.18)
To summarize, we assume that the expectation-value equations (2.11) hold with the form factor (2.18) in which the resolvent $`1/(\mu ^2\mathrm{})`$ is retarded , and $`\mathrm{\Gamma }(\mu ^2)`$ satisfies the conditions
$$\mathrm{\Gamma }(\mu ^2)0,\mathrm{\Gamma }(\mathrm{})=1.$$
(2.19)
In the underlying quantum field theory, Eqs. (2.13)-(2.15) assume (i) positivity of the metric of the physical Hilbert space, (ii) the presence of an energy threshold for pair creation, and (iii) a logarithmic divergence of the charge renormalization. Eq. (2.16) is a condition that $`\kappa ^2`$ is the renormalized coupling constant. For example, in the case of the electron-positron vacuum, the equations above hold with
$$\mathrm{\Gamma }(\mu ^2)=\left(1\frac{4m^2}{\mu ^2}\right)^{1/2}\left(1+\frac{2m^2}{\mu ^2}\right)+\text{ multi-loop contributions}$$
(2.20)
and
$$\kappa ^2=8\alpha +O(\alpha ^2)$$
(2.21)
where $`\alpha `$ is the fine-structure constant. In the case where $`S_{\text{vac}}`$ is the standard loop with the abelian commutator curvature<sup>4</sup><sup>4</sup>4For the standard loop with arbitrary metric, connection, and potential the calculations can be carried out in the general form, and the results tabulated. The one-loop action for any model is then obtained by combining the standard loops. (See Refs. and references therein.)
$$\widehat{}_{\mu \nu }=\widehat{\mathrm{\Omega }}F_{\mu \nu },$$
(2.22)
the coupling constant is defined by the matrix $`\widehat{\mathrm{\Omega }}`$ :
$$\kappa ^2=\text{tr}\widehat{\mathrm{\Omega }}^2,$$
(2.23)
and the spectral-mass function is
$$\mathrm{\Gamma }(\mu ^2)=\left(1\frac{4m^2}{\mu ^2}\right)^{3/2}.$$
(2.24)
## 3 The charged shell expanding in the self field
To avoid complications connected with an emission of the electromagnetic waves , the source will be chosen spherically symmetric, and, moreover, the particles in the action (2.2) will be assumed to pack a thin spherical shell. This amounts to choosing the solution of the form
$$x_i^\mu (s)=\{t(s),r(s),\theta _i,\phi _i\},\frac{d\theta _i}{ds}=0,\frac{d\phi _i}{ds}=0$$
(3.1)
with $`t(s)`$ and $`r(s)`$ independent of $`i`$, and identifying $`i`$ with the set $`\{\theta _i,\phi _i\}`$ . Then the motion of the source boils down to a radial motion in an electric field of a single particle with mass $`M`$ and charge $`e`$ . The electric field is the self field of the shell. An important fact is that the electric field is discontinuous on a charged surface, and the force exerted by the shell on itself is determined by one half of the sum of the electric fields on both sides of the shell . Writing the law of motion of the shell in the form
$$r=\rho (t)$$
(3.2)
one obtains
$$M\frac{d}{dt}\left(\frac{\dot{\rho }}{\sqrt{1\dot{\rho }^2}}\right)=e\frac{E_++E_{}}{2}|_{\text{shell}}$$
(3.3)
where $`E_+`$ and $`E_{}`$ are the electric fields outside and inside the shell.
Any spherically symmetric electromagnetic field is determined by a single function $`𝐞(t,r)`$ which is the charge contained at the time instant $`t`$ inside the sphere of area $`4\pi r^2`$ :
$$𝐞(t,r)=𝑑\overline{x}\overline{g}^{1/2}\theta (r\overline{r})\delta (t\overline{t})\overline{}_\mu \overline{t}j^\mu (\overline{x}).$$
(3.4)
In terms of this function the solution of the conservation equation $`_\alpha j^\alpha =0`$ is
$$4\pi r^2j^\mu =\left(^\mu t\frac{}{r}+^\mu r\frac{}{t}\right)𝐞(t,r),$$
(3.5)
and the solution of the Maxwell equations (2.8) is
$$F_{\mu \nu }=\left(_\mu r_\nu t_\mu t_\nu r\right)E$$
(3.6)
with the electric field
$$E=\frac{𝐞(t,r)}{r^2}.$$
(3.7)
The function $`𝐞(t,r)`$ must satisfy the condition of regularity of the electric field at $`r=0`$
$$𝐞(t,0)=0$$
(3.8)
and the normalization condition
$$𝐞(t,\mathrm{})=e.$$
(3.9)
Since $`j_{\text{bare}}^\alpha `$ in Eq. (2.10) is conserved, it is also of the form (3.5):
$$4\pi r^2j_{\text{bare}}^\mu =\left(^\mu t\frac{}{r}+^\mu r\frac{}{t}\right)𝐞_{\text{bare}}(t,r),$$
(3.10)
and, for $`𝐞_{\text{bare}}(t,r)`$ , Eq. (2.10) yields the obvious result
$$𝐞_{\text{bare}}(t,r)=e\theta (r\rho (t)).$$
(3.11)
Thus, owing to the conservation of the current $`j^\alpha `$ , which is a corollary of the expectation-value equations (2.11), only one of these equations is independent. Finally, on account of Eq. (3.7), the equation of motion of the shell (3.3) takes the form
$$M\frac{d}{dt}\left(\frac{\dot{\rho }}{\sqrt{1\dot{\rho }^2}}\right)=e\frac{𝐞_+(t)+𝐞_{}(t)}{2\rho ^2}$$
(3.12)
where
$$𝐞_\pm (t)=𝐞(t,\rho (t)\pm 0).$$
(3.13)
Eqs. (2.11), (3.5), and (3.10)-(3.13) make a closed set of equations for $`𝐞(t,r)`$ and $`\rho (t)`$ .
The setting of the problem with the in-vacuum of quantum fields implies that the external or mean fields generated by the source are asymptotically static in the past . Accordingly, it will be assumed that, before some time instant $`t=t_{\text{start}}`$ , the shell was kept at some constant value of $`r`$, $`r=r_{\text{min}}`$ , and next was let go. Eq. (3.12) will thus be solved with the initial conditions
$$\rho |_{t_{\text{start}}}=r_{\text{min}},\dot{\rho }|_{t_{\text{start}}}=0.$$
(3.14)
The energy of this initial state is already affected by the static vacuum polarization. However, for
$$mr_{\text{min}}\stackrel{>}{}1$$
(3.15)
this effect is negligible (see Eq. (5.18) below), and it does not make sense to consider $`r_{\text{min}}`$ smaller than the Compton size of the vacuum particle. Then, up to a small correction, the energy of the shell (with the rest energy subtracted) retains its classical value
$$=\frac{e^2}{2r_{\text{min}}},$$
(3.16)
and so does the acceleration of the shell at $`t=t_{\text{start}}`$
$$\ddot{\rho }|_{t_{\text{start}}}=\frac{}{M}\frac{1}{r_{\text{min}}}.$$
(3.17)
Since the shell moves with acceleration, it creates particles from the vacuum provided that its typical frequency exceeds the threshold of pair creation: $`\mathrm{}\nu >2mc^2`$. At the high-frequency limit $`\mathrm{}\nu mc^2`$ its vacuum radiation stops depending on the mass $`m`$ . As seen from Eq. (3.17), the typical frequency $`\nu `$ is proportional to $`/M`$ :
$$\nu =\frac{}{M}\frac{1}{r_{\text{min}}}.$$
(3.18)
The bigger the ratio $`/M`$ , the bigger is the acceleration at $`t_{\text{start}}`$ , and the more violent is the creation of particles. Therefore, it is interesting to consider the ultrarelativistic shell $`(/M)1`$. The latter condition can be enhanced to provide for the high-frequency regime:
$$\frac{}{M}mr_{\text{min}}.$$
(3.19)
At the same time, under condition (3.15) the shell does not probe the small scales where the present description may break down. Assuming both Eqs. (3.15) and (3.19) one switches over from the consideration of the static vacuum polarization to studying the vacuum reaction on a rapidly moving source creating pairs. This is the purpose of the present work.
Without predetermining the law of motion $`\rho (t)`$ one may assume that, beginning with $`t=t_{\text{start}}`$ , the shell expands monotonically with an increasing velocity $`\dot{\rho }(t)`$ which at $`t=\mathrm{}`$ reaches some finite value $`\dot{\rho }(\mathrm{})`$ . Then $`\dot{\rho }(\mathrm{})`$ may serve as a measure at late time of the acceleration at $`t_{\text{start}}`$ . The world line of the shell is shown in Fig. 1. As $`(/M)\mathrm{}`$ , the velocity $`\dot{\rho }(t)`$ approaches $`1`$ at all $`t`$ except in a small sector near $`t=t_{\text{start}}`$ . The world line of the shell approaches then the broken line $`N`$ in Fig. 1. These assumptions are valid for the classical motion of the shell
$$\frac{M}{\sqrt{1\dot{\rho }^2}}+\frac{1}{2}\frac{e^2}{\rho }=M+,$$
(3.20)
and they cannot be invalidated by the quantum corrections if the coupling constant $`\kappa ^2`$ is small.
## 4 The retarded resolvent
As in Refs. , it is convenient to express the resolvent in the expectation-value equations through the operator
$$_q=\sqrt{\frac{2q}{\mathrm{}}}K_1(\sqrt{2q\mathrm{}}),q<0$$
(4.1)
depending on the parameter $`q`$, whose retarded kernel is of the form
$$_qX(x)=\frac{1}{4\pi }\underset{\text{past of}x}{}𝑑\overline{x}\overline{g}^{1/2}\delta \left(\sigma (x,\overline{x})q\right)X(\overline{x}).$$
(4.2)
Here $`K_1`$ is the order-1 Macdonald function, and $`\sigma (x,\overline{x})`$ is the world function: one half of the square of the geodetic distance between the points $`x`$ and $`\overline{x}`$ . The integration in Eq. (4.2) is over the past sheet of the hyperboloid of equal geodetic distance $`\sqrt{2q}`$ from the observation point $`x`$.
The needed expression is provided by the formula
$$\frac{1}{\mu ^2\mathrm{}}=\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})K_0(\sqrt{2q\mathrm{}})$$
(4.3)
involving the Macdonald and Bessel functions, and the result is the following expression for the kernel of the retarded resolvent:
$$\frac{1}{\mu ^2\mathrm{}}X=\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\frac{d}{dq}_qX$$
(4.4)
with $`_qX`$ in Eq. (4.2). If the test function is a tensor
$$X=X^{\mu _1\mathrm{}\mu _n},$$
(4.5)
Eq. (4.2) is an abbreviation of
$$_qX^{\mu _1\mathrm{}\mu _n}(x)=\frac{1}{4\pi }\underset{\text{past of}x}{}𝑑\overline{x}\overline{g}^{1/2}\delta \left(\sigma (x,\overline{x})q\right)g_{\overline{\mu }_1}^{\mu _1}(x,\overline{x})\mathrm{}g_{\overline{\mu }_n}^{\mu _n}(x,\overline{x})X^{\overline{\mu }_1\mathrm{}\overline{\mu }_n}(\overline{x})$$
(4.6)
where $`g_{\overline{\mu }}^\mu (x,\overline{x})`$ is the propagator of the geodetic parallel transport .
As an example of the use of Eq. (4.4) one may derive the retarded kernel of the massless operator $`1/\mathrm{}`$ . For a test function $`X`$ asymptotically static in the past, the operator $`_q`$ decreases as $`q\mathrm{}`$ :
$$_qX|_q\mathrm{}\frac{1}{\sqrt{q}}.$$
(4.7)
Thus one obtains
$$\frac{1}{\mathrm{}}X(x)=_qX|_{q=0}=\frac{1}{4\pi }\underset{\text{past of}x}{}d\overline{x}\overline{g}^{1/2}\delta \left(\sigma (x,\overline{x})\right)X(\overline{x}).$$
(4.8)
When $`X`$ is a spherically symmetric scalar $`X=X(t,r)`$, and the coordinates (2.3) are used, one has
$$_qX(x)=\frac{1}{2}\underset{\text{past of}x}{}𝑑\overline{t}𝑑\overline{r}\overline{r}^2\underset{1}{\overset{1}{}}d(\mathrm{cos}\omega )\delta (\sigma q)X(\overline{t},\overline{r}),$$
(4.9)
$$2\sigma =(t\overline{t})^2+(r\overline{r})^2+2r\overline{r}(1\mathrm{cos}\omega )$$
(4.10)
where $`\omega `$ is the arc length between the points $`(\theta ,\phi )`$ and $`(\overline{\theta },\overline{\phi })`$ on the unit two-sphere. Denote $`{}_{}{}^{2}\sigma `$ the world function of the two-dimensional Lorentzian section:
$${}_{}{}^{2}\sigma =\frac{1}{2}(t\overline{t})^2+\frac{1}{2}(r\overline{r})^2.$$
(4.11)
It follows from Eq. (4.10) that, on the hyperboloid $`\sigma =q`$, the range $`1<\mathrm{cos}\omega <1`$ is equivalent to the following range of variation of $`{}_{}{}^{2}\sigma `$ :
$$q2r\overline{r}<{}_{}{}^{2}\sigma <q.$$
(4.12)
Therefore, the result of the angle integration in Eq. (4.9) is
$$_qX(x)=\frac{1}{2r}\underset{\text{past of}x}{}𝑑\overline{t}𝑑\overline{r}\overline{r}X(\overline{t},\overline{r})\theta (q{}_{}{}^{2}\sigma )\theta ({}_{}{}^{2}\sigma +2r\overline{r}q),$$
(4.13)
$$\frac{d}{dq}_qX(x)=\frac{1}{2r}\underset{\text{past of}x}{}𝑑\overline{t}𝑑\overline{r}\overline{r}X(\overline{t},\overline{r})\left[\delta ({}_{}{}^{2}\sigma q)\delta ({}_{}{}^{2}\sigma +2r\overline{r}q)\right].$$
(4.14)
In the past of the observation point $`x`$, the boundaries specified by the $`\theta `$-functions in Eq. (4.13) are of the form
$`{}_{}{}^{2}\sigma q=0:`$ $`\overline{t}=t\sqrt{(r\overline{r})^22q},`$ (4.15)
$`{}_{}{}^{2}\sigma +2r\overline{r}q=0:`$ $`\overline{t}=t\sqrt{(r+\overline{r})^22q},`$ (4.16)
and Eq. (4.13) can be rewritten as
$$_qX=\frac{1}{2r}\underset{0}{\overset{\mathrm{}}{}}𝑑\overline{r}\overline{r}\underset{t\sqrt{\left(r+\overline{r}\right)^22q}}{\overset{t\sqrt{\left(r\overline{r}\right)^22q}}{}}𝑑\overline{t}X(\overline{t},\overline{r}).$$
(4.17)
Eq. (4.17) yields a simple expression in the case where the source $`X`$ is static: $`X(t,r)=X(r)`$. In this case one obtains
$$_qX=\frac{1}{2r}\underset{0}{\overset{\mathrm{}}{}}𝑑\overline{r}\overline{r}X(\overline{r})\left(\sqrt{(r+\overline{r})^22q}\sqrt{(r\overline{r})^22q}\right)$$
(4.18)
and
$$\frac{1}{\mu ^2\mathrm{}}X=\frac{1}{2\mu r}\underset{0}{\overset{\mathrm{}}{}}𝑑\overline{r}\overline{r}X(\overline{r})\left[\mathrm{exp}\left(\mu |r\overline{r}|\right)\mathrm{exp}\left(\mu (r+\overline{r})\right)\right]$$
(4.19)
where use is made of the integral
$$\underset{\mathrm{}}{\overset{0}{}}𝑑q\frac{J_0(\mu \sqrt{2q})}{\sqrt{a^22q}}=\frac{1}{\mu }\mathrm{exp}\left(\mu |a|\right).$$
(4.20)
Of course, for a static $`X`$, expression (4.18) could be obtained simpler by integrating in Eq. (4.9) first over $`\overline{t}`$
$$_qX=\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}𝑑\overline{r}\overline{r}^2\underset{1}{\overset{1}{}}d(\mathrm{cos}\omega )\frac{X(\overline{r})}{\sqrt{(r\overline{r})^2+2r\overline{r}(1\mathrm{cos}\omega )2q}}$$
(4.21)
and next over $`\mathrm{cos}\omega `$ .
## 5 The static vacuum polarization
The propagator of parallel transport $`g_{\overline{\mu }}^\mu (x,\overline{x})`$ for the metric (2.3) projects on the basis vectors as follows:
$`_\mu tg_{\overline{\mu }}^\mu (x,\overline{x})`$ $`=`$ $`\overline{}_{\overline{\mu }}\overline{t},`$ (5.1)
$`_\mu rg_{\overline{\mu }}^\mu (x,\overline{x})`$ $`=`$ $`\mathrm{cos}\omega \overline{}_{\overline{\mu }}\overline{r}+\overline{r}\overline{}_{\overline{\mu }}\mathrm{cos}\omega .`$ (5.2)
One can choose any of the two projections to convert Eq. (2.11) into a scalar equation. Specifically, one can use the fact that, owing to Eq. (5.1), the operation of projecting on $`t`$ commutes with the action of any nonlocal form factor. Hence
$$_\alpha tj^\alpha =_\alpha tj_{\text{bare}}^\alpha \gamma (\mathrm{})_\alpha tj_{\text{bare}}^\alpha ,$$
(5.3)
and by Eqs. (3.5) and (3.11)
$$_\alpha tj^\alpha =\frac{1}{4\pi r^2}\frac{}{r}𝐞(t,r),$$
(5.4)
$$_\alpha tj_{\text{bare}}^\alpha =\frac{e}{4\pi r^2}\delta \left(r\rho (t)\right).$$
(5.5)
Strictly outside and inside the shell, the function (5.5) vanishes. Therefore, in each of these regions, the local terms on the right-hand side of Eq. (5.3), i.e., the terms in which the current $`j_{\text{bare}}^\alpha `$ appears at the observation point can be omitted. Specifically, the subtraction term in the spectral formula (2.18) can be omitted. As a result, one obtains the equation
$$\frac{1}{r^2}\frac{}{r}𝐞(t,r)=\frac{\kappa ^2}{6}\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)\frac{1}{\mu ^2\mathrm{}}\left(_\alpha tj_{\text{bare}}^\alpha \right)$$
(5.6)
which is valid separately in two regions for the point $`(r,t)`$ : outside and inside the shell. Below, the notation $`\epsilon `$ is used for the function
$$\epsilon (t,r)=\theta \left(r\rho (t)\right)\theta \left(\rho (t)r\right).$$
(5.7)
The broken lines in Fig. 1 bound the future light cone of the point of start. Denote $`P`$ (for past) the exterior of this cone. By causality, the region $`P`$ can be affected only by the static sector of the evolution of the shell. Therefore, when calculating $`𝐞(t,r)`$ for the point $`(r,t)`$ in $`P`$, the law of motion of the shell can be taken $`\rho (t)=r_{\text{min}}`$ . With this law, Eqs. (4.19) and (5.5) yield straight away
$`{\displaystyle \frac{1}{\mu ^2\mathrm{}}}\left(_\alpha tj_{\text{bare}}^\alpha \right)`$ $`=`$ $`{\displaystyle \frac{e}{r}}{\displaystyle \frac{1}{8\pi \mu r_{\text{min}}}}\left[\mathrm{exp}\left(\mu |rr_{\text{min}}|\right)\mathrm{exp}\left(\mu (r+r_{\text{min}})\right)\right],`$
$`(r,t)P`$
and one obtains
$`{\displaystyle \frac{}{r}}𝐞(t,r)`$ $`=`$ $`e{\displaystyle \frac{r}{r_{\text{min}}}}{\displaystyle \frac{\kappa ^2}{24\pi }}{\displaystyle \underset{2m}{\overset{\mathrm{}}{}}}𝑑\mu \mathrm{\Gamma }(\mu ^2)\left[\mathrm{exp}\left(\mu |rr_{\text{min}}|\right)\mathrm{exp}\left(\mu (r+r_{\text{min}})\right)\right],`$
$`(r,t)P.`$
Here the expression in the square brackets is positive. Therefore, in view of the condition $`\mathrm{\Gamma }(\mu ^2)0`$, one has
$$\frac{1}{e}\frac{}{r}𝐞(t,r)<0$$
(5.10)
both outside and inside the shell.
Eqs. (3.8) and (3.9) appear now in the role of boundary conditions for the regions inside and outside the shell respectively, and in both regions they fix the solution of Eq. (5.9). The solution for $`(r,t)P`$ is
$`𝐞(t,r)`$ $`=`$ $`e{\displaystyle \frac{1+\epsilon }{2}}+{\displaystyle \frac{e}{r_{\text{min}}}}{\displaystyle \frac{\kappa ^2}{24\pi }}{\displaystyle \underset{2m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\mu }{\mu ^2}}\mathrm{\Gamma }(\mu ^2)`$
$`\times \left[(1+\epsilon \mu r)\mathrm{exp}\left(\mu \epsilon (rr_{\text{min}})\right)(1+\mu r)\mathrm{exp}\left(\mu (r+r_{\text{min}})\right)\right]`$
and is, of course, static.
As shown below, a number of features of the static solution above persists also beyond $`P`$, i.e., for a moving shell. These features are summarized in Fig. 2. First of all, inside the shell there is charge, and this charge is negative. This is a consequence of the positivity of the spectral-mass function and the nonlocal nature of the expectation-value equations. The retarded nonlocal form factor collects the bare charge from the whole interior of the past light cone of the observation point. Since the world line of the shell crosses this cone at any location of the observation point, the vacuum inside the shell gets polarized.
On the other hand, the total charge inside every sphere surrounding the shell is positive. This occurs owing to a jump of $`𝐞(t,r)`$ across the shell, i.e., owing to the positive charge of the shell itself. The jump is, however, infinite, and so are the values of $`𝐞(t,r)`$ on both sides of the shell. This infinity, of different signs inside and outside, develops in the Compton neighbourhood of the shell. At a large distance from the shell, the polarization falls off exponentially owing to the presence of the threshold $`\mu 2m`$ in the spectral integral. Qualitatively, at each given $`t`$, the $`𝐞(t,r)`$ for a moving shell has a similar shape.
The mechanism by which the singularity on the shell’s surface emerges is noteworthy since it appears repeatedly in the consideration below. This mechanism is connected with the convergence of the spectral integral in Eq. (5.11) at the upper limit. The integrand in Eq. (5.11) provides an exponential cut off at large spectral mass but only for $`rr_{\text{min}}`$. At $`r=r_{\text{min}}`$ , in view of the condition $`\mathrm{\Gamma }(\mathrm{})=1`$, the integral becomes logarithmically divergent. This is none other than the ultraviolet divergence of the charge renormalization. For an observer at infinity, the shell appears as an electric monopole screened by the vacuum. With $`𝐞(t,\mathrm{})`$ normalized as in Eq. (3.9), the unscreened monopole
$$𝐞_+=𝐞(t,r_{\text{min}}+0)(tt_{\text{start}})$$
(5.12)
should be infinite.
However, in the present case the source is not a pointlike object. The total charge inside the shell
$$𝐞_{}=𝐞(t,r_{\text{min}}0)(tt_{\text{start}})$$
(5.13)
is also infinite and has the opposite sign. Owing to this fact, the force moving the shell is finite. Indeed, with Eq. (5.11) one is able to calculate the acceleration of the shell at $`t=t_{\text{start}}`$
$$\ddot{\rho }|_{t_{\text{start}}}=\frac{e}{Mr_{\text{min}}^2}\frac{𝐞_++𝐞_{}}{2}.$$
(5.14)
Making the sum $`𝐞_++𝐞_{}`$ in the spectral integral one obtains unambiguously
$$\frac{𝐞_++𝐞_{}}{2}=\frac{e}{2}+\frac{e}{r_{\text{min}}}\frac{\kappa ^2}{24\pi }\underset{2m}{\overset{\mathrm{}}{}}\frac{d\mu }{\mu ^2}\mathrm{\Gamma }(\mu ^2)\left[1(1+\mu r_{\text{min}})\mathrm{exp}(2\mu r_{\text{min}})\right].$$
(5.15)
This way of subtracting infinities is physically equivalent to giving the shell a Compton width (see Section 9 for a refinement of this point).
The function in the square brackets in Eq. (5.15)
$$f(\mu r_{\text{min}})=1(1+\mu r_{\text{min}})\mathrm{exp}(2\mu r_{\text{min}})$$
(5.16)
is positive since
$$\frac{d}{dx}f(x)>0\text{for}x>0$$
(5.17)
and $`f(0)=0`$. Therefore, the force in Eq. (5.14) is in all cases repulsive. In the case $`mr_{\text{min}}1`$ it even acquires an extra amplifying factor $`|\mathrm{log}mr_{\text{min}}|`$. However, under condition (3.15) this force differs negligibly from its classical value:
$$\frac{𝐞_++𝐞_{}}{2}=\frac{e}{2}+e\frac{\kappa ^2}{24\pi }\left[\frac{1}{mr_{\text{min}}}\underset{2}{\overset{\mathrm{}}{}}\frac{dx}{x^2}\mathrm{\Gamma }(m^2x^2)+O\left(\mathrm{exp}(2mr_{\text{min}})\right)\right].$$
(5.18)
Also, the charge inside the shell is then concentrated almost entirely in the Compton neighbourhood of the shell, and so is the excess of charge over $`e`$ outside the shell. In this way the correspondence principle is fulfilled. On the other hand, no large scales or low energies can save one from the development of the singularity within the Compton neighbourhood of the shell. Its appearance may be understood as a signal that a charge cannot be localized more accurately than within a Compton neighbourhood. The charges of the shell immersed in the real vacuum are always annihilated and created anew in a slightly different place. As a result, the shell gets smeared to a Compton width. In this way the quantum uncertainty manifests itself.
For the sake of comparison consider also a pointlike source. This is the limiting case of the charged ball
$`𝐞_{\text{bare}}(t,r)=\{\begin{array}{ccc}e{\displaystyle \frac{r^3}{r_0^3}}& ,& r<r_0\\ e& ,& r>r_0\end{array}`$ (5.21)
as $`r_00`$. In this case, assuming that the observation point is outside the ball, Eqs. (4.19) and (3.10) yield
$$\frac{1}{\mu ^2\mathrm{}}\left(_\alpha tj_{\text{bare}}^\alpha \right)=\frac{e}{4\pi r}\mathrm{exp}(\mu r).$$
(5.22)
With the boundary condition (3.9) one then obtains from Eq. (5.6)
$$𝐞(t,r)=e+e\frac{\kappa ^2}{12\pi }\underset{2m}{\overset{\mathrm{}}{}}\frac{d\mu }{\mu }\mathrm{\Gamma }(\mu ^2)(1+\mu r)\mathrm{exp}(\mu r).$$
(5.23)
The electric field (3.7) with this $`𝐞(t,r)`$ can be written down as
$$E=\frac{}{r}U,$$
(5.24)
$$U=\frac{e}{r}\left(1+\frac{\kappa ^2}{12\pi }\underset{2m}{\overset{\mathrm{}}{}}\frac{d\mu }{\mu }\mathrm{\Gamma }(\mu ^2)\mathrm{exp}(\mu r)\right).$$
(5.25)
With the spectral-mass function in Eq. (2.20) and $`\kappa ^2`$ in Eq. (2.21), this reproduces the textbook result for the ”modified Coulomb law”. In fact, the Coulomb law is not modified as seen from Eq. (3.7). What gets modified is the charge distribution. The pointlike charge induces the same infinite screening as the charged shell does.
## 6 Solution for the moving shell
Consider Eq. (4.13) in which $`X(t,r)`$ is identified with the charge density (5.5). The lines (4.15) and (4.16) on the $`\overline{r},\overline{t}`$ plane bound the mapping on this plane of the past (sheet of the) hyperboloid $`\sigma (x,\overline{x})=q`$ of the observation point $`x`$. The observation point has the coordinates $`r,t`$ and is shown in Fig. 3 along with the two boundaries of the mapping of its past hyperboloid (the bold lines). At $`q=0`$ the past hyperboloid becomes the past light cone whose mapping on the $`\overline{r},\overline{t}`$ plane is bounded by the light lines in Fig. 3. The world line of the shell $`\overline{r}=\rho (\overline{t})`$ which is also shown in Fig. 3 crosses the upper boundary of the past hyperboloid at some point $`r_+,t_+`$ and the lower boundary at some point $`r_{},t_{}`$. These points are determined by the equations
$`\{\begin{array}{ccc}\hfill r_+& =& \rho (t_+),\hfill \\ \hfill t_+& =& t\sqrt{(rr_+)^22q},\hfill \end{array}`$ (6.3)
$`\{\begin{array}{ccc}\hfill r_{}& =& \rho (t_{}),\hfill \\ \hfill t_{}& =& t\sqrt{(r+r_{})^22q},\hfill \end{array}`$ (6.6)
and their locations on the world line of the shell depend on the location of the observation point $`r,t`$ and on the value of $`q`$. At $`q=0`$ the coordinates solving Eqs. (6.1) and (6.2) will be denoted
$$r_+^0,t_+^0,r_{}^0,t_{}^0$$
(6.7)
respectively.
With the notation above, Eqs. (4.13) and (5.5) yield
$$_q\left(_\alpha tj_{\text{bare}}^\alpha \right)=\frac{e}{8\pi r}\underset{t_{}}{\overset{t_+}{}}\frac{dt}{\rho (t)}.$$
(6.8)
Using Eqs. (6.1), (6.2) one can calculate
$$\frac{t_\pm }{q}\frac{1}{A_\pm },$$
(6.9)
$$A_+=(tt_+)(rr_+)\dot{\rho }(t_+),$$
(6.10)
$$A_{}=(tt_{})+(r+r_{})\dot{\rho }(t_{}),$$
(6.11)
and hence
$$\frac{d}{dq}_q\left(_\alpha tj_{\text{bare}}^\alpha \right)=\frac{e}{8\pi r}\left(\frac{1}{r_+A_+}\frac{1}{r_{}A_{}}\right).$$
(6.12)
With this expression, Eq. (4.4) yields
$$\frac{1}{\mu ^2\mathrm{}}\left(_\alpha tj_{\text{bare}}^\alpha \right)=\frac{e}{8\pi r}\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\left(\frac{1}{r_+A_+}\frac{1}{r_{}A_{}}\right),$$
(6.13)
and then from Eq. (5.6) one obtains
$$\frac{}{r}𝐞(t,r)=r\frac{e\kappa ^2}{48\pi }\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\left(\frac{1}{r_+A_+}\frac{1}{r_{}A_{}}\right).$$
(6.14)
Eq. (6.10) should now be integrated over $`r`$ along the line $`t=\text{const.}`$ with the boundary conditions (3.8) and (3.9) inside and outside the shell respectively. This integration can be done explicitly, and the final result is
$$𝐞(t,r)=e\frac{1+\epsilon }{2}+e\frac{\kappa ^2}{24\pi }\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)w(\mu ,t,r),$$
(6.15)
$$w(\mu ,t,r)=\frac{1}{\mu ^2}\frac{1+\epsilon }{2}+\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})F(q,t,r)$$
(6.16)
with $`\epsilon `$ in Eq. (5.7), and
$$F(q,t,r)=\frac{1}{2}\left[\underset{t_{}}{\overset{t_+}{}}\frac{dt}{\rho (t)}+\mathrm{log}\left((tt_{})(r+r_{})\right)\mathrm{log}\left((tt_+)+(rr_+)\right)\right].$$
(6.17)
It follows from Eqs. (6.1), (6.2) that the arguments of both $`\mathrm{log}`$’s in Eq. (6.13) are nonnegative.
For the proof of the result above, first use Eqs. (6.1), (6.2) to show that the derivative of the function (6.13) is
$$\frac{}{r}F(q,t,r)=\frac{r}{2}\left(\frac{1}{r_+A_+}\frac{1}{r_{}A_{}}\right),$$
(6.18)
and thereby expression (6.11) satisfies the equation (6.10). Next note that at $`r=0`$ the points $`r_+,t_+`$ and $`r_{},t_{}`$ coincide. Hence
$$F(q,t,0)=0,$$
(6.19)
and thereby expression (6.11) with $`\epsilon =1`$ satisfies the boundary condition (3.8).
Finally, consider expression (6.11) with $`\epsilon =+1`$ at the limit where the observation point $`r,t`$ moves to spatial infinity: $`r\mathrm{}`$ at a fixed $`t`$. At any $`t`$ and a sufficiently large $`r`$, the point $`r,t`$ will enter the region $`P`$ which is affected only by the static sector of the evolution of the shell. Indeed, as the observation point moves to spatial infinity, both points $`r_+,t_+`$ and $`r_{},t_{}`$ shift to the past along the world line of the shell and turn out to be on its static sector. Therefore,
$$F(q,t,r)|_r\mathrm{}=F_{\text{stat}}(q,r)|_r\mathrm{}$$
(6.20)
where
$`F_{\text{stat}}(q,r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}[{\displaystyle \frac{1}{r_{\text{min}}}}(\sqrt{(r+r_{\text{min}})^22q}\sqrt{(rr_{\text{min}})^22q})`$ (6.21)
$`+\mathrm{log}\left(\sqrt{(r+r_{\text{min}})^22q}(r+r_{\text{min}})\right)`$
$`\mathrm{log}(\sqrt{(rr_{\text{min}})^22q}+(rr_{\text{min}}))].`$
With this expression the integral in Eq. (6.12) can be calculated:
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑qJ_0(\mu \sqrt{2q})F_{\text{stat}}(q,r)={\displaystyle \frac{1}{\mu ^2}}{\displaystyle \frac{1+\epsilon _0}{2}}+{\displaystyle \frac{1}{2\mu ^3r_{\text{min}}}}`$
$`\times \left[(1+\epsilon _0\mu r)\mathrm{exp}\left(\mu \epsilon _0(rr_{\text{min}})\right)(1+\mu r)\mathrm{exp}\left(\mu (r+r_{\text{min}})\right)\right],`$ (6.22)
$$\epsilon _0=\theta (rr_{\text{min}})\theta (r_{\text{min}}r),$$
(6.23)
and it is seen why the explicit term in $`1/\mu ^2`$ is introduced in Eq. (6.12). Here use is made of the integrals
$$\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\left(\sqrt{a^22q}\sqrt{2q}\right)=\frac{1}{\mu ^3}\left[1(1+\mu |a|)\mathrm{exp}(\mu |a|)\right],$$
(6.24)
$$\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\mathrm{log}\left(\frac{\sqrt{a^22q}\pm |a|}{\sqrt{2q}}\right)=\pm \frac{1}{\mu ^2}\left(1\mathrm{exp}(\mu |a|)\right),$$
(6.25)
and the result agrees with Eq. (5.11). Thus one obtains from Eqs. (6.11) and (6.16)
$`𝐞(t,r)|_r\mathrm{}=e+{\displaystyle \frac{e}{r_{\text{min}}}}{\displaystyle \frac{\kappa ^2}{24\pi }}{\displaystyle \underset{2m}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\mu }{\mu ^2}}\mathrm{\Gamma }(\mu ^2)(\mathrm{exp}(\mu r_{\text{min}})\mathrm{exp}(\mu r_{\text{min}}))`$
$`\times (1+\mu r)\mathrm{exp}(\mu r)|_r\mathrm{}`$ (6.26)
whence
$$𝐞(t,r)|_r\mathrm{}=e+O\left(\mathrm{exp}(2mr)\right)$$
(6.27)
and thereby the boundary condition (3.9) is satisfied.
## 7 Convergence of the spectral integral
The integral (6.12) in $`q`$ always converges. Indeed, when the observation point $`r,t`$ is fixed, and $`q\mathrm{}`$, the points $`r_+,t_+`$ and $`r_{},t_{}`$ again shift to the past and turn out to be on the static sector of the world line of the shell. Therefore,
$$F(q,t,r)|_q\mathrm{}=F_{\text{stat}}(q,r)|_q\mathrm{}$$
(7.1)
whence
$$F(q,t,r)|_q\mathrm{}=\frac{rr_{\text{min}}^2}{(2q)^{3/2}}.$$
(7.2)
Since also
$$F(q,t,r)|_{q0}=O\left(\mathrm{log}(2q)\right)$$
(7.3)
as discussed below, the integral in Eq. (6.12) converges even at $`\mu =0`$ and at any location of the observation point $`r,t`$.
The behaviour (7.3) is in all cases calculable directly from Eqs. (6.13) and (6.1), (6.2) but its coefficient is different for different locations of the observation point. For the observation point outside the shell one obtains
$$F(q,t,r)|_{r>\rho (t)}=\frac{1}{2}\mathrm{log}(2q)+\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}a_n(t,r)(2q)^n,q0$$
(7.4)
with
$$a_0(t,r)=\mathrm{log}[4(rr_+^0)(r+r_{}^0)]+\underset{t_{}^0}{\overset{t_+^0}{}}\frac{dt}{\rho (t)}$$
(7.5)
whereas, for the observation point inside the shell,
$$F(q,t,r)|_{r<\rho (t)}=\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}b_n(t,r)(2q)^n,q0$$
(7.6)
with
$$b_0(t,r)=\mathrm{log}\frac{(r_+^0r)}{(r_{}^0+r)}+\underset{t_{}^0}{\overset{t_+^0}{}}\frac{dt}{\rho (t)}.$$
(7.7)
The result for the observation point on the shell is different from both Eqs. (7.4) and (7.6):
$$F(q,t,r)|_{r=\rho (t)}=\frac{1}{4}\mathrm{log}(2q)+\frac{1}{2}\underset{n=0}{\overset{\mathrm{}}{}}c_n(t)(2q)^{n/2},q0.$$
(7.8)
Here the expansion is in half-integer powers with
$$c_0(t)=\mathrm{log}2[\rho (t)+\rho (t_{}^0)]+\mathrm{log}\sqrt{\frac{1\dot{\rho }(t)}{1+\dot{\rho }(t)}}+\underset{t_{}^0}{\overset{t}{}}\frac{d\overline{t}}{\rho (\overline{t})},$$
(7.9)
$$c_1(t)=\frac{1}{\sqrt{1\dot{\rho }^2(t)}}\left(\frac{1}{\rho (t)}+\frac{1}{2}\frac{\ddot{\rho }(t)}{1\dot{\rho }^2(t)}\right),$$
(7.10)
and $`t_{}^0`$ taken at $`r=\rho (t)`$. The cause of this discontinuity is in the fact that at $`q=0`$ the upper boundary of the hyperboloid acquires a conic singularity (Fig. 3). When the observation point crosses the shell, the point $`r_+^0,t_+^0`$ passes through the vertex of the cone.
The behaviour of $`F(q,t,r)`$ at $`q0`$ determines the leading (power) behaviour of the integral (6.12) at $`\mu \mathrm{}`$. The latter behaviour can be obtained by integrating by parts with the aid of the relation
$$J_0(\mu \sqrt{2q})=\frac{2}{\mu ^2}\frac{}{q}q\frac{}{q}J_0(\mu \sqrt{2q}).$$
(7.11)
In the case (7.4) or (7.6) one may write
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}dqJ_0(\mu \sqrt{2q})F(q,t,r)=(q{\displaystyle \frac{2}{\mu ^2}}q{\displaystyle \frac{}{q}}){\displaystyle \underset{n=o}{\overset{N}{}}}\left({\displaystyle \frac{2}{\mu ^2}}\right)^n\left({\displaystyle \frac{}{q}}q{\displaystyle \frac{}{q}}\right)^nF(q,t,r)|_{q=0}`$
$`+\left({\displaystyle \frac{2}{\mu ^2}}\right)^{N+1}{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑qJ_0(\mu \sqrt{2q})\left({\displaystyle \frac{}{q}}q{\displaystyle \frac{}{q}}\right)^{N+1}F(q,t,r)`$ (7.12)
where use is made of Eq. (7.2). Since, for any $`N`$, the integral on the right-hand side of Eq. (7.12) converges and decreases as $`\mu \mathrm{}`$, one obtains
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}dqJ_0(\mu \sqrt{2q})F(q,t,r)|_{r>\rho (t)}`$ $`=`$ $`{\displaystyle \frac{1}{\mu ^2}}+O\left({\displaystyle \frac{1}{\mu ^{2N}}}\right),\mu \mathrm{}`$ (7.13)
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}dqJ_0(\mu \sqrt{2q})F(q,t,r)|_{r<\rho (t)}`$ $`=`$ $`O\left({\displaystyle \frac{1}{\mu ^{2N}}}\right),\mu \mathrm{}`$ (7.14)
where the remainder decreases faster than any power of $`1/\mu ^2`$. (It decreases exponentially, see below.)
In the case (7.8) the integration by parts as above can be done only once. One may write
$$\frac{}{q}q\frac{}{q}F(q,t,\rho (t))=\frac{1}{\sqrt{2q}}\mathrm{\Phi }(\sqrt{2q})$$
(7.15)
where $`\mathrm{\Phi }(x)`$ is analytic at $`x=0`$, and
$`{\displaystyle \underset{\mathrm{}}{\overset{0}{}}}𝑑qJ_0(\mu \sqrt{2q}){\displaystyle \frac{}{q}}q{\displaystyle \frac{}{q}}F(q,t,\rho (t))={\displaystyle \frac{1}{\mu }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑xJ_0(x)\mathrm{\Phi }\left({\displaystyle \frac{x}{\mu }}\right)={\displaystyle \frac{1}{\mu }}\left(\mathrm{\Phi }(0)+𝒪\right),`$ (7.16)
$`𝒪0,\mu \mathrm{}`$
where
$$\mathrm{\Phi }(0)=\frac{1}{4}c_1(t)$$
(7.17)
by Eq. (7.8). Hence one obtains
$$\underset{\mathrm{}}{\overset{0}{}}dqJ_0(\mu \sqrt{2q})F(q,t,r)|_{r=\rho (t)}=\frac{1}{2\mu ^2}\frac{c_1(t)}{2\mu ^3}+\frac{𝒪}{\mu ^3},\mu \mathrm{}$$
(7.18)
with $`c_1(t)`$ in Eq. (7.10).
As a result, in both regions outside and inside the shell, the behaviour of the function $`w(\mu ,t,r)`$ in Eq. (6.12) is
$$w(\mu ,t,r)|_{r\rho (t)}=O\left(\frac{1}{\mu ^{2N}}\right),N,\mu \mathrm{}$$
(7.19)
whereas, on the shell,
$$w(\mu ,t,\rho (t)\pm 0)=\pm \frac{1}{2\mu ^2}+O\left(\frac{1}{\mu ^3}\right),\mu \mathrm{}.$$
(7.20)
Recalling the condition $`\mathrm{\Gamma }(\mathrm{})=1`$, one infers that the spectral-mass integral in Eq. (6.11) converges and defines the function $`𝐞(t,r)`$ in all cases except in the case where the point $`r,t`$ is on the world line of the shell. In the latter case the spectral integral diverges logarithmically, the divergent terms having the same coefficients but different signs on the two sides of the shell. It follows that the distribution $`𝐞(t,r)`$ is singular on the shell’s surface, and the next task is obtaining the form of this singularity.
## 8 Singularity of the electric field on the shell’s surface
For obtaining the behaviour of $`𝐞(t,r)`$ on the shell’s surface, the behaviour of the function $`w(\mu ,t,r)`$ at $`\mu \mathrm{}`$ should be known including the exponentially decreasing terms. These are determined by the singularities of the function $`F(q,t,r)`$ in the complex plane of the variable $`z=\sqrt{2q}`$.
The function $`F(q,t,r)`$ will now be considered only off the shell. It is convenient first to integrate by parts
$$\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})F(q,t,r)=\frac{2}{\mu }\underset{0}{\overset{\mathrm{}}{}}𝑑\sqrt{2q}J_1(\mu \sqrt{2q})q\frac{}{q}F(q,t,r)$$
(8.1)
and next use the analyticity of $`q(F/q)`$ at $`q=0`$ to write
$$\underset{0}{\overset{\mathrm{}}{}}dxJ_1(\mu x)\left(q\frac{}{q}F\right)|_{2q=x^2}=\frac{1}{\mu }\left(q\frac{}{q}F\right)|_{q=0}+\frac{1}{2}\text{ Re }\underset{\mathrm{}}{\overset{\mathrm{}}{}}dxH_1^{(1)}(\mu x+\mathrm{i}ϵ)\left(q\frac{}{q}F\right)|_{2q=x^2}$$
(8.2)
where $`H_1^{(1)}`$ is the Hankel function. For the function $`w(\mu ,t,r)`$ both outside and inside the shell this gives
$$w(\mu ,t,r)=\frac{1}{\mu }\text{ Re }\underset{𝒞}{}dzH_1^{(1)}(\mu z)\left[q\frac{}{q}F(q,t,r)\right]|_{2q=z^2}$$
(8.3)
where the contour $`𝒞`$ passes above the real axis and closes counter-clockwise in the upper half-plane. Here
$`q{\displaystyle \frac{}{q}}F(q,t,r)`$ $`=`$ $`{\displaystyle \frac{1}{4A_+}}\left[{\displaystyle \frac{2q}{r_+}}(r_+r)\left(1\dot{\rho }^2(t_+)\right)\right]{\displaystyle \frac{1}{4}}\dot{\rho }(t_+)`$ (8.4)
$``$ $`{\displaystyle \frac{1}{4A_{}}}\left[{\displaystyle \frac{2q}{r_{}}}(r_{}+r)\left(1\dot{\rho }^2(t_{})\right)\right]+{\displaystyle \frac{1}{4}}\dot{\rho }(t_{})`$
with $`A_\pm `$ in Eqs. (6.6) and (6.7).
Of all singularities of the function (8.4) in the variable $`z=\sqrt{2q}`$, we are presently interested in the ones that have the least $`|\text{Im }z|`$ as the point $`r,t`$ approaches the shell. These are easily identified with the solutions of the equation $`A_+=0`$. Indeed, as $`r\rho (t)`$, these solutions shift to $`q=0`$ and, thereby, to $`\text{Im }z=0`$ whereas the remaining singularities stay at $`\text{Im }z0`$. This can be seen from the fact that, apart from the factor $`1/A_+`$ , expression (8.4) with $`\text{Im }z=0`$, i.e., with real $`q0`$ is nonsingular including at $`r=\rho (t)`$.
The equation
$$A_+=0$$
(8.5)
along with Eq. (6.1) determines both $`q`$ and the point $`r_+,t_+`$. Denote $`q^{}`$ the solution for $`q`$, and $`r^{},t^{}`$ the solution for $`r_+,t_+`$. The solution for $`r_+,t_+`$ proves to be real. Indeed, the point $`r^{},t^{}`$ is defined by the equations
$`tt^{}`$ $`=`$ $`(rr^{})\dot{\rho }(t^{}),`$ (8.6)
$`r^{}`$ $`=`$ $`\rho (t^{})`$ (8.7)
and is thus a point at which the world line of the shell crosses the line specified by Eq. (8.6). The latter line is shown in Fig. 4 (line $`L`$). It passes through the observation point $`r,t`$ and, at least in some neighbourhood of this point, is spacelike. This can be checked by calculating along $`L`$
$$\frac{dr^{}}{dt^{}}=\frac{1+(rr^{})\ddot{\rho }(t^{})}{\dot{\rho }(t^{})}.$$
(8.8)
It follows that, at least when the observation point is sufficiently close to the shell, the intersection at $`r^{},t^{}`$ exists and is unique (Fig. 4). The solution for $`q`$ is then real and positive:
$$2q^{}=\eta ^2,\eta =|rr^{}|\sqrt{1\dot{\rho }^2(t^{})}$$
(8.9)
whence for $`z`$ one obtains two complex conjugate solutions
$$z=\pm \mathrm{i}\eta .$$
(8.10)
Introducing a notation for the coefficient of $`1/(4A_+)`$ in Eq. (8.4), one has
$$q\frac{}{q}F(q,t,r)|_{A_+0}=\frac{1}{4A_+}(\beta |_{A_+=0}),$$
(8.11)
and one may calculate
$$\frac{A_+^2}{q}=2\alpha $$
(8.12)
with
$`\alpha `$ $`=`$ $`1\dot{\rho }^2(t_+)(r_+r)\ddot{\rho }(t_+),`$ (8.13)
$`\beta `$ $`=`$ $`{\displaystyle \frac{2q}{r_+}}(r_+r)\left(1\dot{\rho }^2(t_+)\right).`$ (8.14)
Both $`\alpha `$ and $`\beta `$ are finite and nonvanishing at $`A_+=0`$, and, moreover,
$$\alpha |_{A_+=0}>0$$
(8.15)
at least when the observation point is sufficiently close to the shell. It follows that the solutions (8.10) are branch points of the function (8.4):
$`A_+^2|_{qq^{}}`$ $`=`$ $`2\alpha |_{q=q^{}}(qq^{})+\mathrm{}`$ (8.16)
$`=`$ $`\alpha |_{q=q^{}}(z^2+\eta ^2)+\mathrm{},`$
$$q\frac{}{q}F(q,t,r)|_{A_+0}=(\frac{\beta }{4\sqrt{\alpha }}|_{A_+=0})\frac{1}{\sqrt{z^2+\eta ^2}}.$$
(8.17)
Of the two branch points, the integral (8.3) picks up the one in the upper half-plane: $`z=+\mathrm{i}\eta `$, and its contribution at large $`\mu `$ is
$$w(\mu ,t,r)\frac{1}{\mu ^2\eta }(\frac{\beta }{2\sqrt{\alpha }}|_{A_+=0})\mathrm{exp}(\mu \eta ).$$
(8.18)
The contribution (8.18) is the leading one as the observation point $`r,t`$ approaches the shell. Summarizing the calculation above, one obtains
$`w(\mu ,t,r)=\left[{\displaystyle \frac{\epsilon }{2\mu ^2}}{\displaystyle \frac{r}{r^{}}}{\displaystyle \frac{\sqrt{1\dot{\rho }^2(t^{})}}{\sqrt{1\dot{\rho }^2(t^{})(r^{}r)\ddot{\rho }(t^{})}}}+O\left({\displaystyle \frac{1}{\mu ^3}}\right)\right]\mathrm{exp}\left(\mu |r^{}r|\sqrt{1\dot{\rho }^2(t^{})}\right),`$
$`r\rho (t),\mu \mathrm{}`$ (8.19)
with $`r^{},t^{}`$ the solution of the equations (8.6), (8.7). It follows immediately from these equations that, when the observation point $`r,t`$ is on the shell, the point $`r^{},t^{}`$ coincides with $`r,t`$ (see Fig. 4). Therefore, as $`r\rho (t)`$, one may expand
$$\rho (t^{})=\rho (t)+\dot{\rho }(t)(t^{}t)+O\left(\rho (t)r\right)^2$$
(8.20)
and in this way obtain the solution
$`r^{}r={\displaystyle \frac{1}{1\dot{\rho }^2(t)}}\left(\rho (t)r\right)+O\left(\rho (t)r\right)^2,`$ (8.21)
$`t^{}t={\displaystyle \frac{\dot{\rho }(t)}{1\dot{\rho }^2(t)}}\left(\rho (t)r\right)+O\left(\rho (t)r\right)^2,`$ (8.22)
$`r\rho (t).`$
This brings Eq. (8.19) to its final form
$$w(\mu ,t,r)=\left(\frac{\epsilon }{2\mu ^2}+O\left(\frac{1}{\mu ^3}\right)\right)\mathrm{exp}\left(\mu \frac{|r\rho (t)|}{\sqrt{1\dot{\rho }^2(t)}}\right),r\rho (t),\mu \mathrm{}.$$
(8.23)
It is seen from the latter expression that, with any law of motion $`\rho (t)`$, the mechanism of formation of the singularity on the shell’s surface is one and the same. At $`r=\rho (t)`$, the integrand in Eq. (6.11) loses the exponential cut off and becomes $`O(1/\mu ^2)`$, $`\mu \mathrm{}`$. The behaviour of $`𝐞(t,r)`$ as $`r\rho (t)`$ can now be obtained by calculating the spectral-mass integral (6.11) with the function (8.23):
$$\frac{}{\eta }𝐞(t,r)=\frac{e\kappa ^2}{24\pi }\frac{\epsilon }{\eta }\underset{2m\eta }{\overset{\mathrm{}}{}}𝑑x\mathrm{\Gamma }\left(\frac{x^2}{\eta ^2}\right)\mathrm{exp}(x),\eta =\frac{|r\rho (t)|}{\sqrt{1\dot{\rho }^2(t)}}0.$$
(8.24)
The same result is obtained with Eq. (6.10). Since $`\mathrm{\Gamma }(\mathrm{})=1`$, one has
$$𝐞(t,r)|_{r\rho (t)\pm 0}=e\frac{\kappa ^2}{24\pi }\mathrm{log}|r\rho (t)|.$$
(8.25)
This behaviour is shown in Fig. 2. The electric field, as given in Eq. (3.7), differs from $`𝐞(t,r)`$ only by the factor $`1/r^2`$ which is finite and continuous across the shell.
## 9 The force exerted by the charged shell on itself
The singularity of the electric field on the shell’s surface does not affect the motion of the shell since it cancels in the sum
$$𝐞_+(t)+𝐞_{}(t)=𝐞(t,\rho (t)+0)+𝐞(t,\rho (t)0)$$
(9.1)
which according to Eq. (3.12) determines the force exerted by the shell on itself. Making the sum (9.1) in the spectral integral (6.11) makes this cancellation unambiguous. Indeed, at $`q<0`$ the points $`r_+,t_+`$ and $`r_{},t_{}`$ move along smooth trajectories as the observation point crosses the shell (see Fig. 3). Therefore, the function $`F(q,t,r)`$ with $`q<0`$ remains continuous and defines
$$(q,t)F(q,t,\rho (t)).$$
(9.2)
The integral
$$\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})F(q,t,r)$$
(9.3)
is also continuous. The function $`w(\mu ,t,r)`$ is discontinuous but finite and defines
$$2𝒲(\mu ,t)w(\mu ,t,\rho (t)+0)+w(\mu ,t,\rho (t)0).$$
(9.4)
Then from Eqs. (6.11), (6.12) one obtains
$$𝐞_+(t)+𝐞_{}(t)=e+e\frac{\kappa ^2}{12\pi }\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)𝒲(\mu ,t),$$
(9.5)
$$𝒲(\mu ,t)=\frac{1}{2\mu ^2}+\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})(q,t),$$
(9.6)
and the function $`(q,t)`$ is given by expression (6.13) with the insertion of $`r=\rho (t)`$.
The behaviour of the function $`(q,t)`$ at $`q\mathrm{}`$ is determined by Eq. (7.2). The behaviour of this function at $`q0`$ is obtained in Eq. (7.8), and the behaviour of the integral (9.6) with this function is obtained in Eq. (7.18). For $`𝒲(\mu ,t)`$ this yields
$$𝒲(\mu ,t)=O\left(\frac{1}{\mu ^3}\right),\mu \mathrm{}.$$
(9.7)
As a result, the spectral-mass integral (9.5) converges, and the force exerted on the shell is finite. The effect of making the sum (9.1) in the spectral integral is a cancellation of the $`1/\mu ^2`$, $`\mu \mathrm{}`$ terms (7.20) in the integrand.
Since the calculation above implies a subtraction of infinities, it should be analysed what regularization does it correspond to physically. The answer is contained in expression (8.23). Calculate $`𝐞(t,r)`$ for two close points $`r_1,t_1`$ and $`r_2,t_2`$ outside and inside the shell respectively, and consider the sum
$$𝐞(t_1,r_1)+𝐞(t_2,r_2).$$
(9.8)
This is given by the spectral integral (6.11) with
$$w(\mu ,t_1,r_1)+w(\mu ,t_2,r_2)=\frac{1}{2\mu ^2}\left[\mathrm{exp}\left(\mu \eta (t_1,r_1)\right)\mathrm{exp}\left(\mu \eta (t_2,r_2)\right)\right]+O\left(\frac{1}{\mu ^3}\right),$$
(9.9)
$$\eta (t,r)=\frac{|r\rho (t)|}{\sqrt{1\dot{\rho }^2(t)}}.$$
(9.10)
The result (9.5) is recovered at the limit where the exponents in Eq. (9.9) tend to zero. Since the scale for $`\mu `$, set up by the spectral integral, is $`m`$, the regularization implied in $`𝐞_+(t)`$ and $`𝐞_{}(t)`$ before their sum is made consists in a fixation of the parameter
$$m\eta (t,r)=\text{const. }0.$$
(9.11)
As soon as the sum (9.8) is made, the points $`r_1,t_1`$ and $`r_2,t_2`$ can be brought to the shell in any succession and along any pathes. Since the $`1/\mu ^2`$ term of expression (9.9) cancels in all cases, the limit for the sum is finite and independent of the way the regularization is removed.
The difference $`|r\rho (t)|`$ that figures in expression (9.10) is the proper distance from the point $`r,t`$ to the shell along the line $`t=\text{const}`$ . But not this distance is made fixed in Eq. (9.11). The function $`\eta (t,r)`$ is the proper distance from the point $`r,t`$ to the shell along the line orthogonal to the world line of the shell. Indeed, this function was introduced in Eq. (8.9), and originally it had the form
$$\eta (t,r)=\sqrt{(rr^{})^2(tt^{})^2}$$
(9.12)
where $`r^{},t^{}`$ is the point on the world line of the shell connected with the point $`r,t`$ by the line $`L`$ (Fig. 4). It is easy to check from Eq. (8.8) that, up to $`O(rr^{})`$, the line $`L`$ is orthogonal to the world line of the shell at the point of their intersection.
The final inference is that the subtraction of infinities in the spectral integral is physically equivalent to giving the shell a Compton width in the direction orthogonal to its world line. The lines on the $`r,t`$ plane specified by Eq. (9.11):
$$\frac{|r\rho (t)|}{\sqrt{1\dot{\rho }^2(t)}}=\frac{\text{const.}}{m}$$
(9.13)
mark the band of quantum uncertainty around the world line of the shell. This band is shown in Fig. 5, and it narrows as the speed of expansion increases.
## 10 The ultrarelativistic limit
It will now be shown that the force exerted by the shell on itself is singular at the ultrarelativistic limit. This is the limit at which the world line of the shell approaches the world line of an outgoing light ray (the line $`N`$ in Fig. 1). To be more precise, we consider a family of functions $`\rho (t)`$, for which
$$1\dot{\rho }(t)=\delta f(\delta ,t),\delta 0.$$
(10.1)
Here $`\delta `$ is a parameter (function of the initial data), and it is assumed that, at all $`t>t_{\text{start}}`$, $`f(\delta ,t)`$ has a finite limit as $`\delta 0`$, whereas, at $`t=t_{\text{start}}`$,
$$f(\delta ,t_{\text{start}})=\frac{1}{\delta }.$$
(10.2)
The function $`f`$ can be normalized as $`f(\delta ,\mathrm{})=1`$, and then
$$\delta =1\dot{\rho }(\mathrm{}).$$
(10.3)
Eq. (10.1) generalizes the form that the classical law of motion has as $`(M/)0`$. Indeed, with $``$ and $`r_{\text{min}}`$ taken for independent data, Eq. (3.20) can be written as
$$\frac{1}{\sqrt{1\dot{\rho }^2}}=\frac{}{M}\left(1\frac{r_{\text{min}}}{\rho }\right)+1.$$
(10.4)
However, Eq. (10.1) does not predetermine the dependence of the velocity on energy. The limiting form of $`\rho (t)`$ at $`\delta =0`$ is
$`\rho _{\text{lim}}(t)=\rho (t)|_{\delta =0}=\{\begin{array}{cc}r_{\text{min}},\hfill & t<t_{\text{start}},\\ r_{\text{min}}+(tt_{\text{start}}),\hfill & t>t_{\text{start}}.\end{array}`$ (10.7)
This world line is shown in Fig. 6.
Consider the function (9.2) for the shell obeying the law of motion (10.1). The line $`N`$ in Figs. 1 and 6 is the boundary of the region $`P`$ considered in Section 5. Therefore, when the point $`r,t`$ in the argument of the function $`F(q,t,r)`$ is on the line $`N`$, the respective points $`r_+,t_+`$ and $`r_{},t_{}`$ are, at all $`q`$, at the static sector of the world line of the shell. For the function (9.2) with $`\rho (t)`$ in Eq. (10.5) this yields the result
$$(q,t)|_{\delta =0}=F_{\text{stat}}(q,\rho _{\text{lim}}(t))$$
(10.8)
where the function $`F_{\text{stat}}(q,r)`$ is given in Eq. (6.17). The integral that figures in Eq. (9.6) is then already calculated in Eq. (6.18). One obtains
$$𝒲(\mu ,t)|_{\delta =0}=\frac{1}{2\mu ^2}+\frac{(1+\mu \rho _{\text{lim}}(t))}{2\mu ^3r_{\text{min}}}\{\mathrm{exp}[\mu (\rho _{\text{lim}}(t)r_{\text{min}})]\mathrm{exp}[\mu (\rho _{\text{lim}}(t)+r_{\text{min}})]\}.$$
(10.9)
For $`tt_{\text{start}}`$ this brings one back to Eq. (5.15), but for $`t>t_{\text{start}}`$ one has
$$𝒲(\mu ,t)|_{\delta =0}=\frac{1}{2\mu ^2}(1+𝒪),𝒪0,\mu \mathrm{}$$
(10.10)
and the integral in Eq. (9.5) diverges at large $`\mu `$ :
$$[𝐞_+(t)+𝐞_{}(t)]|_{\delta =0}=\mathrm{},t>t_{\text{start}}.$$
(10.11)
The force exerted on the shell is infinite at the ultrarelativistic limit.
The null limit (10.5) for $`\rho (t)`$ is never reached exactly even with the classical motion, and the next task is obtaining the asymptotic behaviour of the force as $`\rho (t)`$ approaches the null limit. For that consider any point with a given $`t>t_{\text{start}}`$ on a timelike world line of the shell. When this point is sufficiently close to the line $`N`$, the respective point $`r_{},t_{}`$ is, at all $`q`$, at the static sector of the evolution of the shell (see Fig.3). However, for the point $`r_+,t_+`$ this is not the case. Rather the range of variation of $`q`$ should be divided into two: the one for which $`t_+<t_{\text{start}}`$ and the one for which $`t_+>t_{\text{start}}`$. It follows from Eq. (6.1) that the former range is
$$\mathrm{}<q<\frac{1}{2}s^2(t),$$
(10.12)
and the latter is
$$\frac{1}{2}s^2(t)<q<0$$
(10.13)
where
$$s(t)=\sqrt{(tt_{\text{start}})^2(\rho (t)r_{\text{min}})^2}.$$
(10.14)
Only in the range (10.10) does one have
$$(q,t)=F_{\text{stat}}(q,\rho (t)).$$
(10.15)
Expression (9.6) first integrated by parts as in Eq. (8.1):
$$𝒲(\mu ,t)=\frac{1}{2\mu ^2}\frac{2}{\mu }\underset{0}{\overset{\mathrm{}}{}}𝑑\sqrt{2q}J_1(\mu \sqrt{2q})q\frac{}{q}(q,t)$$
(10.16)
may now be written in the form
$`𝒲(\mu ,t)={\displaystyle \frac{1}{2\mu ^2}}`$ $``$ $`{\displaystyle \frac{2}{\mu }}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑\sqrt{2q}J_1(\mu \sqrt{2q})q{\displaystyle \frac{}{q}}F_{\text{stat}}(q,\rho (t))`$ (10.17)
$`+`$ $`{\displaystyle \frac{2}{\mu }}{\displaystyle \underset{0}{\overset{s\left(t\right)}{}}}𝑑\sqrt{2q}J_1(\mu \sqrt{2q})q{\displaystyle \frac{}{q}}F_{\text{stat}}(q,\rho (t))`$
$``$ $`{\displaystyle \frac{2}{\mu }}{\displaystyle \underset{0}{\overset{s\left(t\right)}{}}}𝑑\sqrt{2q}J_1(\mu \sqrt{2q})q{\displaystyle \frac{}{q}}(q,t).`$
It will be noted that $`s(t)`$ is the two-dimensional geodetic distance between the point of start and a point on the shell. Therefore, when the latter point is on the null line $`N`$, $`s(t)`$ vanishes. Indeed, the insertion of the limiting form (10.5) for $`\rho (t)`$ in Eq. (10.12) yields
$$s(t)|_{\delta =0}=0,t>t_{\text{start}}.$$
(10.18)
With $`s(t)=0`$, the last two integrals in Eq. (10.15) vanish, and one recovers the result (10.7). Thus $`s(t)`$ serves in Eq. (10.15) as a parameter of proximity of the law $`\rho (t)`$ to its ultrarelativistic limit.
For obtaining the asymptotic behaviour of $`𝒲(\mu ,t)`$ as $`s(t)0`$, rewrite the last two integrals in Eq. (10.15) as
$`{\displaystyle \frac{2}{\mu }}s(t){\displaystyle \underset{0}{\overset{1}{}}}dxJ_1(x\mu s(t))\left[q{\displaystyle \frac{}{q}}F_{\text{stat}}(q,\rho (t))\right]|_{2q=x^2s^2(t)}`$ (10.19)
$``$ $`{\displaystyle \frac{2}{\mu }}s(t){\displaystyle \underset{0}{\overset{1}{}}}dxJ_1(x\mu s(t))\left[q{\displaystyle \frac{}{q}}(q,t)\right]|_{2q=x^2s^2(t)}`$
and recall that $`𝒲(\mu ,t)`$ is needed at large $`\mu `$. The approximation of interest is, therefore,
$$s(t)0,\mu s(t)=\text{finite}.$$
(10.20)
At this limit, the behaviours of the integrals (10.17) are obtained by expanding
$`\left[q{\displaystyle \frac{}{q}}F_{\text{stat}}(q,\rho (t))\right]|_{q0}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+O(q),t>t_{\text{start}}`$ (10.21)
$`\left[q{\displaystyle \frac{}{q}}(q,t)\right]|_{q0}`$ $`=`$ $`{\displaystyle \frac{1}{4}}+O(\sqrt{2q}).`$ (10.22)
Here use is made of Eqs. (7.4) and (7.8). Using also the explicit form (10.7) of the first two terms in Eq. (10.15), one obtains finally
$$𝒲(\mu ,t)=\frac{J_0(\mu s(t))}{2\mu ^2}\left(1+𝒪\right),𝒪0,\mu \mathrm{},s(t)0.$$
(10.23)
Eq. (10.21) is the sought for asymptotic formula for the ultrarelativistic motion. Setting in it $`s(t)=0`$, one recovers the limiting behaviour (10.8) which caused the divergence of the integral in Eq. (9.5). With the function (10.21) this integral converges:
$$[𝐞_+(t)+𝐞_{}(t)]|_{s(t)0}=ee\frac{\kappa ^2}{24\pi }\underset{4m^2}{\overset{\mathrm{}}{}}\frac{d\mu ^2}{\mu ^2}\mathrm{\Gamma }(\mu ^2)J_0(\mu s(t)),$$
(10.24)
and its behaviour as $`s(t)0`$ can be found by calculating
$$\frac{}{s}\underset{4m^2}{\overset{\mathrm{}}{}}\frac{d\mu ^2}{\mu ^2}\mathrm{\Gamma }(\mu ^2)J_0(\mu s)=\frac{2}{s}\underset{2ms}{\overset{\mathrm{}}{}}𝑑x\mathrm{\Gamma }\left(\frac{x^2}{s^2}\right)J_1(x).$$
(10.25)
Since $`\mathrm{\Gamma }(\mathrm{})=1`$, one obtains
$$𝐞_+(t)+𝐞_{}(t)=e+e\frac{\kappa ^2}{12\pi }\mathrm{log}(ms(t))+\kappa ^2O(1)$$
(10.26)
where $`O(1)`$ denotes the terms that remain finite at the ultrarelativistic limit.
The $`s(t)`$ in Eq. (10.12) can be represented in the form
$$s(t)=(tt_{\text{start}})\sqrt{1\dot{\rho }^2(\stackrel{~}{t})},t>t_{\text{start}}$$
(10.27)
where $`\stackrel{~}{t}`$ is some time instant between $`t_{\text{start}}`$ and $`t`$. By Eq. (10.1),
$$\mathrm{log}\left(1\dot{\rho }^2(\stackrel{~}{t})\right)=\mathrm{log}\left(1\dot{\rho }^2(t)\right)+O(1),$$
(10.28)
and, therefore, expression (10.24) may finally be written as
$$𝐞_+(t)+𝐞_{}(t)=e+e\frac{\kappa ^2}{24\pi }\mathrm{log}\left(1\dot{\rho }^2(t)\right)+\kappa ^2O(1).$$
(10.29)
By derivation, the remainder $`O(1)`$ in Eq. (10.27) is bounded uniformly in energy but not necessarily in time. Because of Eq. (10.2), one may worry about the vicinity of $`t=t_{\text{start}}`$. However, also at $`t=t_{\text{start}}`$, expression (10.27) yields the correct result, Eq. (5.18), provided that condition (3.15) is fulfilled. Since $`\kappa ^2`$ is small, all bounded terms of order $`\kappa ^2`$ may be regarded as negligible corrections. The term $`\kappa ^2O(1)`$ in Eq. (10.27) can then be discarded for all times and energies.
## 11 Vacuum back-reaction on the motion of the shell
The expression (10.27) with the term $`\kappa ^2O(1)`$ discarded is to be inserted in Eq. (3.12). Then the equation of motion of the shell closes and takes the form
$$M\frac{d}{dt}\left(\frac{\dot{\rho }}{\sqrt{1\dot{\rho }^2}}\right)=\frac{e^2}{2\rho ^2}\left[1+\frac{\kappa ^2}{24\pi }\mathrm{log}(1\dot{\rho }^2)\right].$$
(11.1)
The last term on the right-hand side of this equation is the force of the vacuum reaction. As will be clear in the next section, this is the force of the back-reaction of a radiation produced by the charged shell in the vacuum.
The force of the vacuum back-reaction depends on the velocity. Nevertheless, the equation of motion (11.1) admits the energy integral:
$$M\underset{1}{\overset{1/\sqrt{1\dot{\rho }^2}}{}}\frac{dx}{1{\displaystyle \frac{\kappa ^2}{12\pi }}\mathrm{log}x}+\frac{1}{2}\frac{e^2}{\rho }=$$
(11.2)
which at $`\kappa ^2=0`$ goes over into the classical law (3.20). That the constant $``$ is indeed the energy of the initial state is seen from the fact that, at $`\dot{\rho }=0`$, one recovers Eq. (3.16).
There is no problem with the singularity of the integral in Eq. (11.2). It is never reached. As in Eq. (3.20), for a given energy, the velocity $`\dot{\rho }`$ reaches its maximum value at $`\rho =\mathrm{}`$ but the value is now different:
$$\underset{1}{\overset{1/\sqrt{1\dot{\rho }^2\left(\mathrm{}\right)}}{}}\frac{dx}{1{\displaystyle \frac{\kappa ^2}{12\pi }}\mathrm{log}x}=\frac{}{M}.$$
(11.3)
As in Eq. (3.20), $`\dot{\rho }(\mathrm{})`$ grows with $`/M`$ but not up to 1:
$$\dot{\rho }(\mathrm{})=1\frac{1}{2}\mathrm{exp}\left(\frac{24\pi }{\kappa ^2}\right),\frac{}{M}\mathrm{}$$
(11.4)
and this is the principal consequence of the vacuum back-reaction.
Eq. (11.2) is surprising. The coupling to the vacuum charges does not change the electric potential<sup>5</sup><sup>5</sup>5Under condition (3.15). as one could expect. It changes the kinematics of motion as relativity theory does. Furthermore, within a given type of coupling, this change is universal. It does not depend on the parameters of the source, only on the coupling constant $`\kappa ^2`$. There emerges a new kinematic bound on the velocity of a charged body. As shown below, this bound is crucial for the maintenance of the conservation laws.
## 12 Emission of charge
The singularity of $`𝐞(t,r)`$ on the shell’s surface, as calculated in Section 8, has an important feature. Namely, the coefficient of the divergent $`\mathrm{log}`$ in Eq. (8.25), or, equivalently, the coefficient of $`1/\mu ^2`$ in Eq. (7.20) is constant in time. This suggests that the singularity comes from the static contribution which is present in $`𝐞_\pm (t)`$ but cancels in the difference
$$𝐞_\pm (t_1)𝐞_\pm (t_2).$$
(12.1)
The flux of charge across the shell should, therefore, be finite.
Indeed, from Eqs. (6.11), (6.12), and (9.2) one obtains
$$𝐞_\pm (t_1)𝐞_\pm (t_2)=e\frac{\kappa ^2}{24\pi }\underset{4m^2}{\overset{\mathrm{}}{}}𝑑\mu ^2\mathrm{\Gamma }(\mu ^2)\left[w(\mu ,t_1,\rho (t_1)\pm 0)w(\mu ,t_2,\rho (t_2)\pm 0)\right],$$
(12.2)
$$w(\mu ,t_1,\rho (t_1)\pm 0)w(\mu ,t_2,\rho (t_2)\pm 0)=\underset{\mathrm{}}{\overset{0}{}}𝑑qJ_0(\mu \sqrt{2q})\left[(q,t_1)(q,t_2)\right].$$
(12.3)
It follows that, first, the quantity (12.3) is continuous across the shell:
$`w(\mu ,t_1,\rho (t_1)+0)w(\mu ,t_2,\rho (t_2)+0)`$ (12.4)
$`=`$ $`w(\mu ,t_1,\rho (t_1)0)w(\mu ,t_2,\rho (t_2)0),`$
and, therefore,
$$𝐞_+(t_1)𝐞_+(t_2)=𝐞_{}(t_1)𝐞_{}(t_2).$$
(12.5)
Second, by Eq. (7.18) the quantity (12.3) is $`O(1/\mu ^3)`$, $`\mu \mathrm{}`$, and, therefore, the difference (12.2) is finite.
For obtaining the flux of charge across the shell, no new calculation is needed. Denote
$`\mathrm{\Delta }e`$ $`=`$ $`𝐞_+(\mathrm{})𝐞_+(\mathrm{})`$ (12.6)
$`=`$ $`𝐞_{}(\mathrm{})𝐞_{}(\mathrm{}).`$
This is the charge emitted by the shell for the whole of its history. Owing to Eq. (12.5), one may write
$`\mathrm{\Delta }e`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[𝐞_+(\mathrm{})+𝐞_{}(\mathrm{})\right]`$ (12.7)
$`{\displaystyle \frac{1}{2}}\left[𝐞_+(\mathrm{})+𝐞_{}(\mathrm{})\right]`$
and thereby relate the radiation of charge to the force of its back-reaction. The latter has already been considered in Sections 9-11. For the ultrarelativistic shell one has from Eq. (10.27)
$$𝐞_+(\mathrm{})+𝐞_{}(\mathrm{})=e+e\frac{\kappa ^2}{24\pi }\mathrm{log}\left(1\dot{\rho }^2(\mathrm{})\right)+\kappa ^2O(1),$$
(12.8)
and, from Eq. (5.18),
$$𝐞_+(\mathrm{})+𝐞_{}(\mathrm{})=e+\kappa ^2O\left(\frac{1}{mr_{\text{min}}}\right).$$
(12.9)
Hence
$$\mathrm{\Delta }e=e\frac{\kappa ^2}{48\pi }\mathrm{log}\left(1\dot{\rho }^2(\mathrm{})\right)+\kappa ^2O(1).$$
(12.10)
Also the instantaneous radiation flux can readily be estimated. Let $`t_1<t_2`$ be two time instants belonging to the epoch of the rapid expansion of the shell. The amount of charge emitted by the shell for the time between $`t_1`$ and $`t_2`$ is the quantity (12.1):
$$𝐞_\pm (t_1)𝐞_\pm (t_2)=e\frac{\kappa ^2}{48\pi }\mathrm{log}\frac{1\dot{\rho }^2(t_1)}{1\dot{\rho }^2(t_2)}+\kappa ^2O(1).$$
(12.11)
From Eq. (10.1) one infers that this is a negligible amount:
$$𝐞_\pm (t_1)𝐞_\pm (t_2)=\kappa ^2O(1),t_{\text{start}}<t_1<t_2.$$
(12.12)
However, if in Eq. (12.11) one takes $`t_{\text{start}}=t_1`$, i.e., if one calculates the amount of charge emitted from the beginning of expansion by the instant $`t`$, the result will be different:
$`𝐞_\pm (t_{\text{start}})𝐞_\pm (t)`$ $`=`$ $`e{\displaystyle \frac{\kappa ^2}{48\pi }}\mathrm{log}\left(1\dot{\rho }^2(t)\right)+\kappa ^2O(1)`$ (12.13)
$`=`$ $`e{\displaystyle \frac{\kappa ^2}{48\pi }}\mathrm{log}\left(1\dot{\rho }^2(\mathrm{})\right)+\kappa ^2O(1),t_{\text{start}}<t.`$
This is easy to understand. The cause of the vacuum particle creation is the acceleration of the source. The shell radiates at a short stage of its evolution near $`t=t_{\text{start}}`$ where its acceleration is maximum . Almost all the emitted charge $`\mathrm{\Delta }e`$ is released at this stage. Therefore, up to a small correction, the quantity (12.13) is constant.
Thus the rate of emission of charge by the ultrarelativistic shell is
$$\frac{\mathrm{\Delta }e}{e}=\frac{\kappa ^2}{48\pi }\mathrm{log}\left(1\dot{\rho }^2(\mathrm{})\right),\dot{\rho }(\mathrm{})1.$$
(12.14)
The rate of emission of energy was calculated in Ref. . Generalized properly, this calculation yields the same result as in Eq. (12.14):
$$\frac{\mathrm{\Delta }}{}=\frac{\kappa ^2}{48\pi }\mathrm{log}\left(1\dot{\rho }^2(\mathrm{})\right),\dot{\rho }(\mathrm{})1.$$
(12.15)
One sees that the radiation rate grows unboundedly as the motion of the shell approaches the ultrarelativistic limit. Inserting in Eqs. (12.14) and (12.15) the $`\dot{\rho }(\mathrm{})`$ calculated from the classical law of motion (3.20), one obtains the result :
$$\frac{\mathrm{\Delta }e}{e}=\frac{\mathrm{\Delta }}{}=\frac{\kappa ^2}{24\pi }\mathrm{log}\frac{}{M},\frac{}{M}\mathrm{}$$
(12.16)
which manifestly contradicts the conservation laws. However, the result (12.16) does not take into account the back-reaction of radiation. The vacuum friction does not allow the velocity of the source to approach the speed of light closer than the limit (11.4). The insertion of the expression (11.4) in Eqs. (12.14) and (12.15) restores the conservation laws:
$$\frac{\mathrm{\Delta }e}{e}=\frac{\mathrm{\Delta }}{}=\frac{1}{2},\frac{}{M}\mathrm{}.$$
(12.17)
Up to 50% of energy and charge can be extracted from the source by raising its initial energy.
## Acknowledgments
The appearance of this article owes to a grant of the Italian Ministry for Foreign Affairs and Landau Network - Centro Volta. G.V. is supported also by the Russian Foundation for Fundamental Research (Grant No. 99-02-18107).
## References
1. G. A. Vilkovisky, Phys. Rev. Lett. 83 (1999) 2297 \[hep-th/9906241\].
2. B. S. DeWitt, Phys. Rep. C 19 (1975) 295.
3. A. G. Mirzabekian and G. A. Vilkovisky, Phys. Lett. B 414 (1997) 123; Ann. Phys. 270 (1998) 391 \[gr-qc/9803006\].
4. G. A. Vilkovisky, Phys. Rev. D 60 (1999) 065012 \[hep-th/9812233\].
5. G. A. Vilkovisky, Class. Quantum Grav. 9 (1992) 895.
6. A. O. Barvinsky, Yu. V. Gusev, G. A. Vilkovisky, and V. V. Zhytnikov, J. Math. Phys. 35 (1994) 3525.
7. A. O. Barvinsky and G. A. Vilkovisky, Nucl. Phys. B 282 (1987) 163.
8. I. E. Tamm, Foundations of Electricity Theory (Gostekhizdat, Moscow, 1956).
9. B. S. DeWitt, Dynamical Theory of Groups and Fields (Gordon and Breach, New York, 1965).
## Figure captions
* The world line of the shell on the $`r,t`$ plane. The broken lines bound the future light cone of the point of start. The broken line $`N`$ is the world line of the outgoing radial light ray.
* The function $`𝐞(t,r)`$ for a given $`t`$.
* The world line of the shell crosses the past hyperboloid of the observation point. For definiteness, the observation point is shown inside the shell.
* $`L`$ is the line specified by Eq. (8.6), and $`r^{},t^{}`$ is the point at which it crosses the world line of the shell. The observation point $`r,t`$ is shown inside the shell, and the broken lines mark its light cone.
* The band of quantum uncertainty around the world line of the shell narrows as the speed of expansion increases.
* The world line of the shell at the ultrarelativistic limit. |
warning/0001/nlin0001033.html | ar5iv | text | # Approximate renormalization for the breakup of invariant tori with three frequencies
## I Introduction
In this paper, we define an approximate renormalization scheme for Hamiltonians with three degrees of freedom in order to study the break-up of invariant tori with the frequency vector $`𝝎_0=(\tau ^2+\tau ,\tau ,1)`$ where $`\tau =2\mathrm{cos}(2\pi /7)1.2469796`$ is the root of modulus larger than one of the polynomial :
$$\tau ^3+\tau ^22\tau 1=0.$$
This approximate renormalization is defined for the following family of Hamiltonians which are quadratic in the actions $`𝑨=(A_1,A_2,A_3)`$ :
$$H(𝑨,𝝋)=H_0(𝑨)+V(𝑨,𝝋),$$
(1)
where the Hamiltonian $`H_0`$ is quadratic :
$$H_0(𝑨)=𝝎_0𝑨+\frac{1}{2}(𝛀𝑨)^2,$$
(2)
and the perturbation is described by two scalar functions of the angles $`𝝋=(\phi _1,\phi _2,\phi _3)`$ :
$$V(𝑨,𝝋)=g(𝝋)𝛀𝑨+f(𝝋).$$
(3)
For $`H_0`$, the invariant torus with frequency vector $`𝝎_0`$ is located at $`𝑨`$ such that $`𝛀𝑨=0`$ and $`𝝎_0𝑨=E`$ where $`E`$ is the total energy of the system. Since $`𝝎_0`$ satisfies a diophantine condition, the KAM theorem for Hamiltonians (1) states that for a sufficiently small and smooth perturbation $`V`$, this invariant torus persists . For a sufficiently large perturbation, this invariant torus is broken (converse KAM ).
The purpose of this paper is to construct a renormalization transformation in order to investigate the properties of invariant tori with the frequency vector $`𝝎_0`$ at criticality. The idea of the renormalization approach is to iterate a transformation in the space of Hamiltonians. For Hamiltonians that have a smooth invariant torus with frequency vector $`𝝎_0`$, the iteration should converge to a trivial fixed point. The set of Hamiltonians that have a non-smooth invariant torus with this frequency vector form a surface (called critical surface) which is invariant under the renormalization, and is expected to be a codimension one stable manifold of a non-trivial fixed set of the renormalization.
The aim is to define a renormalization as was done for Hamiltonians with two degrees of freedom in Refs. . This renormalization is a combination of a partial elimination of the perturbation (we eliminate the Fourier modes of the perturbation which are sufficiently far from resonance) and a rescaling transformation which consists of a shift of the resonances, and a rescaling of time and of the actions.
An approximate renormalization scheme for the frequency vector $`𝝎_0`$ was defined in Ref. following the renormalization defined for the spiral mean torus constructed in Ref. . The main result of Ref. was that the renormalization has a fixed point on the critical surface but it is not hyperbolic. This renormalization dynamics is structurally unstable.
We propose to improve on Ref. , just as was done for Ref. in Refs. for the spiral mean torus. The main difference with the scheme constructed in Ref. is that we use a different normalization condition for the rescaling in the actions and we include a term that was previously neglected. The main result we find is that a large part of the critical surface is the codimension-one stable manifold of a hyperbolic fixed point of this renormalization. We expect the dynamics of the exact renormalization to be the same as this approximate renormalization, at least locally.
Following the approach of Escande and Doveil for systems with two degrees of freedom, we perform two approximations in this renormalization :
$`(1)`$ A three resonance approximation : we keep only the three main Fourier modes of the perturbation $`V`$ at each step of the transformation.
$`(2)`$ A mean-value quadratic approximation : we neglect the dependence on the angles of the quadratic term in the actions.
The frequency vector $`𝝎_0`$ is an eigenvector of the matrix $`\stackrel{~}{N}`$ with the eigenvalue $`(\tau +1)^10.445`$, where $`\stackrel{~}{N}`$ is the transposed matrix of the following matrix $`N`$ with integer coefficients and determinant -1 :
$$N=\left(\begin{array}{ccc}0& 1& 1\\ 1& 1& 1\\ 0& 1& 2\end{array}\right).$$
(4)
From $`\stackrel{~}{N}`$ one generates a sequence of periodic orbits (with frequency vector $`\{𝛀_k\}`$) approximating the motion with frequency vector $`𝝎_0`$, i.e. $`𝝋(t)=𝝎_0t+𝝋_0\text{ mod }2\pi `$ :
$$𝛀_k=(p_k/r_k,q_k/r_k,1)_n\mathrm{}𝝎_0,$$
where $`p_k`$, $`q_k`$ and $`r_k`$ are determined by the following recursion relations :
$`p_{k+1}=p_k+2q_k+r_k,`$
$`q_{k+1}=p_k,`$
$`r_{k+1}=q_k+r_k,`$
and $`p_0=3`$, $`q_0=1`$, $`r_0=1`$, i.e. $`𝛀_k=\stackrel{~}{N}^k𝛀_0`$. These relations define a sequence of simultaneous rational approximations $`(p_k/r_k,q_k/r_k)`$ to the pair $`(\tau ^2+\tau ,\tau )`$.
The sequence of resonances is generated by the matrix $`N`$ from the vector $`𝝂_1=(1,0,0)`$ :
$$𝝂_k=N^{k1}𝝂_1.$$
The small denominators decrease geometrically to zero with the ratio $`(\tau +1)^10.445`$ :
$$𝝎_0𝝂_k=\tau (\tau +1)^{2k}.$$
The modes $`𝝂_k`$ are in resonance with the periodic motion with frequency vector $`𝛀_k`$ (with period $`r_k`$) :
$$𝛀_k𝝂_{k+4}=0.$$
The matrix $`N`$ can be decomposed in the following way :
$$N=LPL^2,$$
(5)
where
$$L=\left(\begin{array}{ccc}0& 0& 1\\ 1& 0& 0\\ 0& 1& 1\end{array}\right),$$
(6)
and
$$P=\left(\begin{array}{ccc}0& 0& 1\\ 0& 1& 0\\ 1& 0& 0\end{array}\right).$$
(7)
We notice that $`\text{det }L=\text{det }P=1`$ and that $`P`$ is a permutation of two elements of the basis; in particular, it satisfies $`P^2=1`$. This decomposition follows from the construction of the Farey sequence for an incommensurate vector $`𝝎^3`$ . With this decomposition of the matrix $`N`$, we construct the renormalization from two operators (one associated with $`L`$ and the other one to $`P`$) in a similar way as MacKay did for Hamiltonian systems with two degrees of freedom in Ref. in order to investigate the break-up of invariant tori with arbitrary winding ratio. The approximate renormalization transformation we define for $`𝝎_0`$ can be constructed for a more general frequency vector from these two operators given the Farey decomposition of the frequency vector. For instance, concerning the spiral mean torus , the renormalization is equal to one of the operators (the one which is associated with $`L`$).
## II Definition of the renormalization transformation
The renormalization transformation $``$ we define for the frequency vector $`𝝎_0`$ is composed of two operators : one associated with the matrix $`P`$ and the other one with $`L`$. We denote $`_P`$ and $`_L`$ these operators. The transformation $``$ is defined for a fixed frequency vector $`𝝎_0`$ and is given by
$$=_L_L_P_L.$$
### A Definition of $`_P`$
The renormalization operator $`_P`$ acts on the following Hamiltonians :
$$H(𝑨,𝝋)=H_0(𝑨)+h(𝛀𝑨,𝝋),$$
(8)
where
$`H_0(𝑨)=𝝎𝑨+{\displaystyle \frac{1}{2}}(𝛀𝑨)^2,`$
$`h(𝛀𝑨,𝝋)={\displaystyle \underset{i=1}{\overset{3}{}}}h_i(𝑨)\mathrm{cos}\phi _i={\displaystyle \underset{i=1}{\overset{3}{}}}(f_i+g_i𝛀𝑨)\mathrm{cos}\phi _i.`$
It contains a shift of the Fourier modes and a rescaling of time and of the actions. The shift of the Fourier modes is constructed such that it exchanges the mode $`(1,0,0)`$ and $`(0,0,1)`$ without changing the mode $`(0,1,0)`$. More precisely, we require that the new angles $`𝝋^{}`$ satisfy :
$`\mathrm{cos}\phi _1^{}=\mathrm{cos}\phi _3,`$
$`\mathrm{cos}\phi _2^{}=\mathrm{cos}\phi _2,`$
$`\mathrm{cos}\phi _3^{}=\mathrm{cos}\phi _1.`$
This is performed by the following linear canonical transformation :
$$(𝑨,𝝋)(𝑨^{},𝝋^{})=(P𝑨,P𝝋),$$
where we recall that $`P`$ is symmetric and orthogonal. The vectors $`𝝎`$ and $`𝛀`$ are changed into $`P𝝎`$ and $`P𝛀`$. We impose that the images $`𝝎^{}`$ and $`𝛀^{}`$ of the vectors $`𝝎`$ and $`𝛀`$ by the renormalization $`_P`$ satisfy the following normalization conditions : $`𝛀^{}`$ must be of (euclidean) norm one, and the third component of $`𝝎^{}`$ must be equal to one. Since $`P`$ is orthogonal, $`P𝛀`$ is again of norm one, and thus $`𝛀^{}=P𝛀`$. The third component of $`P𝝎`$ is equal to $`\omega _1`$. We rescale the time by a factor $`\omega _1`$, i.e. we multiply the Hamiltonian by $`1/\omega _1`$. Then $`𝝎^{}`$ is given by $`𝝎^{}=𝝎/\omega _1=(\omega _3/\omega _1,\omega _2/\omega _1,1)`$.
The quadratic part of the Hamiltonian $`H_0`$ becomes $`(𝛀^{}𝑨)^2/(2\omega _1)`$. In order that this quadratic term is equal to $`(𝛀^{}𝑨)^2/2`$, we rescale the actions by a factor $`\lambda _P=1/\omega _1`$, i.e. we change the Hamiltonian $`H`$ into $`\lambda _PH(𝑨/\lambda _P,𝝋)`$.
In summary, a Hamiltonian $`H`$ given by Eq. (8) is mapped into
$$H^{}(𝑨,𝝋)=𝝎^{}𝑨+\frac{1}{2}(𝛀^{}𝑨)^2+\underset{i=1}{\overset{3}{}}(f_i^{}+g_i^{}𝛀^{}𝑨)\mathrm{cos}\phi _i,$$
where $`𝝎^{}=(\omega _3/\omega _1,\omega _2/\omega _1,1)`$, $`𝛀^{}=(\mathrm{\Omega }_3,\mathrm{\Omega }_2,\mathrm{\Omega }_1)`$, and
$`f_1^{}=f_3/\omega _1^2,`$
$`f_2^{}=f_2/\omega _1^2,`$
$`f_3^{}=f_1/\omega _1^2,`$
$`g_1^{}=g_3/\omega _1,`$
$`g_2^{}=g_2/\omega _1,`$
$`g_3^{}=g_1/\omega _1.`$
The renormalization operator $`_P`$ is equivalent to the 10-dimensional map given by the above equations (we recall that we impose a normalization condition on $`𝝎`$ and $`𝛀`$) :
$$(𝝎,𝛀,f_i,g_i;i=1,2,3)(𝝎^{},𝛀^{},f_i^{},g_i^{};i=1,2,3).$$
### B Definition of $`_L`$
The renormalization operator $`_L`$ is associated to the matrix $`L`$ and acts on the family of Hamiltonians (8). It contains an elimination of the mode $`𝝂_1`$ of the perturbation, and a rescaling procedure (shift of the resonances, rescaling of time and of the actions) such that the image of a Hamiltonian $`H`$ given by Eq. (8) is of the same general form as $`H`$ and describes the system on a smaller scale in phase space and at a longer time scale. This operator is similar to the approximate renormalization constructed for the spiral mean torus in Refs. .
We eliminate the mode $`𝝂_1`$ of the scalar functions $`f`$ and $`g`$, by a near-identity canonical transformation. We perform a Lie transformation generated by a function $`S(𝑨,𝝋)`$. The image of a Hamiltonian $`H`$ is given by :
$$H^{}=\mathrm{exp}(\widehat{S})H=H+\{S,H\}+\frac{1}{2}\{S,\{S,H\}\}+\mathrm{},$$
(9)
where $`\{,\}`$ denotes the Poisson bracket :
$$\{f,g\}=\frac{f}{𝝋}\frac{g}{𝑨}\frac{g}{𝝋}\frac{f}{𝑨},$$
and the operator $`\widehat{S}`$ acts on $`H`$ like $`\widehat{S}H=\{S,H\}`$. The generating function is chosen linear in the actions and of the form :
$$S(𝑨,𝝋)=(z+y𝛀𝑨)\mathrm{sin}\phi _1.$$
(10)
We choose $`z`$ and $`y`$ in the following way :
$`z=f_1/\omega _1,`$
$`y=(g_1f_1\mathrm{\Omega }_1/\omega _1)/\omega _1.`$
Then the generating function $`S`$ satisfies :
$$\{S,H_0\}=(Q(𝑨)h_1(𝑨))\mathrm{cos}\phi _1,$$
with $`Q(𝑨)=y\mathrm{\Omega }_1(𝛀𝑨)^2`$. The image of a Hamiltonian $`H`$ given by Eq. (8) is :
$`H^{}=`$ $`H_0+h_1\mathrm{cos}\phi _1+h_2\mathrm{cos}\phi _2+h_3\mathrm{cos}\phi _3`$
$`+\{S,H_0\}+\{S,h_1\mathrm{cos}\phi _1\}+\{S,h_2\mathrm{cos}\phi _2\}`$
$`+\{S,h_3\mathrm{cos}\phi _3\}+{\displaystyle \frac{1}{2}}\{S,\{S,H_0\}\}+O(\epsilon ^3).`$
We neglect the term $`Q(𝑨)\mathrm{cos}\phi _1`$ produced by the terms $`\{S,H_0\}+h_1\mathrm{cos}\phi _1`$, i.e. we neglect the dependence on the angles of the quadratic terms in the actions. The term $`\{S,h_1\mathrm{cos}\phi _1\}`$ is of degree one in the actions and contains the modes $`\mathrm{𝟎}`$ and $`2𝝂_1`$. We neglect the mode $`2𝝂_1`$ and we eliminate the mode $`\mathrm{𝟎}`$ of the linear term in the actions by a shift : $`𝑨^{}=𝑨+𝒂`$ where $`𝒂`$ is of order $`\epsilon ^2`$ (in the direction of $`𝛀`$). We neglect the modes with frequency vectors $`𝝂_1\pm 𝝂_2`$ produced by the term $`\{S,h_2\mathrm{cos}\phi _2\}`$. The term $`\{S,h_3\mathrm{cos}\phi _3\}`$ generates the modes $`𝝂_1\pm 𝝂_3`$. We neglect the mode $`𝝂_1+𝝂_3`$ and keep the next relevant Fourier mode $`𝝂_1𝝂_3`$ whose amplitude is denoted $`h_3^{}=f_3^{}+g_3^{}𝛀𝑨`$, where
$`f_3^{}=(zg_3\mathrm{\Omega }_1+yf_3\mathrm{\Omega }_3)/2,`$ (11)
$`g_3^{}=(\mathrm{\Omega }_1+\mathrm{\Omega }_3)yg_3/2.`$ (12)
The term $`\{S,Q\mathrm{cos}\phi _1\}/2`$ produced by $`\{S,\{S,H_0\}\}/2`$ gives a contribution to the quadratic part of the Hamiltonian. The new quadratic part is equal to $`m(𝛀𝑨)^2/2`$ where $`m`$ is given by :
$$m=1+\frac{3}{2}y^2\mathrm{\Omega }_1^2.$$
(13)
Then after the elimination of the mode $`𝝂_1`$, the new Hamiltonian is equal to :
$`H^{}=`$ $`𝝎𝑨+m(𝛀𝑨)^2/2+h_2\mathrm{cos}\phi _2+h_3\mathrm{cos}\phi _3`$ (15)
$`+h_3^{}\mathrm{cos}(\phi _1\phi _3),`$
where $`h_3^{}`$ is given by Eqs. (11) and (12).
We shift the Fourier modes according to the following linear canonical transformation :
$$(𝑨,𝝋)(𝑨^{},𝝋^{})=(L^1𝑨,\stackrel{~}{L}𝝋),$$
such that $`\mathrm{cos}\phi _2=\mathrm{cos}\phi _1^{}`$, $`\mathrm{cos}\phi _3=\mathrm{cos}\phi _2^{}`$ and $`\mathrm{cos}(\phi _1\phi _3)=\mathrm{cos}\phi _3^{}`$. The vectors $`𝛀`$ and $`𝝎`$ are changed in the following way :
$`𝛀^{}={\displaystyle \frac{\stackrel{~}{L}𝛀}{\stackrel{~}{L}𝛀}}=(\mathrm{\Omega }_2,\mathrm{\Omega }_3,\mathrm{\Omega }_1\mathrm{\Omega }_3)/(1+\mathrm{\Omega }_3(\mathrm{\Omega }_32\mathrm{\Omega }_1)),`$
$`𝝎^{}={\displaystyle \frac{\stackrel{~}{L}𝝎}{(\stackrel{~}{L}𝝎)_3}}=(\omega _2/(\omega _11),1/(\omega _11),1).`$
We rescale time in such a way that the linear term in the actions of $`H_0`$ is equal to $`𝝎^{}`$ : we multiply the Hamiltonian $`H^{}`$ by a factor $`1/(\omega _11)`$. The quadratic term of the new Hamiltonian is equal to $`m\stackrel{~}{L}𝛀^2(𝛀^{}𝑨)^2/(2\omega _12)`$. In order to map this quadratic part into $`(𝛀^{}𝑨)^2/2`$, we rescale the actions by a factor
$$\lambda _L=m\stackrel{~}{L}𝛀^2/(\omega _11),$$
by changing the Hamiltonian $`H^{}`$ into $`\lambda _LH^{}(𝑨/\lambda _L,𝝋)`$.
After this elimination and rescaling procedures, a Hamiltonian $`H`$ is mapped into
$$H^{\prime \prime }=𝝎^{}𝑨+\frac{1}{2}(𝛀^{}𝑨)^2+h_1^{\prime \prime }\mathrm{cos}\phi _1+h_2^{\prime \prime }\mathrm{cos}\phi _2+h_3^{\prime \prime }\mathrm{cos}\phi _3,$$
where $`h_i^{\prime \prime }=f_i^{\prime \prime }+g_i^{\prime \prime }𝛀^{}𝑨`$ and
$`f_1^{\prime \prime }=f_2c_f,`$
$`f_2^{\prime \prime }=f_3c_f,`$
$`f_3^{\prime \prime }=f_3^{}c_f,`$
$`g_1^{\prime \prime }=g_2c_g,`$
$`g_2^{\prime \prime }=g_3c_g,`$
$`g_3^{\prime \prime }=g_3^{}c_g,`$
with $`c_f=m\stackrel{~}{L}𝛀^2/(\omega _11)^2`$ and $`c_g=\stackrel{~}{L}𝛀/(\omega _11)`$.
The renormalization transformation is equivalent to the 10-dimensional map
$$(𝝎,𝛀,f_i,g_i;i=1,2,3)(𝝎^{},𝛀^{},f_i^{\prime \prime },g_i^{\prime \prime };i=1,2,3),$$
given by the above equations.
### C Renormalization transformation $``$
In summary, the renormalization $`=_L^2_P_L`$ acts in the following way : a canonical transformation eliminates the three main Fourier modes and produce the next three resonances which are the main Fourier modes at a smaller scale in phase space. A rescaling procedure shift these resonances into the original Fourier modes, and normalize the new Hamiltonian by rescaling the time and the actions.
The rescaling procedure acts on $`𝝎_0`$ in the following way : $`𝝎_0`$ is changed into $`𝝎^{(1)}=(1/(3\tau ^2),1/(3\tau \tau ^3),1)`$ by $`_L`$, then into $`𝝎_0^{(2)}=(3\tau ^2,\tau ^1,1)`$ by $`_P`$, then into $`𝝎_0^{(3)}=(1/(2\tau \tau ^3),1/(2\tau ^2),1)`$ by $`_L`$, and after the final step $`_L`$, into $`𝝎_0^{(4)}=𝝎_0`$. Thus the renormalization $``$ is defined for a fixed frequency vector $`𝝎_0`$. In the same way, the action of $``$ on $`𝛀`$ reduces to the map :
$$𝛀^{}=\frac{\stackrel{~}{N}𝛀}{\stackrel{~}{N}𝛀}.$$
(16)
The spectrum of $`\stackrel{~}{N}`$ is real and consists of the following eigenvalues $`\tau 1.2469`$, $`(\tau +1)^10.4450`$, $`1+\tau ^11.8019`$. Thus by iterating the map (16), the vector $`𝛀`$ converges to the eigenvector (denoted $`𝛀^{(3)}`$, with euclidean norm one) associated with the largest eigenvalue of $`\stackrel{~}{N}`$. With this reduction, the renormalization $``$ reduces to a 6-dimensional map $`(f_i,g_i;i=1,2,3)(f_i^{},g_i^{};i=1,2,3)`$.
The renormalization transformation we define can be also constructed for the following family of Hamiltonians :
$$H(𝑨,𝝋)=𝝎_0𝑨+\frac{1}{2}𝑨M𝑨+𝒈(𝝋)𝑨+f(𝝋),$$
(17)
where $`M`$ is a $`3\times 3`$ symmetric matrix with non-zero mean-value, and $`𝒈`$ is a three dimensional vector. Applying the renormalization changes the matrix $`M`$ into :
$$(M)=\frac{\stackrel{~}{N}MN}{\text{tr }(\stackrel{~}{N}MN)},$$
(18)
where $`\text{tr }(\stackrel{~}{N}MN)`$ is the trace of the matrix $`\stackrel{~}{N}MN`$. By iterating the map (18), the matrix converges to $`𝛀^{(3)}𝛀^{(3)}`$. Moreover, $`𝒈`$ is renormalized into $`\stackrel{~}{N}𝒈`$ which tends to be aligned to $`𝛀^{(3)}`$ by iteration. Then the Hamiltonians (17) tend to the degenerate Hamiltonians (1) under the iterations of renormalization. We expect that the Hamiltonians (17) belong to the same universality class as the Hamiltonians (1).
## III Renormalization flow
The renormalization transformation $``$ has the following properties : $``$ has an attractive integrable fixed point $`H_0`$ given by
$$H_0(𝑨)=𝝎_0𝑨+\frac{1}{2}(𝛀^{(3)}𝑨)^2,$$
and another non-integrable fixed point $`H_{}`$ which lies on the boundary of the domain of attraction of $`H_0`$. Outside the closure of the domain of attraction of $`H_0`$, the iterations of renormalization diverge to infinity. The renormalization dynamics has the same qualitative features as the renormalization for the golden mean torus for Hamiltonian systems with two degrees of freedom . The properties of critical tori are given by the analysis of the renormalization around the non-trivial fixed point $`H_{}`$. In particular, the existence of a non-trivial fixed point for an exact renormalization would imply self-similarity of critical invariant tori. The stable manifold of $`H_{}`$ is of codimension one; the linearized renormalization around $`H_{}`$ has only one eigenvalue of modulus larger than one : $`\delta 3.4414`$. The value of the total rescaling coefficient in the actions is $`\lambda 11.2726`$, and the rescaling coefficient of time is $`\tau +12.2469`$.
We also find a critical fixed cycle with period 7 on the critical surface. This periodic cycle is obtained from the non-trivial fixed point $`H_{}`$ by changing the sign of the Fourier coefficients $`f_i`$ and $`g_i`$. The existence of this cycle is explained by symmetry reasons as it was done for the critical cycle with period 3 for the golden mean case . In particular, it involves the same critical exponents and scaling factors as the non-trivial fixed point, and belongs to the same universality class.
## IV Conclusion
We have defined two elementary operators $`_P`$ and $`_L`$. With these operators, we have defined an approximate renormalization in order to study invariant tori with frequency vector $`𝝎_0=(\tau ^2+\tau ,\tau ,1)`$. The renormalization has a hyperbolic fixed point with codimension one stable manifold. Consequently, we expect critical invariant tori with this frequency vector to be self-similar at criticality. In order to give a firm basis of this statement, it will be interesting to build an exact renormalization transformation without the drastic approximations we used.
We notice also that the approximate renormalization we define in this paper can be generalized to an arbitrary incommensurate frequency vector $`𝝎_0`$ given by its Farey sequence, with the operators $`_P`$ and $`_L`$. For a periodic Farey sequence the hyperbolic invariant sets are expected to be fixed points, periodic orbits or strange non-chaotic attractors (such as the spiral mean torus ). For a non-periodic Farey sequence, this fixed sets can be strange chaotic attractors, but this point has not been investigated yet.
## acknowledgments
We acknowledge useful discussions with G. Benfatto, G. Gallavotti, H.R. Jauslin, H. Koch, and J. Laskar. Support from EC Contract No. ERBCHRXCT94-0460 for the project “Stability and Universality in Classical Mechanics” is acknowledged. CC thanks support from the Fondation Carnot and from the British Council – Ministère des Affaires Etrangères Alliance Program. |
warning/0001/nlin0001065.html | ar5iv | text | # Three-dimensional pattern formation, multiple homogeneous soft modes, and nonlinear dielectric electroconvection
## I Introduction
Spontaneous formation of spatially periodic structures on a homogeneous background is ubiquious in nature, fascinating to look at, and often hard to understand in detail. The periodic structures are almost never ideal. Irregularities may be generated in the transient after the pattern formation is initiated and anneal after some time. Or, in particular in dissipative systems far from equilibrium, they may be the result of an instability of the regular, spatially periodic state itself, then often leading to a state which exhibits persistent spatio-temporally chaotic dynamics. In some systems this is the case arbitrary close to the threshold of pattern formation in control parameter space. With the help of reduced descriptions like phase-diffusion and Ginzburg-Landau like amplitude equations, which are to some extent universal (i.e., independent of physical details), several phenomena associated with these deviations from the simple periodic structure can be explained .
For mostly practical reasons, experimental and theoretical research on pattern formation and dynamics has concentrated on quasi one- or two-dimensional systems, but most of the results obtained should have a direct correspondence also in genuinely three dimensional patterns. By a genuinely three dimensional pattern (below simply 3D pattern) I do here mean a spatially periodic structure for which (i) the spatial period(s) are not determined by the spatial extension of the sample (referred to as the class of patterns formed by “competing interaction” in Ref. ) and (ii) for which the spatial extension of the sample is in all directions large compared to the period(s) of the pattern. This is a stronger conception of a “3D pattern” than the one used in Refs. , which is based only on (ii).
Additional complications in 3D patterns, as compared to 1D or 2D patterns, arise from the structure and dynamics of defects (dislocations as well as disclinations), which are point like in 2D but line like in 3D. The implications for dissipative, nonpotential systems, mainly those described by the complex Ginzburg-Landau equation, have been addressed by several authors . But there is another particularity of 3D patterns, which has so far found little attention: the massive occurrence of *homogeneous soft modes*, which couple to the pattern and can drastically change its dynamics.
By homogeneous soft modes I mean marginally stable or slowly decaying homogeneous or long-wavelength perturbations of the homogeneous basic state from which the pattern arises. In the abstract sense of the word, they are hydrodynamic modes of the basic state. But for the sake of clarity the terms “hydrodynamic mode” and “hydrodynamics” shall here be reserved for slowly relaxing deviations from the *thermodynamic* equilibrium in an unbounded, homogeneous medium, and their dynamics (e.g., the velocity field in a convective flow is a hydrodynamic variable). The pattern-forming basic state is itself a non-equilibriums state. Thus, although there is some correspondence between homogeneous soft modes and hydrodynamic modes (in the narrow sense), the notions are not identical.
As it has become clear by the investigation of several 1D and 2D model systems, homogeneous soft modes are the key for understanding many of the phenomena occurring at, or close to, the onset of pattern formation. The most prominent example is the mode associated with a homogeneous perturbation of the pressure field in Rayleigh-Bénard convection which leads to a “singular mean flow”. It is worth noticing that, since it is usually possible to construct self-consistent amplitude equations which do *not* include the effect of homogeneous soft modes, their relevance is easily underestimated in the theoretical analysis.
In 1D and 2D pattern-forming systems, most hydrodynamic modes are damped by the boundaries enclosing the system. For example, momentum and heat can usually diffuse freely through the boundaries and are stabilized by large external reservoirs. Obviously, this mechanism is ineffective in systems which are extended in all three spatial dimensions. On the other hand, some coupling to a reservoir will also be required in 3D in order to sustain non-equilibrium pattern formation. This could be through electromagnetic fields, some matrix embedding the active, pattern-forming medium, or some chemical reactant provided in excess. But the couplings to the reservoirs are highly specific in these cases and stabilize only a few hydrodynamic variables. The remaining fields do then lead to homogeneous soft modes. *As a result, several homogeneous soft modes should be considered as the rule in 3D, pattern-forming systems.*
For example, when studying the 3D structures formed by chemical waves in the Belousov-Zhabotinsky (BZ) reaction, the dynamics of the plain BZ reagent does also involve convective fluid motion. Since these hydrodynamic modes are usually considered to be a nuisance, they are suppressed by embedding the reagent in a gel . But another homogeneous soft mode excited by the pattern, the temperature field (gradients of which are probably driving the convection) remains. Since temperature gradients have a strong influence on the dynamics of the pattern , a complete description of the 3D BZ reaction should explicitly involve this mode.
The work presented here is a case study of 3D pattern formation in the dielectric regime of electroconvection (EC) in nematic liquid crystals. The system was chosen because of its easy experimental accessibility. In particular, the electric nature of the instability allows to obtain patterns with several hundred periods extension in cells of a fingernail’s size, evolving on the time scale of seconds. A closely related variant, the conduction regime of EC, which always leads to quasi 2D patterns, is currently one of the best understood experimental pattern-forming systems, on the phenomenological as well as on the quantitative level (see the reviews ). These advantages compensate the inconvenience of dealing with rather complicated (electro-)hydrodynamic equations.
Rather than trying to understand the complex pattern dynamics itself, this paper is mainly devoted to the development of consistent reduced descriptions of the dynamics. Section II sketches the experimental phenomenon and the hydrodynamics of dielectric EC, emphasizing its 3D nature. Approximations used for an analytic or semi-analytic description of dielectric EC are introduced in Section III, thereby discussing the linear stability problem. In Section IV the 3D amplitude formalism for dielectric EC is derived. Close to the threshold of EC and in a liquid-crystal slab of large but finite thickness, the pattern dynamics becomes essentially 2D. The corresponding equations of motion are derived from the 3D formalism in Section V. In Section VI the stability of ideal periodic patterns is investigated and in Section VII a general scenario for the transition from the onset of dielectric EC to fully 3D pattern dynamics with increasing external stress is developed. Section VIII discusses possible experiments based on electric Nusselt number measurements and Section IX summarizes the results. Appendix A contains some analytic and numerical results for coupling coefficients, Appendix B compares two different methods for integrating multiple, homogeneous soft modes into the amplitude formalism in a general framework; one method is used in the main text.
## II Some phenomenology of electroconvection
Notice that below some points are oversimplified in order to ease intuition. For comprehensive reviews of EC see Refs. , for introductions into nemato-hydrodynamics Refs. .
### A Basic phenomena
In the typical experiment a nematic liquid crystal with negative dielectric anisotropy is sandwiched between a pair of transparent, parallel electrodes (separation $`d2050\mu \mathrm{m}`$, area $`1\mathrm{c}\mathrm{m}^2`$). By a special treatment of the electrode surfaces, the nematic director $`\stackrel{}{n}`$ (the locally averaged molecular orientation; $`|\stackrel{}{n}|=1`$) is forced to align parallel to the electrodes in some preferred direction which shall here be identified with the $`x`$-direction ($`z`$ be normal to the electrodes, $`y`$ normal to $`x`$ and $`z`$). An ac voltage $`E_0\widehat{z}d\mathrm{cos}\omega t`$ is applied at the electrodes. In the *conduction regime* at frequencies below the *cut-off frequency* $`\omega _c`$, the first instability to be observed as the voltage is increased is towards a pattern of convection rolls called Williams domains .
At higher frequencies a different kind of structure periodic along $`x`$ is found. Compared to Williams domains it has shorter wavelength and decays faster after switching of the voltage (fast turnoff mode). At least two concurring mechanisms have been proposed for this high frequency mode: the *dielectric EC* , which depends essentially on the anisotropy of the nematic (its threshold diverges at the nematic-isotropic phase transition ), and the *isotropic mechanism* where the liquid crystal’s anisotropy is not essential for the convection mechanism itself but only for selecting a preferred modulation direction. It has a finite threshold at the nematic-isotropic phase transition as its characteristic signature . The two linear modes have the same symmetry and do in principle mix, but generally the corresponding thresholds can be assumed to be sufficiently separated to consider the mechanisms isolatedly. The isotropic mode is thought to be located mainly near the electrodes, while the dielectric mode is maximal at mid plane. Unfortunately, it is not always clear which mode is actually observed. At least in some cases the dielectric mode could be identified by the good match of the threshold curve with theoretical predictions (e.g. ). The isotropic mechanisms will not be considered here.
For voltages slightly higher than the threshold of dielectric EC, the formation of the *chevron* superstructure is observed: defects (dislocations) in the pattern of convection rolls accumulate along lines oriented in $`y`$ direction, such that the topological charge of the defects alternates from line to line. Between the lines, the convection rolls are rotated and the nematic director is twisted, alternately clock- and counterclockwise . The observation of chevron patterns in the conduction regime of EC with homeotropic director alignment shows that this scenario is not restricted to a particular convection mechanism.
### B Hydrodynamic equations and material parameters
EC in both the conduction and the dielectric regime result from the interaction of electric field, space charges, mass flow, and the nematic director via the Carr-Helfrich mechanism: Spatial modulations of the director orientation are amplified by an inhomogeneous mass flow generated by electric volume forces on space charges which accumulate due to inhomogeneous electric currents in the inhomogeneous director field. Thus Maxwell’s equations <sup>*</sup><sup>*</sup>*In MKSA units. (in the quasi-static approximation We will only allow for homogeneous magnetic fields. $`\text{curl}\stackrel{}{E}=\text{curl}\stackrel{}{H}=0`$) and the balance equations for charge, momentum, mass (continuity equation), and the torque acting on $`\stackrel{}{n}`$ have to be taken into account. They contain several material parameters: the conductivities $`\sigma _{}\sigma _{}=O(10^9\mathrm{}10^5\mathrm{\Omega }^1\text{m}^1)`$ for electric currents parallel ($``$) and perpendicular ($``$) to $`\stackrel{}{n}`$ respectively (they vary on a large range depending on purity and doping, while $`\sigma _{}/\sigma _{}`$ changes only little), the dielectric constants For convenience, the factor $`ϵ_0`$ is absorbed into $`ϵ_{}`$ and $`ϵ_{}`$. $`ϵ_{},ϵ_{}=O(ϵ_0)`$ (the quantities $`\sigma _a:=\sigma _{}\sigma _{}=O(\sigma _{})`$, $`ϵ_a:=ϵ_{}ϵ_{}=O(ϵ_{})`$ measure their anisotropies), the flexoelectric constants $`e_1,e_3=O(10^{12}10^{11}\text{C}\text{m}^1)`$ ($`e_+:=e_1+e_2`$, $`e_{}:=e_1e_2`$), the diffusion constants for (ionic) charge carriers $`O(10^{11}\text{m}^2\text{s}^1)=:D_\rho `$, the mass density $`\rho _\text{m}=O(10^3\text{kg}\text{m}^3)`$, the five independent viscosities $`\alpha _1,\mathrm{},\alpha _5=O(0.1\text{N}\text{m}^2\text{s})`$ ($`\alpha _6=\alpha _2+\alpha _3+\alpha _5`$, $`\gamma _1=\alpha _3\alpha _2`$, $`\gamma _2=\alpha _3+\alpha _2`$, $`2\eta _1=\alpha _2+\alpha _4+\alpha _5`$, $`\eta _2=\gamma _2+\eta _1`$), and the curvature elasticities of the director field $`k_{22}k_{11}k_{33}=O(10^{11}\text{N})`$. I also include the “dynamic flexoelectric effect”, which was predicted on the basis of a systematic rederivation of nematohydrodynamics, but, to the authors knowledge, has not been detected, yet. It is characterized by a parameter $`\zeta ^E`$ and leads to additional dissipative contributions in the charge, momentum, and torque balance equations.
### C Dimensional analysis
With the exception of the charge relaxation time $`\tau _0:=ϵ_{}/\sigma _{}=O(10^610^1\text{s})`$ the nematohydrodynamic equations (without external fields), being derived as a limit of large time and length scales (though typically valid down to molecular scales), do not set any time or length scale by themselves. Instead, one finds basically three types of diffusivities: for charge ($`D_\rho `$), director orientation \[e.g. $`D_{d,\text{stat}}=k_{33}/\gamma _1=O(10^{10}\text{m}^2\text{s}^1)`$ for static, $`D_{d,\text{dyn}}=(k_{33}\eta _1)/(\gamma _1\eta _1\alpha _2^2)=O(10^9\text{m}^2\text{s}^1)`$ for dynamic deformations; notice that (static) flexoelectric effects do not introduce a new diffusive scale since $`e_{1/3}^2/ϵ_0k_{33}`$\], and momentum \[e.g. $`D_p=(\gamma _1\eta _1\alpha _2^2)/(\gamma _1\rho _m)=O(10^5\text{m}^2\text{s}^1)`$ along $`\stackrel{}{n}`$\]. The overdamped limit $`\rho _m0`$, $`D_p\mathrm{}`$ is generally a good approximation. Charge diffusion is not essential for the Carr-Helfrich mechanism and is typically screened out. Then orientational diffusion sets the only diffusive scale.
With an externally generated electric ac field $`E_0\widehat{z}\mathrm{cos}\omega t`$ two additional time scales are introduced: The period $`2\pi \omega ^1`$ and the “director time” $`\tau _d=\gamma _1/(ϵ_{}E_0^2)`$ \[the more intuitive choice $`\tau _d:=\gamma _1/(ϵ_aE_0^2)`$ would suggest that $`ϵ_a=0`$ is singular for EC, which is, for the convective modes themselves, not the case\].
In the conduction regime the charge densities oscillate with the frequency of the applied field, while the director orientation is mostly constant. This imposes a condition
$`\tau _d\omega ^1\tau _0`$ (1)
on the three time scales. The finite sample thickness $`d`$ determines the wavelength $`\lambda `$ of the convection pattern and leads through a condition $`d\lambda (D_{d,\text{stat}}\tau _d)^{1/2}`$ to a voltage threshold $`V_c^2=d^2E_c^2k_{33}/ϵ_{}`$ for the onset of convection (see Ref. for a good analytic formula). A lower limit for the sample thickness is given via relation (1) by $`d^2D_{d,\text{stat}}\tau _0`$. At frequencies higher then the cutoff frequency $`\omega _c\tau _0^1`$ the conduction mechanism is also disabled or at least supersede by the dielectric mode.
In the dielectric regime director and fluid flow oscillate with $`\omega `$, which leads to a condition
$$\tau _d\omega ^1.$$
(2)
The charge distribution is at high enough frequencies ($`\omega ^1\tau _0`$) mostly constant in time. This is not actually necessary for the dielectric mechanism to be effective , but typically dielectric EC is supersede by the conductive mode at lower $`\omega `$. The threshold for the onset of EC is now given by a condition $`\tau _d\omega ^1`$ (or $`E_c^2ϵ_{}/\gamma _1\omega `$), i.e., the lowest $`E_0`$ compatible with relation (2). The wavelength $`\lambda `$ of the critical mode can under some conditions be $`d`$ , but for typical materials used it is $`\lambda (D_{d,\text{dyn}}/\omega )^{1/2}`$, at least as long as this length is smaller than $`d`$ and larger than the Debye screening length $`(D_\rho \tau _0)^{1/2}`$, i.e., $`\omega \tau _0<D_{d,\text{dyn}}/D_\rho =O(10^2)`$, where charge diffusion becomes important.
Thus, the length scales given by the spatial period of the pattern $`\lambda `$ and the sample thickness $`d`$ are usually independent and easily separated ($`\lambda d`$) in the dielectric regime, either by increasing $`d`$ or by simultaneously increasing $`\omega `$ and the conductivities, while leaving the secondary control parameter $`\omega \tau _0`$ constant. Since strong doping may affect the nematic material parameters and the nematodynamics at high frequencies is not fully understood, the program carried out below is best seen as the *limit of thick cells*. The theory should accurately describe typical experiments in cells with $`d10\lambda `$. To observe fully three dimensional patterns, thicker cells might be required (see also Section VII).
## III Approximation methods and linear theory
### A 2D vs. 3D amplitude formalism
Below we will develop the amplitude formalism for the pattern dynamics in the dielectric regime, i.e., obtain the laws of motion of amplitude and phase of the spatial modulations as described by the complex pattern amplitudes $`A^{}(x,y,t)`$ or $`A(x,y,z,t)`$, respectively.
The basic state in the experimental cell is anisotropic and inversion symmetric and the primary bifurcation is supercritical (forward) towards a steady-state pattern with a single critical wave vector. Hence, the most elementary description of the pattern dynamics is give by the time-dependent Ginzburg-Landau equation in 2D,
$$\tau _tA^{}=\left(ϵ^{}+\xi _x^2_x^2+\xi _y^2_y^2g^{}|A^{}|^2\right)A^{},$$
(3)
important physical properties of which are reviewed in Refs. . The real, positive coefficients $`\tau `$, $`\xi _x`$, $`\xi _y`$, and $`g^{}`$ have magnitudes corresponding to natural scales of the system (e.g. the pattern wavelength for $`\xi _x`$, $`\xi _y`$) and can be calculated from the underlying hydrodynamic equations. The small, dimensionless parameter $`ϵ^{}`$ measures the distance from the threshold of pattern formation in the control-parameter space of the underlying system.
It will be shown in Sec. V that, as a direct consequence of the separation of length scales in the dielectric regime, the range of validity of Eq. (3) is highly restricted. Already for values of $`ϵ^{}`$ of the order $`\lambda ^4/d^4`$ corrections to Eq. (3) must be taken into account. For $`ϵ^{}\lambda ^2/d^2`$ the 2D description breaks down completely. But the convective dynamics can then still be described in terms of the 3D modulations of the complex pattern amplitude $`A(x,y,z,t)`$ (defined by Eq. (4) below), which is coupled nonlinearly to several homogeneous soft modes. It is therefore natural to derive first the 3D amplitude dynamics which can then be reduced further to a 2D description in a subsequent step.
Julien, Knobloch, and Tobias were the first to implement the idea of deriving a reduced description for the $`z`$-dependence of the amplitude of patterns with $`\lambda /d1`$ as an intermediate step in the theory, and also the first to observe that this method significantly eases the restriction of the control parameter to values close to threshold. Their calculation do, however, not involve in-plane modulations of the pattern and the resulting excitation of homogeneous soft modes. Several results concerning the 3D description of dielectric EC and its reduction to 2D are derived in an unpublished work by Lindner , which is quoted here whenever necessary.
### B Linear stability in 3D
The starting point for setting up the 3D amplitude equations is to calculate the linear threshold $`E_c`$, critical wave number $`q_c`$ and critical eigenvector (i.e., the $`2\pi /\omega `$-periodic time dependence of the hydrodynamic fields at threshold) ignoring any spatial variations along $`z`$. Several linear stability calculations of this type have been carried out .
In experiments, the critical wave vector is always found to be parallel to the orientation of the nematic director in the basic state (the $`x`$ direction). The linear problem is thus effectively one-dimensional, with trivial, sinusoidal variations along the remaining $`x`$ direction, and is much easier to solve than the 2D problem including variations and boundary conditions along $`z`$. This technical advantage of the 3D approach, which is of course not restricted to EC, remains effective also in the subsequent calculations of the coupling coefficients in the 3D amplitude equation.
For analytic as well as numerical calculations it is convenient to use a truncated Fourier expansion of the time dependences of the hydrodynamic fields, assuming them to be $`2\pi /\omega `$ periodic. In the simplest cases truncated at lowest order (i.e., including constant and $`\mathrm{sin}`$/$`\mathrm{cos}\omega t`$ contributions) or including the $`\mathrm{sin}`$/$`\mathrm{cos}2\omega t`$ modulation of the induced electric potential in order to better model the interplay between electric charges and fields. With these truncations, the stability problem for sinusoidal excitation can be solved explicitly (see Appendix A). Notice, however, that this “lowest order” or, respectively, “second lowest order Fourier approximation” involves some arbitrariness in the choice of variables and does not correspond to any physical limit. Numerical convergence ($`5\%`$ accuracy) requires inclusion of at least the third harmonic. The actual time dependence of, e.g., the director field, depends on $`\omega `$ and is non-trivial even as $`\omega \mathrm{}`$ .
Some results presented here rely on the first or second lowest order Fourier expansion of the induced electric potential $`\mathrm{\Phi }`$, the nematic director expressed by $`n_z`$ and $`\phi `$ such that $`\stackrel{}{n}=n_z\widehat{z}+(1n_z^2)^{1/2}\widehat{c}`$, $`\widehat{c}:=(\mathrm{cos}\phi ,\mathrm{sin}\phi ,0)`$, the velocity field $`\stackrel{}{v}`$, and the pressure $`P`$. The other hydrodynamic fields are treated implicitly. For the representation of the dielectric mode itself, $`\stackrel{}{v}`$ is expressed in the divergence-free form $`\stackrel{}{v}=(_y,_x,0)g+(_x_z,_y_z,_x^2_y^2)f`$ and the pressure is eliminated.
As a natural consequence of $`\tau _d\omega ^1`$, the relative phases of electric, director and velocity fields in the linear eigenvector are shifted by angles $`O(1)`$. Remarkably, the phase shift between the lowest Fourier mode of director oscillations and external field is $`\pi /4`$, for the first and second-lowest-order Fourier approximation (see Appendix A) exactly and only slightly perturbed ($`<1\%`$ for $`\omega \tau _0>2`$) when higher Fourier modes are included. No simple physical explanation for this result should be expected, since it holds only at the critical (most unstable) wavenumber. Experimental observations seem to agree with a value $`\pi /4`$ for the phase shift even better than the comparison with Galerkin calculations including $`z`$-dependence , which had been carried out as a test of the dielectric model of EC.
We define the pattern amplitude $`A`$ such that the amplitude $`n_{z,c}`$ of the $`\mathrm{cos}\omega t`$ Fourier component of the director tilt oscillations (which is in phase with the applied voltage) has a spatial dependence
$`n_{z,c}=A(x,y,z,t)\mathrm{exp}(iq_cx)+c.c.`$ (4)
Assuming as usual $`ϵ`$, $`_x`$, $`_y`$, $`_z`$ to be small and discarding contributions beyond the lowest nontrivial order, the linear part of the 3D amplitude equation assumes the form
$$\tau _tA=\left(ϵ+\xi _x^2_x^2+\xi _y^2_y^2+\xi _z^2_z^2\right)A.$$
(5)
As conventional, $`ϵ:=(E_0^2E_c^2)/E_c^2`$. The coherence lengths $`\xi _x`$, $`\xi _y`$, $`\xi _z`$ turn out to be $`\lambda `$ and the relaxation time $`\tau `$ is of the order of the charge relaxation time $`\tau _0`$ (s. Appendix A).
The horizontal boundary conditions for $`A`$ are simply
$$A=0\text{at}z=\pm d/2.$$
(6)
Contributions from derivatives of $`A`$ and nonlinear contributions to the boundary conditions are of higher order and can be discarded. In particular, as is well known, the distinction between free and no-slip boundary conditions for the velocity field plays no role at this point. With the realistic no-slip boundary conditions for the velocities, the relative magnitude of the hydrodynamic fields in the 2D linear eigenvector (including $`x`$ and $`z`$ variations) locally deviates from the 1D eigenvector (only $`x`$ variation) only in a boundary layer of thickness $`\lambda `$, an example of which is shown in Fig. 1. This boundary layer might provide a problem for numerical approaches directly using 2D eigenvectors in thick cells with no-slip boundaries, notably when Galerkin approximations are used.
### C Linear stability in a cell of finite thickness
Assuming as usual the lateral ($`x`$, $`y`$) extensions of the cell to be large compared to its thickness $`d`$, the trivial solution $`A0`$ of Eq. (5) with boundary conditions (6) becomes unstable at $`ϵ=ϵ_d:=(\pi \xi _z/d)^2`$ \[i.e., $`ϵ^{}=ϵϵ_d`$ in Eq. (3)\] with a critical mode $`A\mathrm{cos}(\pi z/d)`$. The small threshold shift $`E_c(1+ϵ_d/2)E_c`$ due to the $`z`$-variation is rather uninteresting by itself (there is also a shift $`\lambda ^2/d^2`$ in $`q_c`$ by a discarded contribution to Eq. (5) of the form $`i_x_z^2A`$), but the effect provides, for example, a simple interpretation of the small gap ($`\mathrm{\Delta }E_0=3ϵ_dE_c/2`$) between the critical mode and the lowest $`z`$-antisymmetric mode (i.e. $`A\mathrm{sin}(2\pi z/d)`$) reported in .
### D Flexoelectric effects
A short remark about flexoelectric effects, which are generally difficult to isolate experimentally, is in place at this point. The high symmetry of the linear problem in 3D does not allow flexoelectric effects: $`E_c`$, $`q_c`$, $`\tau `$, and $`\xi _x`$ are independent of the flexoelectric coefficients. Most of the remaining linear and nonlinear coefficients contain flexoelectric contributions. With our choice of variables these contributions are, except for some dynamic-flexoelectric terms, all indirect: flexoelectric effects excite additional, “slaved” contributions in the subspace orthogonal to the critical eigenvector, which, again by flexoelectric effects, feed back into the dynamics of the amplitude of the eigenvector. In contrast, the contributions not depending on flexoelectric coefficients are, except for $`\xi _x`$, all direct: no excitation of slaved degrees of freedom is involved. Therefore *flexoelectric effects are separated in a natural way* from the standard dynamics. This might provide methods for measuring the flexoelectric coefficients in a way not sensitive to parasitic boundary effects. Previous calculations involving flexoelectric effects where restricted to the conventional, “static” flexoelectric contributions and concentrated on the determination of critical mode and voltage.
## IV The 3D amplitude equations
### A Method
Before discussing the homogeneous soft modes relevant for dielectric EC, some comments on methodology are required.
Multiple homogeneous soft modes excited by a patterning mode have, to my knowledge, first been introduced by Plaut and Pesch . But their description requires the soft-mode amplitudes to be constant along all but one spatial direction. This limitation seems to be partly due to the procedure by which the equations were derived. There are two popular philosophies for this procedure (see below and Appendix B), which shall here be labeled as “order parameter” method and “center manifold” method. The two methods usually (for at most one homogeneous soft mode) differ only in the way in which the problem is formulated and solved, but lead to the same results. The association of existing general prescriptions for deriving amplitude equations with the former (e.g. ) and the latter philosophy (e.g. , the Chapman-Enskog approach for the derivation of hydrodynamics from statistical mechanics is also of this type, see Refs. ) is therefore not always conclusive (see also ). When using the “order parameter” method to obtain reduced equations, each spatial Fourier mode of the physical state is projected onto the (adjoint) slowly decaying linear eigenmodes of the basic state *with the corresponding wave vector*. Using the “center manifold” method, the projection is always onto the (adjoint) eigenvectors for (typically neutrally stable) *homogeneous perturbations*. When there are multiple slow modes at a single wave vector (usually $`\stackrel{}{q}=0`$), the resulting reduced equations differ, as is shown in Appendix B. Plaut and Pesch seem to be using the “order parameter” method, which leads to problems in more than one spatial dimension. Here, the “center manifold” method is used to derive the nonlinear extensions of the amplitude equation (5) including homogeneous soft modes.
### B Derivation of the soft mode equations
Some particularities of the problem under consideration have to be taken into account: The quasi-static approximation of electrodynamics $`curl\stackrel{}{E}=0=curl\stackrel{}{B}`$ and the approximation of an incompressible fluid $`\stackrel{}{v}=0`$ both lead to additional homogeneous “soft modes”: Since no time derivatives of electric potential (electric field) or pressure occur in the basic equations as they are used here, all their temporal Fourier modes are in the kernel of the linear operator $`L(0,0)`$ (s. Appendix B). When the viscid limit $`\rho _m0`$ becomes effective, the same applies for the oscillating part of the velocity field. This is the case when spatial variations occur on scales smaller than $`(\omega D_p)^{1/2}=(D_p/D_{\text{dyn}})^{1/2}\lambda =O(10^{5/2})\lambda `$. Formally we shall assume length scales to be larger than this. But, since, at least in this case, velocity and pressure oscillations do not feed back into the remaining dynamics at lowest order in the derivatives, i.e., there are no contributions $`O(\rho _m^1)`$, the description should be good also on smaller scales.
To simplify the problem further, the equations for the soft modes are here calculated only in the lowest order Fourier approximation. The “slaved” modes being eliminated are then the slowly varying average director tilt $`n_z`$, which is stabilized by the applied electric field through the dielectric anisotropy $`ϵ_a`$ (assumed to be negative hereafter), and oscillations of the director, which are viscously damped. Finally, in anticipation of corresponding boundary conditions, only small deviations from the basic state $`\stackrel{}{v},\phi ,n_z,\mathrm{\Phi }=0`$ shall be considered for now.
In order to obtain a consistent truncation of the soft-mode equations, recall that the only non-diffusive scale in the hydrodynamic equations is the charge relaxation time $`\tau _0`$. Assume $`\omega \tau _0`$ to be fixed. This also determines $`E_c\omega ^{1/2}`$ for given $`\tau _0`$ and we will assume $`E_0E_c`$. The elimination the fast modes ($`n_z`$ and oscillations of $`\stackrel{}{n}`$) becomes more efficient when $`\omega `$ and $`E_0`$ increase or, respectively, $`\tau _0`$ decreases. Then, in the limit of small $`\tau _0`$ (or large $`\sigma _{}`$, with fixed $`\sigma _{}/\sigma _{}`$), the time scale $`\tau _0`$ drops out of the equations. A purely diffusive scaling for the derivatives ($`_x^2_y^2_z^2_t`$) is retained, without making any a priori assumptions about the actual scaling laws of typical lengths and times, which may be different. This approximation breaks down when length scales become shorter than $`(D_{d,\mathrm{stat}}\tau _0)^{1/2}`$ or $`(D_{d,\mathrm{stat}}/\omega )^{1/2}`$.
It turns out that with these approximations the only *relevant* modes in the fast subspace (i.e. $`R`$ in Appendix B 3) are $`n_z`$ and its temporal modulations as given by
$$\begin{array}{cc}\hfill n_z& (x,y,z,t)=\hfill \\ & 2_z\frac{\left(k_{22}k_{11}\right)_yn_y+\alpha _3v_x}{E_0^2ϵ_a}\hfill \\ \hfill +& 2_x\frac{ϵ_aE_0\mathrm{\Phi }_r\left(e_+2\gamma _1\zeta ^E\right)_z\mathrm{\Phi }_0+\alpha _2v_z}{E_0^2ϵ_a}\hfill \\ \hfill +& \frac{2b\left(ϵ_a_x\mathrm{\Phi }_0e_{}_yn_y\right)}{E_0^3ϵ_a^2}\times \hfill \\ & \left(E_0^2ϵ_a\mathrm{cos}(\omega t)4\gamma _1\omega \mathrm{sin}(\omega t)\right).\hfill \end{array}$$
(7)
The electric potential has been decomposed as $`\mathrm{\Phi }=\mathrm{\Phi }_0+2\mathrm{\Phi }_r\mathrm{cos}\omega t2\mathrm{\Phi }_i\mathrm{sin}\omega t`$. The parameter $`b:=E_0^4ϵ_a^2/(3E_0^4ϵ_a^2+16\gamma _1^2\omega ^2)`$ measures the strength of the excitation of the oscillatory part of $`n_z`$. It is numerically small \[in the standard material MBBA ($`p`$-methoxybenzilidene-$`p^{}`$-$`n`$-butylaniline) $`b0.01`$\]. The resulting description for the soft mode dynamics is given by the Eq. (39) on page 39 (for the terms containing $`A`$, see Section IV F below).
### C Comments on the soft mode equations
Most of the terms in Eq. (39) reproduce linearized nematohydrodynamics. Equations (39a-39c) derive from the charge balance equation \[Eqs. (39b,39c) have been multiplied with $`E_0`$\], Eq. (39d) from the angular momentum balance on $`\stackrel{}{n}`$, Eqs. (39e-39g) from the Navier-Stokes equation and Eq. (39h) is the unchanged continuity equation. In Eq. (39d) a term $`\chi _aH_y^2\phi `$ has been included, which describes the action of a magnetic field in $`y`$ direction. It will be used in Section V. The contributions resulting from the elimination of fast modes are underlined.
Remarkably, the equations are mostly independent of the strength of the external field (a factor $`E_0`$ can be absorbed into the definitions of $`\mathrm{\Phi }_r`$ and $`\mathrm{\Phi }_i`$), although it is the cause for the excitations of the slaved modes. As an example, consider the mechanism for the reduction of viscosity by the term $`\alpha _2_x^2v_z`$ in Eq. (39g) (recall $`\alpha _2<0`$): Shear forces $`\alpha _2_xv_z`$ excite $`n_z`$. This leads to polarization charges $`_xE_0ϵ_an_z`$. The electric field $`E_0`$ acting on these charges generates bulk forces on the fluid. On the other hand, the excitation of $`n_z`$ is damped by electric forces $`ϵ_aE_0^2`$ and the factor $`E_0^2`$ cancels out.
The soft mode equations reflect the non-equilibrium character of the basic state. For example, if Onsager’s relations would hold, the coefficients of $`_x_zv_x`$ in Eq. (39g) and of $`_x_zv_z`$ in Eq. (39e) would be the same. At low $`\omega `$ the basic state is even unstable. The mechanism corresponds to EC in the conduction regime. Ignoring flexoelectric effects, the cutoff frequency $`\omega _c`$ above which the basic state stabilizes is given by
$$\omega _c^2=\frac{\sigma _{}\left(\alpha _2ϵ_{}\sigma _{}\alpha _2ϵ_{}\sigma _{}ϵ_a\eta _1\sigma _{}\right)}{ϵ_aϵ_{}ϵ_{}\eta _1},$$
(8)
which reproduces the result of direct stability calculations using (effectively) the same Fourier truncation . Only the threshold field for the Williams domains is too small to be resolved by Eq. (39).
### D Nonlinear extensions
Of course, constant values can always be added to any of the soft modes by a Galilei transformation, a rotation, or a gauge transformation. The problem of adding nonlinear contributions (e.g. advection terms) to Eqs. (39) such that they become formally invariant under these transformation is easily solved. The solution is not unique, but it can be seen by inspection that the precise form of the nonlinearities does not matter under the following conditions:
1. Dynamics is such that, in fact, diffusive scaling holds. In particular this implies that, if $`L\lambda `$ is the typical length scale (i.e. $`_x,_y,_zL^1`$), the variations of $`\stackrel{}{v}`$, $`\mathrm{\Phi }_r`$, $`\mathrm{\Phi }_i`$, $`\phi /L`$, and $`\mathrm{\Phi }_0/L`$ over $`L`$ scale in the same manner as $`L\mathrm{}`$.
2. Variations in $`\phi `$ over $`L`$ are much smaller than one. It is not necessary to specify the scaling relation of $`\phi `$ and $`L`$. When $`L`$ is determined through the bulk dynamics, $`\phi `$ may actually vary by $`O(1)`$ over the sample.
It should be noticed that, although the second condition is satisfied for many problems of pattern dynamics, it is too strong for disclinations (line defects) in the director field at any distance $`R`$ from the core of the disclination: On the typical length scale $`L=R`$ variations of $`\phi `$ are $`O(1)`$. Then, for example, nontrivial nonlinear contributions from the velocity field $`\stackrel{}{v}L^1`$, like a term of the form $`(_xv_z)(_y\phi )`$ in Eq. (39g), might have to be included.
For the part describing the curvature elasticity in Eq. (39d), the fully nonlinear corrections in $`\phi `$ have been calculated by extending the “center manifold” method to nonlinear contributions and requiring rotation invariance. The result has the same form as close to equilibrium:
$$\begin{array}{c}\gamma _1_t\phi =k_{33}\widehat{c}_{}(_x^2+_y^2)\widehat{c}\hfill \\ \hfill +(k_{11}^{}k_{33})(\widehat{c}_{})(\widehat{c})+k_{22}\widehat{c}_{}_z^2\widehat{c}+\mathrm{}\end{array}$$
(9)
\[$`\widehat{c}_{}:=(\mathrm{sin}\phi ,\mathrm{cos}\phi ,0)`$\], however, with $`k_{11}^{}:=k_{11}+2be_{}^2/ϵ_a`$. Although an additional term proportional to $`(\widehat{c}_{})(\widehat{c})`$ would be thinkable, it does not occur in the present approximation. For a rotationally invariant description of the dynamics of the convection pattern, it is useful to go over to a representation in terms of $`\stackrel{~}{A}(\stackrel{}{r})=A(\stackrel{}{r})\mathrm{exp}(iq_cx)`$ as in Ref. .
### E Stability of the twist mode with respect to $`x`$ modulations
Another point to notice is that below the threshold of dielectric EC (i.e. with $`A=0`$) the system can, in our approximation, never destabilize in such a way that $`\phi `$ is excited but $`_y\phi =0`$, even when allowing for flexoelectric effects and arbitrary $`z`$ dependencies. Without $`y`$ modulations, the in-plane director couples only to $`v_x`$, and this interaction is not affected by the elimination of the fast modes and hence relaxational. This is remarkable because such modulation instabilities of $`\phi `$ below the EC threshold are apparently observed in experiments (the “inertial mode” which was proposed as an explanation has the wrong symmetry). A far-fetched but possible explanation would be that the modulations in $`\phi `$ are generated through a mechanism which is similar to the one that generates the chevron superstructure in the dielectric regime, however, invoked by the convection rolls of “isotropic” EC (see Section II A). The isotropic mechanism is expected to become active in the respective experimental situations but does itself not involve excitations of $`\phi `$. The convection rolls might themselves be smaller than the optical resolution and therefore remain unobserved. Then the onset of $`\phi `$ modulation would *appear* to be the threshold of a primary instability of the homogeneous basic state, although the actual primary (“isotropic”) threshold is at slightly lower voltages.
### F Interaction with the pattern
The equation of motion for the pattern amplitude itself is given by
$$\begin{array}{c}\tau (_t+\stackrel{}{v}+iv_xq_c)A=[\kappa _x_xv_x+\kappa _z_zv_z\hfill \\ \hfill +\xi _x^2_x^2+\xi _y^2(_y^22iq_c\phi _yq_c^2\phi ^2)+\xi _z^2_z^2\\ \hfill +i\alpha _x_x\mathrm{\Phi }_0+i\beta _y_y\phi \frac{4_z\mathrm{\Phi }_r}{E_c}+\epsilon g|A|^2]A.\end{array}$$
(10)
All coefficients are real. The differential operators $`_x`$, $`_y`$, $`_z`$ are used to indicate that only the immediately following expression should be differentiated, not $`A`$. With the exceptions of $`v_x_xA`$ and $`_xv_xA`$, which are included for symmetry, only terms up to lowest nontrivial order have been included, assuming time, all lengths and all fields to scale independently. The l.h.s. is given by Galilean invariance, the form of the expression following $`\xi _y^2`$ from rotation invariance . The term involving $`_z\mathrm{\Phi }_r`$ represents simply an additional contribution to the electric driving field. Symmetry would also allow a term $`\kappa _y_yv_yA`$, but the hydrodynamic equations do not generate it.
When $`\beta _y=\xi _y^2q_c`$ it is possible to absorb the corresponding term into the parenthesis after $`\xi _y^2`$ and to rewrite the complete expression in the “potential” form $`\xi _y^2(_yiq_c\phi )(_yiq_c\phi )A`$. In this case the phase of the pattern does not drift in a weakly deformed $`\phi `$ field. In fact, $`\beta _y`$ is close to this value (see , Appendix A), which is largely due to the conservation of charge and momentum and the flux-divergence form of the resulting expressions.
The Landau coefficient $`g`$ is of order unity in our normalization . As in the conduction regime , it is dominated by “geometric” effects, i.e., inhibition of the Carr-Helfrich mechanism for large $`n_z`$: besides flexoelectric effects, about 99% of $`g`$ come from contributions quadratic or cubic in $`n_z`$. This indicates that a breakdown of the weakly nonlinear expansion should be expected for $`An_z=O(1)`$.
The nonlinear excitation of the soft modes by the convection pattern is less intuitive than their feedback onto the rolls discussed above. The principle of truncation for the contributions of $`A`$ in Eq. (39) is again to keep only terms of lowest nontrivial order, however, allowing for phase gradients of $`A`$ without gradients of the modulus. As the relative scaling of $`\phi `$ vs. length scales is left undetermined at this stage (see Section IV D), the nonlinear term $`A^{}iq_c\phi A`$, which is given by rotation symmetry, must always come along with $`A^{}_yA`$.
Since Eq. (39a) must have flux-divergence form, the largest contributions from $`A`$ are $`O(^2A^2)`$, which is too small to be relevant. The coefficients $`I_r`$ and $`I_i`$ in Eqs. (39b,39c) have the dimensions of a current density and measure the strength of an alternating current $`\stackrel{}{j}_{|A|^2}=2|A|^2(I_r\mathrm{cos}\omega tI_i\mathrm{sin}\omega t)`$ generated by the convection pattern. Important contributions to the coefficients $`I_r`$,…,$`I_{iz}`$ come from charge advection.
The strongest contributions to $`\mathrm{\Gamma }`$ in Eq. (39d) are simple potential effects. In MBBA, the most important one describes a relaxation of the bend of the director modulations by twist and is given by $`4(k_{22}k_{33})`$ (see and Appendix A), which is generally negative. The second most important contribution comes from the dielectric torques from applied and induced field on the director. In MBBA this contribution to $`\mathrm{\Gamma }`$ is positive and in materials with large negative $`ϵ_a`$ (e.g. $`ϵ_{}`$) it could compensate the elastic one and reverse the sign of $`\mathrm{\Gamma }`$. Notice also that $`\mathrm{\Gamma }`$ is paricularly sensitive to flexoelectric effects (s. Table I). A negative value of $`\mathrm{\Gamma }`$ is required for the occurrence of abnormal rolls (see Sec. VI A) and chevron patterns.
The terms associated with $`S_x`$ and $`S_z`$ in Eqs. (39e,39g) represent internal stresses of the convection pattern. As for the coupling of $`A`$ to $`\stackrel{}{v}`$ in Eq. (10), a corresponding term for the $`y`$ direction or a term of the form $`S_{yx}_y\mathrm{Im}\{A^{}_xA\}`$ do not enter Eq. (39f), although they are allowed by symmetry. $`S_{xx}`$, $`S_{yy}`$, …, $`S_{zx}`$ can be interpreted as surface tensions of the planes of equal phase of the convection pattern.
The high number of soft modes and the rich, non-potential coupling almost certainly lead to spatio-temporally chaotic states already at onset, provided the spatial extensions of the sample are large enough.
## V Effectively 2D pattern dynamics near threshold
Currently more interesting than the fully 3D chaotic state are, from the experimental point of view, the quasi 2D pattern dynamics in a restricted geometry near threshold \[$`ϵ^{}(\lambda /d)^2`$\]. The dynamics is here derived for a slight generalization of the usual setup: A magnetic field $`H_y`$ might be applied along the $`y`$ direction (i.e. in-plane, normal to the rubbing direction). With $`H_y`$ slightly below the twist Freédericksz field $`\chi _aH_F^2=k_{22}(\pi /d)^2`$, the amplitude of the twist mode, which is known to be important for the pattern dynamics in any case, becomes a slow variable and must be included explicitly in the 2D formalism. The conventional setup is described by $`H_y=0`$.
### A Derivation of the 2D description
The appropriate boundary conditions for the soft modes at the enclosing electrodes are
$`\mathrm{\Phi }_0,\mathrm{\Phi }_r,\mathrm{\Phi }_i,\phi ,\stackrel{}{v}=0\text{at}z=\pm d/2.`$ (11)
Thus, all components of the electric potential and the velocity field are damped by the boundaries.
Again the “center manifold” method is used (s. Appendix B), now to reduce the 3D equations to 2D. The dynamically active part $`S`$ of the state vector $`U`$ is now given by the sum of
$`A(x,y,z)`$ $`=A^{}(x,y)\mathrm{cos}(\pi z/d),`$ (12)
$`\phi (x,y,z)`$ $`=\phi ^{}(x,y)\mathrm{cos}(\pi z/d),`$ (13)
$`P(x,y,z)`$ $`=P^{}(x,y),`$ (14)
with the amplitudes of the active modes $`A^{}`$, $`\phi ^{}`$, and $`P^{}`$.
To be specific, associate $`A`$ with the complex conjugate of Eq. (10) and $`\mathrm{\Phi }_0`$, $`\mathrm{\Phi }_r`$, $`\mathrm{\Phi }_i`$, $`\phi `$, $`v_x`$, $`v_y`$, $`v_z`$, and $`P`$ with successive equations in the system (39), and define the scalar product $`|`$ as usual as the equally weighted sum over $`z`$ integrals over products of the two components. The projector onto the slow dynamics is constructed from the bi-orthonormalized linear functionals $`(2/d)_{d/2}^{d/2}\mathrm{cos}(\pi z/d)A𝑑z`$, $`(2/d)_{d/2}^{d/2}\mathrm{cos}(\pi z/d)\phi 𝑑z`$, and $`(1/d)_{d/2}^{d/2}P𝑑z`$.
We proceed with a calculation of the excitations in the fast subspace $`R`$ by a term-by-term solution of Eq. (B12). The truncation is chosen such that the distinguished limit
$$P^{},H_F^2H_y^2,_t\epsilon ^{},$$
$$\begin{array}{c}\hfill \phi ^{},A^{},_x,_y\epsilon _{}^{}{}_{}{}^{1/2},\end{array}$$
$$\epsilon ^{}=O(\lambda /d)^4,\text{and}d\mathrm{}$$
is correctly described.
At linear order in the amplitudes and to linear order in $`_x,_y`$, the slaved part of the state vector contains only the contributions
$`v_x^{(1)}`$ $`={\displaystyle \frac{1}{\eta _2}}\left({\displaystyle \frac{d^2}{8}}{\displaystyle \frac{z^2}{2}}\right)_xP^{},`$ (15)
$`v_y^{(1)}`$ $`={\displaystyle \frac{2}{\alpha _4}}\left({\displaystyle \frac{d^2}{8}}{\displaystyle \frac{z^2}{2}}\right)_yP^{}.`$ (16)
At order $`O(|A^{}|^2)`$ and without any $`x`$ or $`y`$ modulations there are excitations of the electric field
$`\mathrm{\Phi }_r^{(2)}+i\mathrm{\Phi }_i^{(2)}=`$ $`{\displaystyle \frac{I_r+iI_i}{\sigma _{}+i\omega ϵ_{}}}{\displaystyle _{d/2}^z}|A(z^{})|^2|A|^2_zdz^{},`$ (17)
and a contribution to the pressure field, orthogonal to the active pressure mode,
$`P^{(2)}=`$ $`\left(S_z2S_E\right)\left(|A|^2|A|^2_z\right),`$ (18)
where $`|A|^2_z=|A^{}|^2/2`$ is the average of $`|A|^2`$ over $`z`$ and
$`S_E:={\displaystyle \frac{E_0ϵ_{}(\omega ϵ_{}I_i+\sigma _{}I_r)}{\omega ^2ϵ_{}^2+\sigma _{}^2}}.`$ (19)
By gradients of $`\mathrm{\Phi }_r^{(2)}`$, $`\mathrm{\Phi }_i^{(2)}`$, $`P^{(2)}`$, and by direct contributions at order $`O(_x|A^{}|^2,_y|A^{}|^2)`$ the mean flow
$`v_x^{(2)}=`$ $`[S_x({\displaystyle \frac{d^2}{16}}{\displaystyle \frac{z^2}{4}})+`$
$`(S_E+S_xS_z){\displaystyle \frac{d^2}{4\pi ^2}}\mathrm{cos}^2\left({\displaystyle \frac{\pi z}{d}}\right)]{\displaystyle \frac{_x|A^{}|^2}{\eta _2}},`$ (20)
$`v_y^{(2)}=`$ $`(S_ES_z){\displaystyle \frac{d^2}{4\pi ^2}}\mathrm{cos}^2\left({\displaystyle \frac{\pi z}{d}}\right){\displaystyle \frac{2_y|A^{}|^2}{\alpha _4}}.`$ (21)
is excited (there is no distinction between “singular” and “non-singular” mean flow in this approach). Excitations of $`\mathrm{\Phi }_0`$ are of the order $`O(\epsilon ^{})`$ and do, as all other remaining corrections, not contribute at leading order.
Projection of the dynamics with the full state vector $`U=S+R`$ onto the slow space yields the equations of motion for the pattern amplitude
$`\tau _tA^{}=[`$ $`\xi _x^2_x^2+\xi _y^2\left(_y^2{\displaystyle \frac{16iq_c}{3\pi }}\phi ^{}_y{\displaystyle \frac{3q_c^2}{4}}\phi _{}^{}{}_{}{}^{2}\right)`$ (22a)
$`+i{\displaystyle \frac{8}{3\pi }}\beta _y_y\phi ^{}+ϵ^{}\left({\displaystyle \frac{3}{4}}g+{\displaystyle \frac{S_E}{ϵ_{}E_0^2}}\right)|A^{}|^2`$ (22b)
$`{\displaystyle \frac{id^2q_c\tau }{48\pi ^2}}{\displaystyle \frac{9S_E+(15+2\pi ^2)S_x9S_z}{\eta _2}}_x|A^{}|^2`$ (22c)
$`+{\displaystyle \frac{id^2q_c\tau }{12\pi ^2}}{\displaystyle \frac{3+\pi ^2}{\eta _2}}_xP^{}]A^{}`$ (22d)
and the twist mode,
$`\gamma _1_t\phi ^{}=`$ $`\left(k_{33}_x^2+k_{11}^{}_y^2\chi _a(H_F^2H_y^2)\right)\phi ^{}`$ (23a)
$`+{\displaystyle \frac{q_c\mathrm{\Gamma }}{2}}\mathrm{Im}\left\{A_{}^{}{}_{}{}^{}\left[{\displaystyle \frac{8}{3\pi }}_y{\displaystyle \frac{3iq_c}{4}}\phi ^{}\right]A^{}\right\}`$ (23b)
$`+{\displaystyle \frac{4d^2}{\pi ^3}}\left[{\displaystyle \frac{\alpha _3}{\eta _2}}+{\displaystyle \frac{2\alpha _2}{\alpha _4}}\right]_x_yP^{}`$ (23c)
$`{\displaystyle \frac{2d^2}{3\pi ^3}}[{\displaystyle \frac{2\alpha _2(S_ES_z)}{\alpha _4}}`$
$`+{\displaystyle \frac{\alpha _3(S_E+4S_xS_z)}{\eta _2}}]_x_y|A^{}|^2,`$ (23d)
and the pressure Poisson equation
$`0=`$ $`{\displaystyle \frac{d^2}{12}}\left[{\displaystyle \frac{1}{\eta _2}}_x^2+{\displaystyle \frac{2}{\alpha _4}}_y^2\right]P^{}`$ (24a)
$`+{\displaystyle \frac{d^2}{8\pi ^2}}{\displaystyle \frac{S_E+(1+\pi ^2/3)S_xS_z}{\eta _2}}_x^2|A^{}|^2`$ (24b)
$`+{\displaystyle \frac{d^2}{4\pi ^2}}{\displaystyle \frac{S_ES_z}{\alpha _4}}_y^2|A^{}|^2.`$ (24c)
### B Remarks on the 2D amplitude equations
By the explicit representation of the coupling coefficients in Eqs. (22-24), the increasing importance of the mean flow contributions in lines (22c, 22d, 23c, 23) as the separation $`d`$ of the damping boundaries increases becomes obvious. On the other hand, for small enough $`d`$, i.e. $`\epsilon ^{}(\lambda /d)^4`$, mean flow is negligible. With some rescaling (indicated by a caret $`\stackrel{ˇ}{\text{ }}`$) dynamics are then described by the system
$`\begin{array}{cc}\hfill \stackrel{ˇ}{\tau }_{\stackrel{ˇ}{t}}\stackrel{ˇ}{A}=[& 1+_{\stackrel{ˇ}{x}}^2+_{\stackrel{ˇ}{y}}^22ic_1\stackrel{ˇ}{\phi }_{\stackrel{ˇ}{y}}c_2\stackrel{ˇ}{\phi }^2\hfill \\ & |\stackrel{ˇ}{A}|^2+i\stackrel{ˇ}{\beta }_{\stackrel{ˇ}{y}}\stackrel{ˇ}{\phi }_{,y}]\stackrel{ˇ}{A},\hfill \end{array}`$ (25a)
$`\begin{array}{cc}\hfill _{\stackrel{ˇ}{t}}\stackrel{ˇ}{\phi }=& _{\stackrel{ˇ}{y}}^2\stackrel{ˇ}{\phi }+\stackrel{ˇ}{K}_3_{\stackrel{ˇ}{x}}^2\stackrel{ˇ}{\phi }\stackrel{ˇ}{H}^2\stackrel{ˇ}{\phi }\hfill \\ & +2\stackrel{ˇ}{\mathrm{\Gamma }}\mathrm{Im}\left\{\stackrel{ˇ}{A}^{}(_{\stackrel{ˇ}{y}}i\stackrel{ˇ}{\phi })\stackrel{ˇ}{A}\right\}.\hfill \end{array}`$ (25b)
This is the “normal form” for the dynamics of a pattern coupled to a soft mode with symmetry under the reflections $`[\stackrel{ˇ}{x}\stackrel{ˇ}{x},\stackrel{ˇ}{\phi }\stackrel{ˇ}{\phi },\stackrel{ˇ}{A}\stackrel{ˇ}{A}^{}]`$ and $`[\stackrel{ˇ}{y}\stackrel{ˇ}{y},\stackrel{ˇ}{\phi }\stackrel{ˇ}{\phi }]`$ or $`[\stackrel{ˇ}{y}\stackrel{ˇ}{y},\stackrel{ˇ}{\phi }\stackrel{ˇ}{\phi },\stackrel{ˇ}{A}\stackrel{ˇ}{A}]`$. In the present case $`c_1=c_2=\frac{4}{3}\left(\frac{8}{3\pi }\right)^20.96`$ are fixed by geometric constraints . These values are quite close to the case with an underlying rotation symmetry $`c_1=c_2=1`$ , which turns out to be somewhat singular in its dynamical properties .
With $`H_y^2=0`$ in line (23a), $`\phi ^{}`$ is damped out and becomes one of many higher order corrections. Equations (22, 24) with $`\phi =0`$ are then sufficient, and with $`\epsilon ^{}(\lambda /d)^4`$ the simple Ginzburg-Landau Equation (3) with
$`g^{}={\displaystyle \frac{3}{4}}g+{\displaystyle \frac{S_E}{ϵ_{}E_0^2}}`$ (26)
(the second contribution is numerically small, see Table I) is retained.
When modulations of $`|A|^2`$ along $`x`$ are strong, like in the chevron pattern, it is instructive to redefine the pressure $`P^{}P^{}\mathrm{const}.\times |A^{}|^2`$ such that its excitation by $`_x^2|A^{}|^2`$ in line (24b) is canceled. It turns out that the remaining excitation of $`P^{}`$ by $`_y^2|A^{}|^2`$ is proportional to $`S_x`$. The substitute for line (22c) incorporating the correction in $`P^{}`$ does then have the form
$`{\displaystyle \frac{id^2q_c\tau (\pi ^26)}{16\pi ^4}}{\displaystyle \frac{S_E+S_xS_z}{\eta _2}}_x|A^{}|^2`$
and accounts for a mean flow with a nontrivial flow profile with zero $`z`$ average.
### C Pressure vs. singular mean flow
As for any incompressible, Newtonian fluid, Eq. (24) is of the form
$`0=_\stackrel{}{r}𝐌_\stackrel{}{r}P^{}_\stackrel{}{r}\stackrel{}{V}.`$ (27)
In the present simple case, the matrix $`𝐌`$ is constant and, in the usual coordinates, diagonal. The inhomogeneity $`\stackrel{}{V}`$ (with dimensions of velocity) depends on time, $`x`$, and $`y`$. Equation (27) can be solved formally by the transformation
$`_\stackrel{}{r}P^{}=𝐌^1\left(\widehat{z}\times _\stackrel{}{r}G+\stackrel{}{V}\right),`$ (28)
which requires a fundamental solution $`G(x,y,t)`$ satisfying
$`0=`$ $`\widehat{z}\times _\stackrel{}{r}_\stackrel{}{r}P^{}`$
$`=`$ $`(\widehat{z}\times _\stackrel{}{r})𝐌^1(\widehat{z}\times _\stackrel{}{r})G+(\widehat{z}\times _\stackrel{}{r})𝐌^1\stackrel{}{V}.`$ (29)
With some rearrangements, Eq. (V C) has the same simple structure as Eq. (27). The quantity $`G`$ can be interpreted as a stream function generating a certain component of the large-scale variations of the velocity field $`(v_x,v_y)(\widehat{z}\times _\stackrel{}{r})G`$, the “singular mean flow” (in early works expressed in terms of the vertical vorticity $`^2G`$). By using Eq. (V C) instead of Eq. (27) and eliminating the pressure $`P^{}`$ via Eqs. (28) also in the remaining equations \[in our case (22,23)\], a description fully in terms of the singular mean flow is obtained.
In principle the three forms (27,28,V C) are equivalent, although the last is often preferred in the literature (for an overview see both Refs. ), perhaps because in some situations with high symmetry $`G`$ is not excited, while $`P^{}`$ is. When disregarding the effects of the additional soft mode $`\phi ^{}`$, application of the less obvious but direct method of Kaiser and Pesch leads to the same result for the mean flow equation (V C) as the route described here.
The flux $`\stackrel{}{V}`$ in Eq. (28) is determined by Eq. (27) only up to a transformation $`\stackrel{}{V}\stackrel{}{V}+(\widehat{z}\times _\stackrel{}{r})G_0(x,y,t)`$ which implies a redefinition $`GGG_0`$. Hence $`(\widehat{z}\times _\stackrel{}{r})G`$ is not generally proportional to the $`z`$ average of the large scale flow. This is not necessary for the formalism to work. In simple cases, like the present, the stream function $`G`$ has this property with the “natural” choice of $`\stackrel{}{V}`$. To guarantee it in general, the method of Newell, Passot and Souli can be used to derive Eq. (28) (Eq. (2.59) in ) directly under this additional constrain.
Here, following Ref. , the formulation as a mass conservation equation (27) is used, since it derives naturally from the general formalism and is thus easily extended to three spatial dimension, to additional homogeneous soft modes, or to relax the assumption of incompressibility (which is not essential for the relevance of mean flow as is sometimes suggested).
In principle, terms proportional to $`_x^2\phi ^{}`$, $`_y^2\phi ^{}`$ could also appear in the $`P^{}`$ equation (24). It is a particularity of the system considered here that they do not. In the conduction regime of EC in cells with homeotropic boundary conditions it can be shown that there is such a, presumably small, excitation of $`P^{}`$ by $`\phi ^{}`$ proportional to the dynamic flexoelectric coefficient $`\zeta ^E`$.
### D Variation of the boundary conditions
Of course, the reduction from 3D to 2D can also be carried out for other boundary conditions than (6,11). Some variants are of practical interest.
An effect similar to applying a magnetic field along $`y`$ can be obtained by a homeotropic (normal) anchoring of the director at the boundaries: For negative, not too small $`ϵ_a`$, and with external electric fields as required for dielectric EC, a homeotropic director alignment is unstable and instead the director becomes planarly oriented everywhere, except for small boundary layers of thickness $`(k_{11}/ϵ_a)^{1/2}E_0^1\lambda `$. Hence, on the length scales $`\lambda `$ relevant for the amplitude formalism, the boundary layer vanishes and instead free boundary conditions for $`\phi `$ can be assumed. For this setup only some geometry factors have to be changed in Eqs. (10-24).
Another variant is a twisted cell, e.g. with $`\phi =0`$ at $`z=d/2`$ and $`\phi =\pi /2`$ at $`z=d/2`$, as it was recently investigated experimentally by Bohatsch and Stannarius . Based on the reduced 3D description derived here, the linear theory for this configuration is developed in Ref. .
### E Higher order contributions
It is worth noticing that, when only the limit of small $`\epsilon ^{}`$ is considered while keeping $`d`$ fixed, there are several other nonlinear and higher order gradient terms besides those in lines (22c,22d) which are formally of the same order of magnitude. A longer list containing more than fifty terms has been calculated numerically by Kaiser and Pesch for the conduction regime of EC. Their form correctly predicts the stability of ideal roll patterns close to threshold but, since it does not separate mean flow and director effects, is unable to describe important effects like the transition to abnormal rolls (see Sec. VI A). These limitations are partly overcome in the less systematic but numerically surprisingly accurate description of Plaut and Pesch .
In fact, since the lowest order mean flow effects entering Eqs. (22-24) all depend on gradients of $`|A^{}|^2`$, they do not contribute to long-wavelength instabilities of the band center (i.e. $`A\mathrm{const}.`$). For the calculation of the thresholds of long-wavelength instabilities of the pattern it may be useful to formally set up amplitude equations including higher order mean flow. But for a systematic quantitative description of general non-ideal patterns containing structures of size $`\lambda ϵ^{1/2}`$ only the truncation used here is justified. When higher order mean flow becomes relevant the 2D amplitude formalism is already breaking down because, for example, the coherence length $`\xi _y\epsilon ^{1/2}`$ becomes of the order of the sample thickness .
## VI Stability of ideal patterns
Rather than deriving stability bounds of ideal patterns $`A=a(z)\mathrm{exp}(iqx+ipy)`$ using a reduced 2D description, it seems more appropriate to do the calculations directly based on the 3D equations, in particular when $`H_y=0`$. This is easily seen from the fact that there is a (numerically small) manifest deviation of $`a(z)`$ from the $`\mathrm{cos}(\pi z/d)`$ profile at the threshold of all instabilities calculated below, indicating that the expansion for small $`A^{}`$ is breaking down. Nevertheless, in order to obtain analytic estimates, Galerkin approximations will be used which correspond effectively to Eqs. (22-24) and their extension to higher order contributions. For simplicity, the stability analysis shall here be performed only at the band center $`q,p=0`$ and only take homogeneous and long-wavelength instabilities with modulations along $`y`$ into account. The latter restriction is justified by experimental observations and by the fact that, with this geometry, the advection of the patter by mean flow is particularly strong. It is then sufficient to consider only the interaction of $`\phi `$, $`v_x`$, and of the phase $`\theta `$ given by $`A=a(z)\mathrm{exp}i\theta (y,z,t)`$ \[$`a=a(z)=O(\epsilon ^{}/g^{})^{1/2}`$ be real and given by Eqs. (10,39b,39c)\], which leads to the following linear problem:
$`\begin{array}{cc}\hfill \gamma _1_t\phi =& \left[\left(k_{11}+2be_{}^2/ϵ_a\right)_y^2+k_{22}_z^2\right]\phi \hfill \\ & \alpha _3_yv_x+(q_c\mathrm{\Gamma }/2)a^2(_y\theta q_c\phi ),\hfill \end{array}`$ (30a)
$`\begin{array}{cc}\hfill \rho _m_tv_x=& \eta _2\left(_y^2+_z^2\right)v_x+\alpha _3_y_t\phi \hfill \\ & +S_{yy}a^2_y(_y\theta q_c\phi )\hfill \\ & +S_{zz}_z\left(a^2_z\theta \right),\hfill \end{array}`$ (30b)
$`\begin{array}{cc}\hfill \tau _t\theta =& \tau q_cv_x+\beta _y_y\phi +\xi _y^2_y^2\theta \hfill \\ & +\xi _z^2a^2_z\left(a^2_z\theta \right).\hfill \end{array}`$ (30c)
Since $`a=0`$ at $`z=\pm d/2`$, the singular last term in Eq. (30c) implies boundary conditions $`_z\theta =0`$. The other boundary conditions are $`\phi ,v_x=0`$ at $`z=\pm d/2`$.
### A Homogeneous destabilization
First, consider homogeneous ($`_y=0`$) destabilizations of the pattern. In this case $`\phi `$ decouples from $`v_x`$ and $`\theta `$. The destabilization of $`\phi `$ is known as the abnormal-roll instability and was investigated in the dielectric regime in Ref. . It was found that
$`\epsilon ^{}=\epsilon _{AR}^{}:={\displaystyle \frac{8\pi ^2g^{}k_{33}}{3d^2q_c^2\mathrm{\Gamma }}},`$ (31)
which can be derived from Eqs. (22,23), is typically a good approximation of the threshold. In MBBA $`\epsilon _{AR}^{}2.8\times (\lambda /d)^2`$ ($`\omega \tau _08`$). The value obtained by numerically calculating eigenmodes of Eq. (30a) directly is $`3\%`$ lower than the value of Eq. (31). Below it will be shown that for MBBA the abnormal-roll instability is precede by a long-wavelength modulation instability. Nevertheless, some phenomena associated with abnormal rolls might be observable around $`\epsilon =\epsilon _{AR}`$, e.g. the tendency of defects in the convection pattern to cluster along lines parallel to the rolls.
For the discussion of homogenous perturbations of $`v_x`$ and $`\theta `$ notice first the neutral mode associated with a translation of the pattern $`\theta \theta +\text{const.}`$. This mode is best dealt with by decomposing $`\theta `$ as $`\theta =\stackrel{~}{\theta }(z,t)+\mathrm{\Theta }(t)`$ such that the $`z`$ average $`\stackrel{~}{\theta }_z`$ vanishes. By multiplying Eq. (30c) by $`a^2`$ and integrating over $`z`$ the dynamics of $`\mathrm{\Theta }`$ is obtained as
$`_t\mathrm{\Theta }=q_c{\displaystyle \frac{a^2v_x_z}{a^2_z}}.`$ (32)
This describes the advection of the pattern by mean flow at large pattern amplitudes.
At the threshold of instability one has $`_t\stackrel{~}{\theta },_tv_x=0`$ and $`_t\mathrm{\Theta }=\text{const.}`$ ($`_t\mathrm{\Theta }0`$ implies an acceleration instability of pattern and liquid crystal). To calculate the threshold, eliminate $`\stackrel{~}{\theta }`$ from Eq. (30b) by Eqs. (30c,32), obtaining the equation
$`0=\eta _2_z^2v_x+{\displaystyle \frac{S_{zz}\tau q_c}{\xi _z^2}}a^2\left[v_x{\displaystyle \frac{a^2v_x_z}{a^2_z}}\right]`$ (33)
from which the critical mode can be determined numerically. When using the low-amplitude approximation $`a^2=(\epsilon ^{}/g^{})\mathrm{cos}^2(\pi z/d)`$ the critical mode is found to be antisymmetric in $`z`$ (i.e. $`_t\mathrm{\Theta }=0`$) at
$`\epsilon ^{}=\epsilon _{\text{drift}}^{}=73.3{\displaystyle \frac{\eta _2g^{}\xi _z^2}{d^2S_{zz}\tau q_c}}.`$ (34)
\[With the large amplitude approximation $`a^2=\epsilon /g`$, the first symmetric and antisymmetric mode both become unstable at the *same* $`\epsilon =(4\pi ^2\eta _2g\xi _z^2)/(d^2S_{zz}\tau q_c)`$.\] The antisymmetric excitation of the phase $`\stackrel{~}{\theta }`$ involved in this instability is obviously inaccessible to a reduced 2D description. Since $`S_{zz}<0`$ in MBBA, the value of $`\epsilon _{\text{drift}}^{}100\times (\omega \tau _0)^1(\lambda /d)^2`$ is negative and the instability does not occur. However, since the electric contribution to $`S_{zz}`$ \[the second term in formula (A24)\] is always positive and comparable in size with the hydrodynamic one, a positive $`S_{zz}`$ is thinkable for other materials. When $`S_{zz}<0`$, the mechanism leading to Eq. (33) is stabilizing – in particular for all perturbations of $`\stackrel{~}{\theta }`$: The lamellae of EC are forced to align normal to the boundaries. This explains why the experimental shadowgraph images, which average the pattern along $`z`$, remain quite sharp even for complicated pattern dynamics.
### B Modulation instabilities
The stability of ideal patterns with respect to perturbations in $`\phi `$, $`v_x`$, and $`\theta `$ modulated with small wave numbers $`k`$ along $`y`$ was investigated numerically using the system (30). It was found that for MBBA the dominating destabilizing feedback loop is based on the excitation of $`v_x`$ by the term containing $`S_{yy}`$ in Eq. (30b) and advection of the phase. There is a good Galerkin approximation for the numerical results. Using the low amplitude approximation for $`a^2`$, Galerkin modes $`\phi ,v_x\mathrm{cos}(\pi z/d)`$ and $`\theta 1`$, and projectors $`\mathrm{cos}(\pi z/d)_z`$ on Eqs. (30a,30b) and $`\mathrm{cos}^2(\pi z/d)_z`$ on Eq. (30c) to reduce the system (30) to algebraic equations, the threshold for long-wavelength modulations instabilities is estimated to be at
$`\begin{array}{cc}\hfill \epsilon ^{}& =\epsilon _{ZZ}^{}:={\displaystyle \frac{72\pi ^4\eta _2g^{}k_{33}\xi _y^2}{d^2q_c}}\times \hfill \\ & \left[512k_{33}S_{yy}\tau \eta _2\mathrm{\Gamma }\left(256\beta _y+27\pi ^2q_c\xi _y^2\right)\right]^1.\hfill \end{array}`$ (35)
With $`\beta _y=q_c\xi _{yy}^2`$ the parenthesis following $`\mathrm{\Gamma }`$ nearly vanish; the remainder is related to the small deviation of $`c_1`$ in Eq. (25a) from unity. Numerically $`\epsilon _{ZZ}^{}9\times (\omega \tau _0)^1(\lambda /d)^2`$ (observe that $`\epsilon _{ZZ}^{}`$ and $`\epsilon _{AR}^{}`$ have different frequency dependence). In particular, at $`\omega \tau _0=8`$, using the second lowest Fourier approximation without flexoelectric effects, Eq. (35) yields $`\epsilon _{ZZ}^{}=0.792(\lambda /d)^2`$ while the numerical solution of system (30) gives a threshold at $`\epsilon ^{}=0.797(\lambda /d)^2`$. Using the Galerkin approximation it is easily verified that the instability is in fact of the long-wavelength ($`k0`$) type.
## VII Qualitative transitions
For the case $`H_F=0`$, order of magnitude estimates shall be used to distinguish regions of qualitatively different pattern dynamics in parameter space - above as well as below the stability bounds of ideal patterns.
As mentioned before, simple Ginzburg Landau dynamics can be expected for small $`ϵ`$ until the last two lines in Eq. (22) become relevant. Assuming $`_x\epsilon _{}^{}{}_{}{}^{1/2}\xi _x^1`$ and $`|A^{}|^2\epsilon ^{}/g^{}`$ , these terms have an effect of the magnitude of $`\epsilon ^{}A^{}`$ when, say, $`\epsilon ^{}[3\pi ^2\eta _2gq_c\xi _x\lambda ^2]^2/[(\pi ^26)(S_E+S_xS_z)\tau d^2]^22.\times 10^4\times (\lambda /d)^4(\omega \tau _0)^2`$ (MBBA, $`\omega \tau _010`$). Preliminary simulations of Eqs. (22,24) with $`\phi ^{}0`$ show that for higher $`\epsilon ^{}`$ defect cores (where $`_x|A^{}|^2`$ is large) are strongly deformed and lines along which the phase $`\mathrm{arg}A^{}`$ “jumps” are often generated and long living (rather long living phase jump lines are also observed experimentally and in simulations of a similar model ; it is not clear, though, whether these are due to lowest order mean flow effects). But the simulations also indicate that these lowest order mean flow effects, although they are formally dominating over the direct nonlinear saturation via the last term in line (22b), do not prevent the system form finally reaching a steady state with $`|A^{}|^2\epsilon ^{}/g^{}`$.
Thus, assuming still $`|A^{}|^2\epsilon ^{}/g^{}`$, there will be a further transition at $`\epsilon ^{}[3^{1/2}\pi (\alpha _4+2\eta _2)gk_{33}q_c^3\xi _y/(128\alpha _2S_x)]^{2/3}\times (\lambda /d)^{8/3}1.\times (\lambda /d)^{8/3}`$ (MBBA, $`\omega \tau _08`$) where contributions from slaved excitations of the in-plane director $`\phi `$ by mean flow become relevant in the 2D pattern dynamics. Semiquantitatively these effects are described by Eq. (23) with $`H_y=0`$, but a restriction of dynamics to a single linear mode of $`\phi `$ is then not justified.
When, with increasing $`ϵ`$, horizontal length scales become of the size of the sample thickness $`d`$ \[at $`\epsilon ^{}(\pi \xi _x/d)^20.2\times (\lambda /d)^2`$ (MBBA), say\] the 2D description breaks down. Because then $`d`$ is not the dominating length scale for the damping of mean flow anymore, the trend in the influence of mean flow on the smallest structures in the $`A`$ field (e.g. defects) is reversed. Now the structures themselves set the length scale. Assuming $`_x\epsilon ^{1/2}\xi _x^1`$ and $`|A|^2\epsilon /g`$, the contribution $`i\tau v_xA`$ in Eq. (10) is large compared to $`\epsilon A`$ up to $`\epsilon (q_cS_x\tau \xi _x/\eta _2g)^21.\times 10^6\times (\omega \tau _0)^2`$ (MBBA, $`\omega \tau _010`$). For larger $`\epsilon `$ the cores of defects are not affected by lowest order mean flow effects. Larger structures, like the phase field of the pattern or variations in the defect density (e.g. in chevron patterns), may still be.
Current experimental resolutions are of the order $`\mathrm{\Delta }\epsilon =10^3\mathrm{}10^2`$. The estimates above suggest that lowest order mean flow effects are best observed near the upper limit of the frequency range for the validity of the hydrodynamic description used here, at $`\omega \tau _0=O(10^2)`$ (see Section II C). With $`d5\lambda `$ they should be observable in the range of validity of the 2D description.
According to the model for the chevron mechanism , chevrons depend essentially on the abnormal-roll mechanism and can only form above the abnormal-roll instability bound \[here given by Eq. (31)\]. Thus, it is plausible to assume a $`\epsilon ^{}(\lambda /d)^2`$ threshold for chevron formation, in accordance with measurements presented in Ref. , where an approximate $`\omega ^1`$ frequency dependence is found.
For very high $`\epsilon ^{}`$ many authors report the formation of disclination loops, which, being singularities in the director field, indicate already a breakdown of *hydrodynamics*. In thicker cells, the chevron pattern might decay along other routes before the 3D amplitude formalism breaks down at some $`\epsilon =O(1)`$.
For convection in most quasi-2D systems, i.e. systems with $`d/\lambda =O(1)`$, all these transitions, from the breakdown of simple 2D Ginzburg-Landau dynamics to the breakdown of the amplitude formalism, do, in principle, collapse at $`\epsilon =O(1)`$. Only by specially designed experiments (as in ) can these transitions be unfolded.
## VIII Electric Nusselt numbers
Recently “electric Nusselt numbers” $`𝒩_r`$, $`𝒩_i`$ have been introduce by Gleeson, Gheorghiu, Plaut as the ratio of the in-phase or, respectively, out-of-phase components of the electric current to the corresponding values expected for the unstructured basic state at a given voltage, minus one. As for the Nusselt number in thermal convection, they are to first approximation proportional to $`|A|^2`$.
Several interesting questions can be addressed by measuring electric Nusselt numbers. First, it follows from the discussion above that, in dielectric EC, $`|A|^2\epsilon /g`$ for thick enough cells even far inside the three-dimensionally chaotic range (the average of $`|A|^2`$ across the pattern is typically only weakly reduced in the presence of defects; see, e.g., Ref. ). Measurements of the electric Nusselt numbers therefore seem to be an effective method to test an essential feature of the theory. Second, Nusselt number measurements may also help to identify the qualitative changes in the dynamics predicted in Section VII, in particular since they do, in contrast to optical methods, not loose their sensitivity in thick cells or with small wavenumbers. Finally, the frequency dependence of the Landau coefficient $`g^{}`$, which enters the Nusselt numbers near threshold in a simple way, provides information on the strength of the dynamic flexoelectric effect (see Appendix A). From the derivation of $`g^{}`$ it is clear that boundary effects do not interfere in these measurements. To obtain the theoretical value for the Nusselt numbers, take the $`x`$ and $`y`$ average (symbol $`_{xy}`$) of Eq. (39b) plus $`i`$ times Eq. (39c), which leads to
$$\begin{array}{c}(\sigma _{}+i\omega ϵ_{})_z^2\mathrm{\Phi }_r+i\mathrm{\Phi }_i_{xy}=\hfill \\ \hfill (I_r+iI_i)_z|A|^2_{xy}\\ \hfill +(I_{rz}+iI_{iz})_z\mathrm{Im}\{A^{}_xA\}_{xy}.\end{array}$$
(36)
Taking the boundary conditions (6,11) into account, this implies, similar as for $`\mathrm{\Phi }_r^{(2)}+i\mathrm{\Phi }_i^{(2)}`$ in Eq. (17),
$$\begin{array}{c}\frac{1}{2}j_z:=(\sigma _{}+i\omega ϵ_{})_z\mathrm{\Phi }_r+i\mathrm{\Phi }_i_{xy}=\hfill \\ \hfill (I_r+iI_i)|A|^2_{xyz}\\ \hfill +(I_{rz}+iI_{iz})\mathrm{Im}\{A^{}_xA\}_{xyz}\end{array}$$
(37)
at $`z=\pm d/2`$. Obviously, $`j_z`$ is the complex amplitude of the average, pattern-induced electric current density through the sample. (The last term on the r.h.s. represents a correction due to a global deviation of the average wavenumber from the critical one, and will be dropped below.) Thus,
$`𝒩_r={\displaystyle \frac{2I_r|A|^2_{xyz}}{E_0\sigma _{}}}`$ (38a)
and
$`𝒩_i={\displaystyle \frac{2I_i|A|^2_{xyz}}{E_0ϵ_{}\omega }}.`$ (38b)
At threshold, $`(d/d\epsilon ^{})|A|^2_{xyz}=(2g^{})^1`$, e.g. $`d𝒩_r/d\epsilon ^{}=0.31`$ in MBBA at $`\omega \tau _0=8`$, dropping flexoelectric contributions. The value increases roughly $`\omega `$ as frequency increases. Typically, it seems to be larger than the corresponding value for the conduction regime .
## IX Conclusion
It has been shown how the 3D dielectric convection pattern interacts with various homogeneous soft modes, which are related to undamped hydrodynamic modes. The method to establish these relations is not unique, but the “center manifold” method seems to be favorable over the “order parameter” method.
The reduction of the 3D pattern dynamics to a quasi-2D form in the layer geometry was derived analytically, thus establishing a description of the interaction of the pattern with the twist mode and the pressure field (or singular mean flow). Scaling analysis suggests that the transition from the simple, quasi-2D Ginzburg-Landau dynamics to manifestly 3D dynamics in thick layers unfolds into several well distinguished steps, the first of which occurs already very close to threshold ($`ϵ^{}=O(\lambda /d)^4`$). These characteristics should generally be expected for 3D patterns.
Ideal, dielectric EC patterns are found to destabilize at some value $`\epsilon ^{}(\omega \tau _0)^1(\lambda /d)^2`$ for which an analytic approximation in terms of material parameters is given. A particular nonlinear mechanism that stabilizes the phase of the pattern to be constant along $`z`$, thus giving the pattern a 2D appearance also at higher $`\epsilon ^{}`$, is identified in Sec. VI. As outlined in Sec. VIII, measurements of the electric Nusselt numbers are suitable for quantitatively testing the theory, probing the dynamic flexoelectric effect in nematic liquid crystals independent of boundary effects, and investigating the route of the transition from simple 2D to fully 3D dynamics.
It is my pleasure to thank L. Kramer, Y. Kuramoto, A. Lindner, W. Pesch, and E. Plaut for valuable discussion, W. Decker and W. Pesch for providing the basic nematodynamic equations in a computer readable form, the Kyoto University for its hospitality and the Japan Foundation for the Promotion of Science (P98285) and the Ministry of Education, Science, Sports and Culture in Japan for their support.
$`\begin{array}{cc}\hfill 0=& \left[\left(\sigma _{}\underset{¯}{2b\sigma _a}\right)_x^2+\sigma _{}\left(_y^2+_z^2\right)\right]\mathrm{\Phi }_0\hfill \\ & \underset{¯}{2be_{}(\sigma _a/ϵ_a)_x_y\phi }\hfill \end{array}`$ (39a)
$`\begin{array}{cc}\hfill 0=& +E_0_x^2\left[ϵ_{}\omega \mathrm{\Phi }_i(\sigma _{}\underset{¯}{\sigma _a})\mathrm{\Phi }_r\right]\hfill \\ & +E_0\left(_y^2+_z^2\right)\left(ϵ_{}\omega \mathrm{\Phi }_i\sigma _{}\mathrm{\Phi }_r\right)\hfill \\ & \left[\underset{¯}{\sigma _a/ϵ_a}+\left(1\underset{¯}{3b}\right)E_0^2ϵ_a/2\gamma _1\right]\left(\underset{¯}{e_+}2\gamma _1\zeta ^E\right)_x^2_z\mathrm{\Phi }_0\hfill \\ & +[\underset{¯}{(\sigma _a/ϵ_a)\left(k_{11}+k_{22}\right)}\hfill \\ & +e_{}E_0^2(1\underset{¯}{3b})(\underset{¯}{e_+/2\gamma _1}\zeta ^E)]_x_y_z\phi \hfill \\ & +\underset{¯}{\left(\sigma _a/ϵ_a\right)_x\left(\alpha _3_zv_x+\alpha _2_xv_z\right)}\hfill \\ & +E_0\left(I_r_z|A|^2+I_{rx}_x\mathrm{Im}\{A^{}_zA\}+I_{rz}_z\mathrm{Im}\{A^{}_xA\}\right)\hfill \end{array}`$ (39b)
$`\begin{array}{cc}\hfill 0=& E_0_x^2\left[\sigma _{}\mathrm{\Phi }_i+(ϵ_{}\underset{¯}{ϵ_a})\omega \mathrm{\Phi }_r\right]\hfill \\ & E_0\left(_y^2+_z^2\right)\left(\sigma _{}\mathrm{\Phi }_i+ϵ_{}\omega \mathrm{\Phi }_r\right)\hfill \\ & \underset{¯}{\omega \left(12b\right)\left(e_+2\gamma _1\zeta ^E\right)_x^2_z\mathrm{\Phi }_0}\hfill \\ & +\underset{¯}{\left[\omega \left(k_{11}+k_{22}\right)(2be_{}/ϵ_a)\left(e_+2\gamma _1\zeta ^E\right)\right]_x_y_z\phi }\hfill \\ & +\underset{¯}{\omega _x\left(\alpha _3_zv_x+\alpha _2_xv_z\right)}\hfill \\ & +E_0\left(I_i_z|A|^2+I_{ix}_x\mathrm{Im}\{A^{}_zA\}+I_{iz}_z\mathrm{Im}\{A^{}_xA\}\right)\hfill \end{array}`$ (39c)
$`\begin{array}{cc}\hfill \gamma _1_t\phi =& \left[k_{33}_x^2+\left(k_{11}+\underset{¯}{2be_{}^2/ϵ_a}\right)_y^2+k_{22}_z^2+\chi _aH_y^2\right]\phi \hfill \\ & +\left(e_+2\gamma _1\zeta ^E\underset{¯}{2be_{}}\right)_x_y\mathrm{\Phi }_0\hfill \\ & \alpha _3_yv_x\alpha _2_xv_y+(q_c\mathrm{\Gamma }/2)\mathrm{Im}\{A^{}(_yiq_c\phi )A\}\hfill \end{array}`$ (39d)
$`\begin{array}{cc}\hfill \rho _m_tv_x=& \left(\alpha _1+\alpha _4+\alpha _5+\alpha _6\right)_x^2v_x+\eta _2\left(_y^2+_z^2\right)v_x\hfill \\ & +\left(\alpha _2+\eta _1\right)_x\left(_yv_y+_zv_z\right)\hfill \\ & +_x\left[PE_0ϵ_{}_z\mathrm{\Phi }_r+2\alpha _3\zeta ^E\left(_y^2+_z^2\right)\mathrm{\Phi }_0\right]\hfill \\ & +\alpha _3_y_t\phi +S_x_x|A|^2+S_{xx}_x\mathrm{Im}\{A^{}_xA\}\hfill \\ & +S_{yy}_y\mathrm{Im}\{A^{}(_yiq_c\phi )A\}+S_{zz}_z\mathrm{Im}\{A^{}_zA\}\hfill \end{array}`$ (39e)
$`\begin{array}{cc}\hfill \rho _m_tv_y=& \left(\eta _1_x^2+\alpha _4_y^2+\alpha _4_z^2/2\right)v_y\hfill \\ & +_y\left[\left(\alpha _2+\eta _1\right)_xv_x+\alpha _4_zv_z/2\right]\hfill \\ & +_y\left[PE_0ϵ_{}_z\mathrm{\Phi }_r+2\alpha _2\zeta ^E_x^2\mathrm{\Phi }_0\right]+\alpha _2_x_t\phi \hfill \\ & +S_{xy}_x\mathrm{Im}\{A^{}(_yiq_c\phi )A\}\hfill \end{array}`$ (39f)
$`\begin{array}{cc}\hfill \rho _m_tv_z=& \left[\left(\underset{¯}{\alpha _2}+\eta _1\right)_x^2+\alpha _4\left(_y^2/2+_z^2\right)\right]v_z\hfill \\ & +_z\left[\left(\underset{¯}{\alpha _3}+\alpha _2+\eta _1\right)_xv_x+\alpha _4_yv_y/2\right]\hfill \\ & +_z\left[P\underset{¯}{\left(k_{11}k_{22}+2be_+e_{}/ϵ_a\right)_x_y\phi }\right]\hfill \\ & \left[\underset{¯}{e_+2be_+}2(\underset{¯}{\gamma _1}+\alpha _2)\zeta ^E\right]_x^2_z\mathrm{\Phi }_0\hfill \\ & E_0\left[(ϵ_{}\underset{¯}{ϵ_a})_x^2+ϵ_{}_y^2+2ϵ_{}_z^2\right]\mathrm{\Phi }_r+S_z_z|A|^2\hfill \\ & +S_{xz}_x\mathrm{Im}\{A^{}_zA\}+S_{zx}_z\mathrm{Im}\{A^{}_xA\}\hfill \end{array}`$ (39g)
$`\begin{array}{cc}\hfill 0=& \stackrel{}{v}\hfill \end{array}`$ (39h)
(Note: is confused here, so we continue with a single column)
## A Results for linear stability and coupling coefficients in 3D
Some analytic and numeric results relating the 3D description of the pattern dynamics to hydrodynamics are presented here, in particular analytic approximations for all coefficients entering the the results in Sections V, VI, and VIII. In the second lowest Fourier approximation the onset of dielectric EC is at
$`E_c^2=`$ $`{\displaystyle \frac{4\omega \sigma _{}\left(\alpha _2^2\gamma _1\eta _1\right)}{X\left(\alpha _2{\displaystyle \frac{2\omega \sigma _{}\left(\alpha _2ϵ_{}+ϵ_a\eta _1\right)}{4ϵ_{}^2\omega ^2+\sigma _{}^2}}\right)ϵ_a\eta _1\sigma _{}}}`$ (A1)
with a critical wavenumber
$$\begin{array}{c}q_c^2=\frac{2\omega \left(\alpha _2^2+\gamma _1\eta _1\right)}{k_{33}\eta _1}\hfill \\ \hfill \times \frac{2ϵ_{}^2\omega ^2+\sigma _{}^2{\displaystyle \frac{2ϵ_aϵ_{}ϵ_{}\eta _1\omega ^2\sigma _{}}{\alpha _2Xϵ_a\eta _1\sigma _{}}}}{4ϵ_{}^2\omega ^2+\sigma _{}^2{\displaystyle \frac{2X\omega \sigma _{}\left(\alpha _2ϵ_{}+ϵ_a\eta _1\right)}{\alpha _2Xϵ_a\eta _1\sigma _{}}}},\end{array}$$
(A2)
where $`X:=ϵ_{}\sigma _aϵ_a\sigma _{}`$. The accurate numerical result for $`E_c^2`$ vs. $`\omega `$ is nearly a perfect straight line (see Fig. 2), which can probably be understood by means of the approximation used in . Formula (A1) nicely estimates the offset of this line for intermediate $`\omega \tau _0`$, but gives a different slope as $`\omega \tau _0\mathrm{}`$. The deviations at the lower end are an artifact occurring with all truncated Fourier approximations. For $`q_c^2`$ the situation is similar.
In the second lowest Fourier approximation, the excitation of the hydrodynamic fields in the convection pattern near threshold (critical mode) is
$`\begin{array}{cc}\hfill \mathrm{\Phi }=& |A|\mathrm{sin}(q_cx+\mathrm{arg}A)\hfill \\ & \times \left(\mathrm{\Phi }_u+2\mathrm{\Phi }_c\mathrm{cos}(2\omega t)2\mathrm{\Phi }_s\mathrm{sin}(2\omega t)\right),\hfill \end{array}`$ (A3)
$`\begin{array}{cc}\hfill n_z=& 2|A|\mathrm{cos}(q_cx+\mathrm{arg}A)\left(\mathrm{cos}(\omega t)+\mathrm{sin}(\omega t)\right),\hfill \end{array}`$ (A4)
$`\begin{array}{cc}\hfill f=& 2|A|\mathrm{sin}(q_cx+\mathrm{arg}A)\hfill \\ & \times \left(f_c\mathrm{cos}(\omega t)f_s\mathrm{sin}(\omega t)\right),\hfill \end{array}`$ (A5)
where $`f`$ generates a velocity field $`\stackrel{}{v}=(_x_z,_y_z,_x^2_y^2)f`$ and the real constants $`f_c`$, $`f_s`$, $`\mathrm{\Phi }_u`$, $`\mathrm{\Phi }_c`$, $`\mathrm{\Phi }_s`$ are given by
$$\begin{array}{c}\hfill 4\eta _1f_cq_c^3=4\alpha _2\omega +\frac{XE_c^2\left(8ϵ_{}^2\omega ^22ϵ_{}\omega \sigma _{}+3\sigma _{}^2\right)}{4ϵ_{}^2\omega ^2\sigma _{}+\sigma _{}^3},\end{array}$$
(A6)
$$\begin{array}{cc}\multicolumn{2}{c}{\begin{array}{cc}\hfill 4\alpha _2f_sq_c^3=& ϵ_aE_c^24\gamma _1\omega +4k_{33}q_c^2\hfill \\ & +\frac{ϵ_aE_c^2\left[2ϵ_a\omega \left(2ϵ_{}\omega \sigma _{}\right)+\sigma _a\left(2ϵ_{}\omega +\sigma _{}\right)\right]}{4ϵ_{}^2\omega ^2+\sigma _{}^2},\hfill \end{array}}\end{array}$$
(A7)
$$\begin{array}{c}\hfill \mathrm{\Phi }_u=\frac{E_c\sigma _a}{q_c\sigma _{}},\end{array}$$
(A8)
$$\begin{array}{cc}\hfill \mathrm{\Phi }_c+i\mathrm{\Phi }_s=& \frac{\left(1i\right)E_c\left(2iϵ_a\omega +\sigma _a\right)}{2q_c(2iϵ_{}\omega +\sigma _{})}.\hfill \end{array}$$
(A9)
With the help of the real constants $`\mathrm{\Phi }_u^+`$, $`\mathrm{\Phi }_c^+`$, $`\mathrm{\Phi }_s^+`$, $`n_z^+`$ which characterize the adjoint eigenvector
$$\begin{array}{c}\hfill \mathrm{\Phi }_u^+=\frac{E_c\left(\alpha _2ϵ_{}+ϵ_a\eta _1\right)q_c}{2\alpha _2\sigma _{}},\end{array}$$
(A10)
$$\begin{array}{cc}& \mathrm{\Phi }_c^++i\mathrm{\Phi }_s^+=\frac{q_c}{4\alpha _2XE_c\omega }\times \hfill \\ & \{+E_c^2[\alpha _2X+ϵ_a\eta _1(2iϵ_{}\omega +\sigma _{})]\hfill \\ & 4[(\alpha _2^2\gamma _1\eta _1)\omega +k_{33}\eta _1q_c^2](2iϵ_{}\omega +\sigma _{})\},\hfill \end{array}$$
$$\begin{array}{c}\hfill n_z^+=\frac{\eta _1q_c^2}{\alpha _2}\end{array}$$
(A11)
and the normalization factor
$`4N=`$ $`\mathrm{\hspace{0.17em}2}(q_c^2n_z^+)ϵ_aE_c^2`$
$`+E_cq_c[(ϵ_an_z^+ϵ_{}q_c^2)(\mathrm{\Phi }_u+\mathrm{\Phi }_c+\mathrm{\Phi }_s)`$ (A12)
$`2ϵ_a\omega (\mathrm{\Phi }_c^++\mathrm{\Phi }_s^+)\sigma _a(2\mathrm{\Phi }_u^++\mathrm{\Phi }_c^+\mathrm{\Phi }_s^+)],`$
analytic results for most coupling coefficients entering the 3D description can be obtained:
$`\begin{array}{cc}\hfill 2N\tau =& ϵ_aE_c\left(2\mathrm{\Phi }_u^++\mathrm{\Phi }_c^+\mathrm{\Phi }_s^+\right)q_c2ϵ_{}\left(\mathrm{\Phi }_u^+\mathrm{\Phi }_u+\mathrm{\Phi }_c^+\mathrm{\Phi }_c+\mathrm{\Phi }_s^+\mathrm{\Phi }_s\right)q_c^2,\hfill \end{array}`$ (A13)
$`\begin{array}{cc}\hfill E_c^2\xi _x^2=& \\ \hfill 8k_{33}\eta _1\{& 6k_{33}\eta _1q_c^2\sigma _{}(4ϵ_{}^2\omega ^2+\sigma _{}^2)+E_c^2[ϵ_a\eta _1(2ϵ_{}ϵ_{}\omega ^2\sigma _{}+2ϵ_{}^2\omega ^2\sigma _{}+\sigma _{}^2\sigma _{})\alpha _2(2ϵ_{}^2\omega ^2+\sigma _{}^2)X]\}\hfill \\ \hfill \times \{& 8\left[2ϵ_{}k_{33}\eta _1\left(\alpha _2ϵ_{}ϵ_a\eta _1\right)\omega ^2q_c^2+\left(\alpha _2ϵ_{}+ϵ_a\eta _1\right)\left(\alpha _2^2\gamma _1\eta _1\right)\omega ^2\sigma _{}+\alpha _2k_{33}\eta _1q_c^2\sigma _{}^2\right]X+3\alpha _2^2E_c^2\sigma _{}X^2\hfill \\ \hfill 8ϵ_ak_{33}\eta _1^2q_c^2& (4ϵ_{}^2\omega ^2+\sigma _{}^2)\sigma _{}+E_c^2[ϵ_a^2\eta _1^2(4ϵ_{}^2\omega ^2\sigma _{}+8ϵ_{}ϵ_{}\omega ^2\sigma _{}+3\sigma _{}\sigma _{}^2)2\alpha _2ϵ_a\eta _1(4ϵ_{}ϵ_{}\omega ^2+3\sigma _{}\sigma _{})X]\}^1,\hfill \end{array}`$ (A14)
$`\begin{array}{cc}\hfill 2N\xi _y^2=& ϵ_{}E_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c+\mathrm{\Phi }_s\right)q_c\left(f_c+f_s\right)\alpha _4q_c^3\hfill \\ & 4ϵ_{}\omega \left(\mathrm{\Phi }_s^+\mathrm{\Phi }_c\mathrm{\Phi }_c^+\mathrm{\Phi }_s\right)2\sigma _{}\left(\mathrm{\Phi }_u^+\mathrm{\Phi }_u+\mathrm{\Phi }_c^+\mathrm{\Phi }_c+\mathrm{\Phi }_s^+\mathrm{\Phi }_s\right)+\mathrm{flexo},\hfill \end{array}`$ (A15)
$`\begin{array}{cc}\hfill 2N\xi _z^2=& ϵ_{}E_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c+\mathrm{\Phi }_s\right)q_c+4\alpha _3\omega 2\left(f_c+f_s\right)\left(\alpha _1+\alpha _3+\alpha _4+\alpha _5+(\alpha _3/\alpha _2)\eta _1\right)q_c^3\hfill \\ & 4ϵ_{}\omega \left(\mathrm{\Phi }_s^+\mathrm{\Phi }_c\mathrm{\Phi }_c^+\mathrm{\Phi }_s\right)2\sigma _{}\left(\mathrm{\Phi }_u^+\mathrm{\Phi }_u+\mathrm{\Phi }_c^+\mathrm{\Phi }_c+\mathrm{\Phi }_s^+\mathrm{\Phi }_s\right)+\mathrm{flexo},\hfill \end{array}`$ (A16)
$`\begin{array}{cc}\hfill (8/3)Ng=& (2q_c^24n_z^+)ϵ_aE_c^2+4q_c^2\left[\omega \left(\alpha _3+\gamma _2\right)2ϵ_an_z^+\mathrm{\Phi }_u\mathrm{\Phi }_c\right]\hfill \\ & +\left(f_c+f_s\right)\left[\left(6\alpha _2+4\alpha _3\right)n_z^+4\left(\alpha _1+\gamma _2\right)q_c^2\right]q_c^32E_cq_c\sigma _a\left(3\mathrm{\Phi }_u^++\mathrm{\Phi }_c^+2\mathrm{\Phi }_s^+\right)\hfill \\ & +ϵ_aE_cq_c\left[\left(5n_z^+2q_c^2\right)\left(\mathrm{\Phi }_u+2\mathrm{\Phi }_c\right)4\omega \left(2\mathrm{\Phi }_c^++\mathrm{\Phi }_s^+\right)\right]+8\omega q_c^2ϵ_a\left(\mathrm{\Phi }_c^+\mathrm{\Phi }_u+2\mathrm{\Phi }_s^+\mathrm{\Phi }_c2\mathrm{\Phi }_c^+\mathrm{\Phi }_s\right)\hfill \\ & +4q_c^2\sigma _a\left[2\mathrm{\Phi }_u^+\left(\mathrm{\Phi }_u\mathrm{\Phi }_s\right)+2\mathrm{\Phi }_c^+\mathrm{\Phi }_c+\mathrm{\Phi }_s^+\left(2\mathrm{\Phi }_s\mathrm{\Phi }_u\right)\right]+\mathrm{flexo},\hfill \end{array}`$ (A17)
$`\begin{array}{cc}\hfill 4I_r+4iI_i=& \left(iϵ_a\omega +\sigma _a\right)\left[\left(42i\right)E_c\left(22i\right)q_c\left(\mathrm{\Phi }_u+i\mathrm{\Phi }_c\mathrm{\Phi }_s\right)\right]\hfill \\ & +q_c^3ϵ_aE_c\left[\left(3+i\right)f_c+\left(1i\right)f_s\right]+2ϵ_{}q_c^4\left[f_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c+i\mathrm{\Phi }_s\right)+if_s\left(\mathrm{\Phi }_u\mathrm{\Phi }_ci\mathrm{\Phi }_s\right)\right],\hfill \end{array}`$ (A18)
$`\begin{array}{cc}\hfill 4I_{rx}+4iI_{ix}=& {\displaystyle \frac{2\left(iϵ_a\omega +\sigma _a\right)}{ϵ_aE_c}}\times \hfill \\ & \left[4\gamma _2\left(f_cf_s\right)q_c^2+2ϵ_aE_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c\mathrm{\Phi }_s\right)ϵ_a\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)q_c+4q_c\left(k_{33}k_{11}\right)\right]\hfill \\ & (22i)\left(iϵ_a\omega +\sigma _a\right)\left(\mathrm{\Phi }_u+i\mathrm{\Phi }_c\mathrm{\Phi }_s\right)q_c^2ϵ_aE_c\left[\left(3+i\right)f_c+(1i)f_s\right]\hfill \\ & 2ϵ_{}q_c^3\left[f_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c+i\mathrm{\Phi }_s\right)+if_s\left(\mathrm{\Phi }_u\mathrm{\Phi }_ci\mathrm{\Phi }_s\right)\right]+\text{flexo},\hfill \end{array}`$ (A19)
$`\begin{array}{cc}\hfill \mathrm{\Gamma }=& \mathrm{\hspace{0.17em}4}k_{22}4k_{33}ϵ_a(E_c/q_c)\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c\mathrm{\Phi }_s\right)+ϵ_a\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)2\alpha _3q_c\left(f_cf_s\right)+\text{flexo}.\hfill \end{array}`$ (A20)
$`\begin{array}{cc}\hfill 2S_x=& q_c\left[2q_c^2\left(\alpha _1+\alpha _5+\gamma _2\right)\left(f_cf_s\right)ϵ_aE_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c\mathrm{\Phi }_s\right)+q_cϵ_{}\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)\right],\hfill \end{array}`$ (A21)
$`\begin{array}{cc}\hfill 2S_z=& \mathrm{\hspace{0.17em}2}\alpha _5q_c^3\left(f_cf_s\right)+ϵ_aE_c\left[2E_cq_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c\mathrm{\Phi }_s\right)\right],\hfill \end{array}`$ (A22)
$`\begin{array}{cc}\hfill 2S_{yy}=& \left(2\eta _2\alpha _4\right)\left(f_cf_s\right)q_c^2+q_cϵ_{}\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)+\text{flexo},\hfill \end{array}`$ (A23)
$`\begin{array}{cc}\hfill 2S_{zz}=& \mathrm{\hspace{0.17em}2}\left(\gamma _2\alpha _1\right)\left(f_cf_s\right)q_c^2+q_cϵ_{}\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)+\text{flexo},\hfill \end{array}`$ (A24)
$`\begin{array}{cc}\hfill 2S_{xy}=& \left(\alpha _2+\alpha _5\right)\left(f_cf_s\right)q_c^2ϵ_aE_c\left(\mathrm{\Phi }_u+\mathrm{\Phi }_c\mathrm{\Phi }_s\right)+q_cϵ_{}\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)+\text{flexo},\hfill \end{array}`$ (A25)
$`\begin{array}{cc}\hfill 2S_{xz}=& \mathrm{\hspace{0.17em}2}\left(\gamma _2\alpha _1\right)\left(f_cf_s\right)q_c^2+q_cϵ_{}\left(\mathrm{\Phi }_u^2+2\mathrm{\Phi }_c^2+2\mathrm{\Phi }_s^2\right)+4q_c\left(k_{33}k_{11}\right)+\text{flexo},\hfill \end{array}`$ (A26)
Here “flexo” stands for flexoelectric corrections, which involve matrix inversions and are hard to express in a compact form. For similar reasons, no formulas are given for $`\kappa _x`$, $`S_{xx}`$, $`S_{zx}`$, $`I_{rz}`$ and $`I_{iz}`$. The coefficients $`\kappa _z`$ vanishes for the second lowest Fourier approximation (and $`\rho _m=0`$) and $`\alpha _x`$ and $`\beta _y+q_c\xi _y^2`$ have only flexoelectric contributions. The expression for $`\xi _x^2`$ has a different structure than the other formulas because it was not calculated with the “center manifold” formalism but by differentiating the determinantal condition for stability of the basic state, which is more effective in this case. There are indirect contributions from excitations of $`n_z`$ entering $`S_{xz}`$, $`I_{ix}`$ and $`I_{rx}`$ \[the bracket in Eq. (A19), in Eq. (A26) they cancel favorably\]. These have been calculated only in the lowest Fourier approximation (i.e. for the non-oscillatory part of $`n_z`$), which introduces a small error.
The quality of these results can be judged by comparing the values obtained for MBBA at $`\omega \tau _0=8`$ with the accurate numerical values in Table I (only the approximation for the indirect contribution from $`n_z`$ entering $`S_{xz}`$, $`I_{ix}`$ and $`I_{rx}`$ is retained). In the combinations in which the results are expressed there, they are, for fixed $`\omega \tau _0`$, independent of the electric conductivity and, with the exception of $`\kappa _x,\kappa _z`$, $`I_{rx}`$, and some contributions involving $`\zeta ^E`$, also largely independent of $`\omega \tau _0`$ for $`\omega \tau _0>8`$. Two exceptions, which should both be experimentally accessible, shall be highlighted: the dynamic flexoelectric contribution to $`E_c\alpha _x`$, which increases linear in $`\omega \tau _0`$, and the contribution to $`g`$ proportional to $`\zeta _{}^{E}{}_{}{}^{2}`$ which is the only one which increases $`(\omega \tau _0)^2`$ as $`\omega \tau _0\mathrm{}`$. The result for $`g`$ in the lowest Fourier approximation
$$\begin{array}{c}g=9.45+0.00252e_+^20.00401\omega \tau _0e_+\zeta ^E\hfill \\ \hfill 5.18\times 10^5\omega \tau _0\zeta _{}^{E}{}_{}{}^{2}+0.00102\omega ^2\tau _0^2\zeta _{}^{E}{}_{}{}^{2}\end{array}$$
(A27)
illustrates the latter effect, although it is correct only in its order of magnitude. The flexoelectric contributions in Eq. (A27) and also in Table I are expressed in terms of $`e_+:=e_1+e_3`$, $`e_{}:=e_1e_3`$, and $`\zeta ^E:=2\gamma _1\zeta ^E`$, in units of $`10^{12}\mathrm{Cm}^1`$ ($`3.00\times 10^5\mathrm{dyn}^{1/2}`$ in Gaussian units). Typical values measured for $`e_1`$ and $`e_3`$ are a few times that much (see the overviews in Refs. ).
As a result of the approximate $`\pi /4`$ phase shift of the director oscillations, $`\kappa _z`$ is so small that the remaining finite viscous effect is comparable in size to the effect of finite mass density $`\rho _m`$, which has been suppressed everywhere else.
## B “Order parameter” vs. “center manifold” method
Here, two general methods for obtaining amplitude equations are compared. It is shown that they give different results in the presence of multiple homogeneous soft modes. In this work, the “center manifold” method is used to derive the reduced equations. In order to keep the formalism simple, it will be restricted to homogeneous ($`\stackrel{}{q}=0`$) modes and their slow modulations. The inclusion of patterning soft modes ($`|\stackrel{}{q}|0`$) is straight forward.
### 1 Formal setting
Let the state vector $`U(\stackrel{}{r})`$ describe the configuration of all relevant degrees of freedom (e.g. hydrodynamic fields) of the system in the (ideally) infinitely extended, $`D`$-dimensional $`\stackrel{}{r}`$ space. Assume the “microscopic equations” to be of the form
$$0=F(U)=L(_t,)U(\stackrel{}{r})+\text{nonlinear terms},$$
(B1)
where the linear operator $`L(_t,)`$ is polynomial in $`_t`$ and $``$, acting on $`U(\stackrel{}{r})`$ locally and translation invariant in space and time. There are several branches $`j`$ of linear modes $`V_j(\stackrel{}{q})`$ which solve the generalized eigenvalue problem
$$L(\sigma _j(\stackrel{}{q}),i\stackrel{}{q})V_j(\stackrel{}{q})=0,$$
(B2)
and, with some suitable scalar product $`|`$ (which does *not* contain an integration over $`\stackrel{}{r}`$), adjoint eigenstates
$$W_j(\stackrel{}{q})L(\sigma _j(\stackrel{}{q}),i\stackrel{}{q})=0.$$
(B3)
For some branches $`jK`$ the growth rates $`\text{Re}\{\sigma _j(\stackrel{}{q})\}`$ vanish (or are small) at their maxima at $`\stackrel{}{q}=0`$, for the others ($`jK`$) they are negatively large in the vicinity of $`\stackrel{}{q}=0`$.
### 2 The order parameter method
Using the “order parameter” method physical states are characterized by weighted sums over slow eigenfunctions of $`L(_t,)`$,
$$U(\stackrel{}{r})=\underset{jK}{}\underset{\mathrm{\Omega }}{}u_j(\stackrel{}{q})V_j(\stackrel{}{q})\mathrm{exp}(i\stackrel{}{q}\stackrel{}{r})𝑑\stackrel{}{q}+R(\{u_k\}).$$
(B4)
The weights $`u_j(\stackrel{}{q})`$ are interpreted as the Fourier transforms of the set of “amplitudes” used in the reduced description. The range of integration $`\mathrm{\Omega }`$ is a region around $`\stackrel{}{q}=0`$, large enough to include all significant contributions from $`u_j(\stackrel{}{q})`$ and small enough to exclude slow (patterning) modes at large wavenumbers. With this ansatz slaved contributions $`R(\{u_k\})`$, which are fully in the fast eigenspace of $`L`$, come in only at nonlinear order. The linear dynamics for each amplitude is simply given by
$$_tu_j(\stackrel{}{q})=\sigma _j(\stackrel{}{q})u_j(\stackrel{}{q}).$$
(B5)
An inverse Fourier transform yields the linear dynamics in physical space, which is usually simplified by truncating the Taylor expansion of $`\sigma _j(\stackrel{}{q})`$ in each component of $`\stackrel{}{q}`$ for small $`|\stackrel{}{q}|`$, such that, in physical space, derivatives of $`u_j`$ are obtained.
However, the situation is different in the case of multiple slow branches. Then $`\sigma _j(\stackrel{}{q})`$ is typically non-analytic in the components of $`\stackrel{}{q}`$ (although it is analytic in $`|\stackrel{}{q}|`$). As a generic example, consider the linear operator
$$L(_t,)=\left(\begin{array}{cc}_x^2+_y^2_t& _x_y\\ _x_y& _x^2+_y^2a_t\end{array}\right)$$
(B6)
with a positive parameter $`a`$. There are two neutral modes at $`(q,p):=\stackrel{}{q}=0`$. It is easily seen that one of the two growth rates is of the form
$$\sigma _1(q,p)=\frac{p^4+p^2q^2+q^4}{p^2+q^2}+O(a),$$
(B7)
i.e., non-analytic at $`q,p=0`$. As a result, the corresponding amplitude equation in physical space
$$\begin{array}{c}_tu_1(\stackrel{}{r})=\frac{7}{8}^2u_1(\stackrel{}{r})\hfill \\ \hfill +K(\stackrel{}{r}^{}\stackrel{}{r})u_1(\stackrel{}{r}^{})d^2r^{}+O(a)+o(u_1,u_2)\end{array}$$
(B8)
is nonlocal. In polar coordinates $`K(\stackrel{}{r})=(3/\pi )r^4\mathrm{cos}4\phi `$. These conclusions do not require $`a`$ to be small, because the additional terms $`O(a)`$ depend on $`a`$ and they can cancel the non-analyticity calculated here at most at particular values of $`a`$.
This transition from local basic equations to amplitude equations with *algebraically* decaying non-localities is counter-intuitive and misleading. This approach has the advantage that in Fourier space the linear dynamics (B5) is simple. This is useful for calculations of patterns stability involving only a few Fourier modes.
Finally, notice that a general method to reobtain local amplitude equations from the Fourier representation (B5), e.g. by redefining the amplitudes, should not be expected. Equation (B5) is general enough to include even non-local interactions in the basic equations, which certainly cannot lead to local amplitude equations.
### 3 The center manifold method
Alternatively, in the “center manifold” method only the slow modes at $`\stackrel{}{q}=0`$, $`V_j(0)`$, are used for the characterization of the physical state
$$U(\stackrel{}{r})=\underset{jK}{}\underset{\mathrm{\Omega }}{}u_j(\stackrel{}{q})V_j(0)\mathrm{exp}(i\stackrel{}{q}\stackrel{}{r})𝑑\stackrel{}{q}+R(\{u_k\}).$$
(B9)
The “slow subspace” spanned by the sum in Eq. (B9) can be extracted by the projection operator
$$P:=\underset{jK}{}P_j,$$
(B10)
with
$$\begin{array}{c}P_jf(\stackrel{}{r}):=\underset{\mathrm{\Omega }}{}𝑑\stackrel{}{q}\frac{d\stackrel{}{r^{}}}{(2\pi )^D}\hfill \\ \hfill W_j(0)|f(\stackrel{}{r}^{})V_j(0)\mathrm{exp}(i\stackrel{}{q}(\stackrel{}{r}\stackrel{}{r^{}})),\end{array}$$
(B11)
where it is assumed without loss of generality that the states $`W_j(0)`$ and $`V_j(0)`$ ($`jK`$) entering $`P`$ form a bi-orthonormal system. The “slaved” contributions $`R(\{u_k\})`$ cover the remaining subspace.
The factor $`V_j(0)`$ in Eq. (B9) can be pulled out of the integral, which is then simply the inverse Fourier transform of $`u_j(\stackrel{}{q})`$ into physical space $`u_j(\stackrel{}{r})`$. It is thus justified to define $`u_j(\stackrel{}{r})`$ as (the local average of) the hydrodynamic variable $`W_j|U(\stackrel{}{r})`$. In particular, if $`W_j|U(\stackrel{}{r})`$ is conserved, so is $`u_j(\stackrel{}{r})`$. In such a local representation, the vicinity of the “center manifold” method to a multiple scale approximation would be more obvious. Similar simplifications would also be possible for the integrals below, but have been suppressed in order to ease the comparison with the “order parameter” method.
The function $`R(\{u_k\})`$ is defined by the perturbative solution of
$$(1P)F(U)=0,$$
(B12)
where $`U`$ is given by Eq. (B9) and the $`u_k`$ are small and vary slowly and smoothly in space and time but are otherwise arbitrary.
At linear order in $`U`$, where Eq. (B12) reduces to
$$(1P)L(_t,)U=0,$$
(B13)
a form
$$R(\{u_k\})=\underset{jK}{}\underset{\mathrm{\Omega }}{}r_j(\stackrel{}{q})V_j(0)\mathrm{exp}(i\stackrel{}{q}\stackrel{}{r})𝑑\stackrel{}{q}$$
(B14)
with contributions $`r_j(\stackrel{}{q})`$ only in the vicinity of $`q=0`$ is sufficient.
The amplitude equations are then given by
$$P_iL(_t,)U=0\text{for each }\text{i∈K}$$
(B15)
where $`R`$ is eliminated from $`U`$ via Eq. (B13).
When Eqs. (B13,B15) are satisfied with slowly varying $`u_k`$ and small $`R(\{u_k\})`$, this implies that $`U`$ contains no fast eigenvectors of $`L(_t,)`$. Hence the resulting linear dynamics for $`U`$ is the same as the one obtained with the “order parameter” method, in particular $`R(\{u_k\})`$ is then given by Eq. (B14) with
$$r_j(\stackrel{}{q})=\underset{k,l}{}W_k(\stackrel{}{q})|V_j(0)^1W_k(\stackrel{}{q})|V_l(0)u_l(\stackrel{}{q}),$$
(B16)
where $`lK`$ is running over all fast modes and $`j,kK`$ are running over all slow modes (for small $`|\stackrel{}{q}|`$ the matrix $`W_k(\stackrel{}{q})|V_j(0)`$ in this expression is generally a perturbed unit matrix and readily inverted).
To see that this method yields local dynamics for the amplitudes, split the linear operator, restricted to the subspace selected by $`Q:=(1P)`$, like
$$QL(_t,)Q=\underset{:=L_0}{\underset{}{QL(0,0)Q}}+l(_t,)$$
(B17)
into a part $`L_0`$ which is regular, and a term which is small for slow temporal and spatial variations of the operand and polynomial in $`_t,`$. Calling the sum on the r.s.h. of Eq. (B9) $`S(\{u_k\})`$, and suppressing the arguments (such that $`U=S+R`$), equation (B13) becomes
$$QLR=QLQR=(L_0+l)R=QLS,$$
(B18)
and is solved by expanding for small $`l`$, i.e.
$$R=\underset{n=0}{\overset{\mathrm{}}{}}\left(L_0^1l(_t,)\right)^nL_0^1QL(_t,)S.$$
(B19)
When eliminating $`R`$ from Eq. (B15) by Eq. (B19) and truncating at some power in the derivatives (i.e., for slow enough variations), linear amplitude equations with local interactions are obtained. The extension to the nonlinear level is straightforward (s. Section V A).
It should be noticed that with this approach all modes in the kernel of $`L(0,0)`$ have to be treated as “soft modes”, some of which, e.g. those resulting from gauge symmetries, may not actually have slowly relaxing modulations associated with them (see e.g. the pressure mode in Sec. IV B). For the simple example (B6) the “center manifold” method leads to amplitude equations identical to the basic equations.
The reason for the difference between the two approaches is that in the multidimensional kernel of $`L(0,0)`$ the choice of the basis vectors characterizing the slow modes is not unique. While they are fixed (with respect to the hydrodynamic variables in the “microscopic equations”) for the “center manifold” method, they point, depending on $`\stackrel{}{q}`$, into arbitrary directions in the slow space for the “order parameter” method. For the same reason, there is no-near identity transformation mapping one representation onto the other. |
warning/0001/hep-th0001024.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Although N=1 supersymmetric field theories in $`3+1`$ dimensions have been extensively investigated for more than twenty five years, most of these investigations have been based on the standard supersymmetry algebra. It has been known for some time, however, that $`p`$-brane solitons in supersymmetric theories carry $`p`$-form charges that appear as central charges in the spacetime supertranslation algebra . Allowing for all such charges, the D=4 N=1 supertranslation algebra is spanned by a four component Majorana spinor charge $`Q`$, the 4-vector $`P_\mu `$ and a Lorentz 2-form charge $`Z_{\mu \nu }`$. The only non-trivial relation is the anticommutator
$$\{Q,Q\}=C\gamma ^\mu P_\mu +\frac{1}{2}C\gamma ^{\mu \nu }Z_{\mu \nu },$$
(1)
where $`C`$ is the charge conjugation matrix and $`\gamma _\mu =(\gamma _0,\gamma _i)`$ are the four Dirac matrices. Our metric convention is ‘mostly plus’ so that we may choose a real representation of the Dirac matrices. In this representation the Majorana spinor charges $`Q`$ are real, so $`\{Q,Q\}`$ is a symmetric $`4\times 4`$ matrix with a total of ten real entries. The number of components of $`P_\mu `$ and $`Z_{\mu \nu }`$ is also ten, so that we have indeed included all possible bosonic central charges. Note that the automorphism group of this algebra is $`GL(4;\text{})`$.
The components of $`Z_{\mu \nu }`$ can be interpreted as charges carried by domain walls , while $`P_\mu `$ is (in general) a linear combination of the momentum and a string charge. In the case of a domain wall, the tension is bounded by the charge, and saturation of this bound implies preservation of 1/2 of the N=1 D=4 supersymmetry. This is one example in the class of ‘1/2 supersymmetric’ configurations allowed by the supersymmetry algebra<sup>1</sup><sup>1</sup>1An analysis of 1/2 supersymmetric combinations of charges in $`N>1`$ D=4 theories, $`N=2`$ in particular, can be found in . Such 1/2 supersymmetric domain walls were shown to occur in in the Wess-Zumino (WZ) model, for an appropriate superpotential, and also arise in the $`SU(n)`$ SQCD because the low-energy effective Lagrangian is that of a WZ model with a superpotential admitting $`n`$ discrete vacua . More recently, it was shown that the WZ model also admits (again for an appropriate superpotential) 1/4 supersymmetric configurations that can be interpreted as intersecting domain walls . More precisely, it was established that such configurations must solve a certain ‘Bogomol’nyi’ equation for which earlier mathematical studies had made the existence of appropriate solutions plausible (especially in view of the results of which were recently brought to our attention). Domain wall junctions of the WZ model have since been studied further in and an explicit 1/4 supersymmetric domain wall junction of a related model has recently been found .
It was pointed out in that the possibility of 1/4 supersymmetric intersecting domain walls is inherent in the supersymmetry algebra. If we choose $`C=\gamma ^0`$ and $`\gamma _5=\gamma ^0\gamma ^1\gamma ^2\gamma ^3`$, then (1) becomes
$$\{Q,Q\}=H+\gamma ^{0i}P_i+\frac{1}{2}\gamma ^{0ij}U_{ij}+\frac{1}{2}\gamma ^{0ij}\gamma _5V_{ij}$$
(2)
where $`H=P^0`$, $`U_{ij}=Z_{ij}`$ and $`V_{ij}=\epsilon _{ijk}Z_{0k}`$. One is thus led to expect ‘electric’ type domain walls with non-zero 2-form $`U_{ij}`$ but vanishing $`V_{ij}`$ and ‘magnetic’ type domain walls with non-zero 2-form $`V_{ij}`$ but vanishing $`U_{ij}`$. In general, a domain wall will be specified not only by its tension and orientation but also by an angle in the electric-magnetic charge space; the domain wall is ‘dyonic’ when this angle is not a multiple of $`\pi `$. It is not difficult to show that the algebra (2) allows for dyonic charge configurations preserving 1/4 supersymmetry. In this paper we determine the model-independent restrictions on such configurations that are implied by the supersymmetry algebra.
As pointed out in , the charge associated with the stringlike junctions of domain walls in the WZ model appears in the supersymmetry algebra in the same way as the 3-momentum, so for a static 1/4 supersymmetric configuration of the WZ model the 3-vector $`𝐏`$ must be interpreted as a string charge carried by the domain wall junction. It was further shown in that this junction charge contributes positively to the energy of the 1/4 supersymmetric configuration as a whole. In contrast, the charge associated to domain wall junctions of the model considered in was shown there to contribute negatively to the total energy. As we shall see, this apparent disagreement is due to a different central charge structure for the two models. There is therefore more than one field theory realization of static intersecting domain walls preserving 1/4 supersymmetry, but as yet no example that exploits the most obvious possibility in which $`𝐏`$ vanishes but $`U_{ij}`$ and $`V_{ij}`$ do not.
These observations underscore the importance of the model-independent analysis of 1/4 supersymmetric configurations based only on the N=1 D=4 supersymmetry algebra, but our aim is to understand the implications of the supersymmetry algebra for all supersymmetric configurations, not just those preserving 1/4 supersymmetry. Since the matrix $`\{Q,Q\}`$ is a positive definite real symmetric one, it can be brought to diagonal form with real non-negative eigenvalues. The number of zero eigenvalues is the number of supersymmetries preserved by the configuration. The ‘supersymmetric’ configurations are those for which this number is $`1,2,3`$ or $`4`$. There is a unique ‘vacuum’ charge configuration preserving all four supersymmetries. Configurations preserving two supersymmetries are 1/2 supersymmetric while those preserving one supersymmetry are 1/4 supersymmetric. Configurations preserving three supersymmetries are 3/4 supersymmetric, but there is no known field theoretic realization of this possibility. Indeed, we will show here that there is no classical field configuration of the WZ model that preserves 3/4 supersymmetry. However, possible string-theoretic realizations of exotic supersymmetry fractions such as 3/4 supersymmetry were recently explored in , and this possibility has been considered previously in a variety of other contexts . In particular, the $`OSp(1|8;\text{})`$-invariant superparticle model of provides a simple realization in the context of particle mechanics. The fundamental representation of $`OSp(1|8;\text{})`$ is spanned by $`(\rho ^\alpha ,\lambda _\alpha ,\zeta )`$, where $`\rho `$ and $`\lambda `$ are two 4-component real commuting spinors of $`Spin(1,3)`$, and $`\zeta `$ is a real anticommuting scalar. The action
$$S=𝑑t\left[\rho ^\alpha \dot{\lambda }_\alpha +\zeta \dot{\zeta }\right]$$
(3)
is manifestly $`OSp(1|8)`$ invariant; in particular, it is supersymmetric with supersymmetry charge $`Q=\lambda \zeta `$. The canonical (anti)commutation relations imply that $`\{Q_\alpha ,Q_\beta \}=\lambda _\alpha \lambda _\beta `$, which is a matrix of rank one, corresponding to 3/4 supersymmetry.
Thus, there exist models of one kind or another in which all possible fractions of D=4 N=1 supersymmetry are preserved. This fact provides further motivation for the general model-independent analysis of the possibilities allowed by the supersymmetry algebra that we present here. As we shall explain, the space of supersymmetric charge configurations, or ‘BPS states’, is the boundary of the convex cone of $`4\times 4`$ real symmetric matrices and this has an interpretation in terms of Jordan algebras. In analogy with the way that the conformal group acts on massless states on the light-cone $`P^2=0`$, there is a group $`Sp(8,\text{})`$ that acts on the ‘BPS cone’ of supersymmetric configurations and which has an interpretation in this context as the Möbius group of the Jordan algebra . Another purpose of this paper is to explore some of the geometrical ideas underlying this interpretation of supersymmetric charge configurations.
It is generally appreciated that BPS states are stable states, this being the main reason for their importance, but some “standard” arguments for stability rely on physical intuition derived from special cases. For example, a massive charged particle that minimises the energy for given charge cannot radiate its energy away in the form of uncharged photons because this would leave behind a particle with the same charge but lower energy, contradicting the statement that the original particle minimised the energy in its charge sector. However, this heuristic argument is not conclusive. For instance, the stability against radiative relaxation to a lower energy state of the same ‘charge vector’ assumes that the radiated energy carries away no momentum because momentum is one of the charges, and this assumption would be violated by a decay in which just one photon is emitted. It is also implicit in the heuristic argument that prior to decay one can go to the rest frame, but the supersymmetry algebra allows BPS states for which this is not possible, a massless particle being an obvious, but by no means the only, example. These considerations show that it is not quite as obvious as generally supposed that BPS states are stable. Here we provide a complete analysis, for the general D=4 N=1 supersymmetry algebra, based on a combination of the Minkowski reverse-triangle inequality for positive-definite matrices and the ordinary triangle inequality for BPS energies.
The supertranslation algebra for which (1) is the only non-trivial (anti)commutator is a contraction of the superalgebra $`osp(1|4;\text{})`$, which is the D=4 N=1 anti-de Sitter (adS) superalgebra. The anticommutator of the 4 real supercharges of the latter is
$$\{Q,Q\}=C\gamma ^\mu P_\mu +\frac{1}{2}C\gamma ^{\mu \nu }M_{\mu \nu },$$
(4)
where $`M_{\mu \nu }`$ are the Lorentz generators. This is formally equivalent to (1), although the charges on the right hand side are no longer central because they generate the adS group $`SO(3,2)`$. However the positivity conditions on these charges are the same, as are the conditions for preservation of supersymmetry. This fact means that much of our analysis of the centrally-extended supertranslation algebra can be immediately applied to the adS case. A related analysis has been considered previously for D=5 in , where the D=4 case was briefly mentioned, and BPS states in D=4 adS have also been analysed by other methods in .
We begin with an analysis of the N=1 D=4 supersymmetry algebra, determining the charge configurations that preserve the various possible fractions of supersymmetry, and we show how the positivity of $`\{Q,Q\}`$ implies the stability of BPS states carrying these charges. We also show how the supersymmetry algebra determines, in a model-independent way, some properties of the 1/4 supersymmetric intersecting domain walls that are realized by the WZ model, but show also that 3/4 supersymmetry is not realized by classical WZ field configurations. We then turn to an exposition of the geometry associated with the supersymmetric configurations, which is that of self-dual homogeneous convex cones, and review their relation to Jordan algebras. We then discuss how our results apply to D=4 N=1 adS supersymmetry, and conclude with comments on implications and generalizations of our work, in particular to M-theory.
## 2 BPS states
The anticommutator (2) can be rewritten as
$$\{Q,Q\}=H+\gamma ^{0i}P_i+\gamma _5\gamma ^iU_i+\gamma ^iV_i$$
(5)
where
$$U_i=\frac{1}{2}\epsilon _{ijk}U_{jk}V_i=\frac{1}{2}\epsilon _{ijk}V_{jk}.$$
(6)
As mentioned above, a charge configuration is supersymmetric if the matrix $`\{Q,Q\}`$ has at least one zero eigenvalue. Thus, supersymmetric charge configurations are those for which $`\{Q,Q\}`$ has vanishing determinant. We see from (5) that this determinant must be expressible in terms of $`H`$ and the three 3-vectors $`𝐏`$, $`𝐔`$ and $`𝐕`$.
Now $`det\{Q,Q\}`$ is manifestly $`SL(4;\text{})`$ invariant, but the subgroup that keeps $`H`$ fixed is its maximal compact $`SO(4)[SU(2)\times SU(2)_R]/\text{}_2`$ subgroup. Ignoring the quotient by $`\text{}_2`$, the first $`SU(2)`$ factor can be identified with the 3-space rotation group while the $`SU(2)_R`$ group rotates the three 3-vectors $`𝐏`$, $`𝐔`$ and $`𝐕`$ into each other, i.e. these three 3-vectors form a triplet of $`SU(2)_R`$. The notation chosen here reflects the fact that $`SU(2)_RU(1)_R`$, where $`U(1)_R`$ is the R-symmetry group<sup>2</sup><sup>2</sup>2This symmetry is usually broken in D=4 N=1 QFTs, either by the superpotential or by anomalies. We shall comment on this fact in the conclusions, but it is not relevant to the purely algebraic analysis presented here. rotating $`𝐔`$ into $`𝐕`$ keeping $`𝐏`$ fixed (this is the automorphism group of the standard supersymmetry algebra, including Lorentz generators). We conclude from this that $`det\{Q,Q\}`$ is a fourth-order polynomial in $`H`$ with coefficients that are homogeneous polynomials in the three algebraically-independent $`SU(2)\times SU(2)_R`$ invariants that can be constructed from $`𝐏`$, $`𝐔`$ and $`𝐕`$. These are
$`a`$ $`=`$ $`U^2+V^2+P^2`$
$`b`$ $`=`$ $`𝐏𝐔\times 𝐕`$
$`c`$ $`=`$ $`|𝐔\times 𝐕|^2+|𝐏\times 𝐔|^2+|𝐏\times 𝐕|^2.`$ (7)
An explicit computation shows that
$$det\{Q,Q\}=P(H)$$
(8)
where $`P(H)`$ is the quartic polynomial
$$P(H)=H^42aH^28bH+a^24c$$
(9)
The fact that $`\{Q,Q\}`$ is a positive real symmetric matrix imposes a bound on $`H`$ in terms of the invariants $`a,b,c`$. Specifically,
$$HE(a,b,c)$$
(10)
where $`E(a,b,c)`$ is the largest root of $`P(H)=(H\lambda _1)(H\lambda _2)(H\lambda _3)(H\lambda _4)`$. Since the sum of the roots vanishes, the largest root $`E`$ is necessarily non-negative. The number of supersymmetries preserved is then the number of roots equal to $`E`$. The vacuum configuration has all roots equal with $`E=0`$. In all other cases $`E`$ is strictly positive and the number of roots equal to it is 1,2 or 3, corresponding to 1/4,1/2 or 3/4 supersymmetry.
Our first task, to be undertaken below, is to analyse the conditions required for the realization of each of these possibilities. We will then show how the stability of states preserving supersymmetry, alias ‘BPS states’, is guaranteed by the supersymmetry algebra. Although all model-independent consequences of supersymmetry are encoded in the supersymmetry algebra, the extraction of these consequences for BPS states is facilitated by methods that involve only the constraints on the Killing spinors associated with these states, and we show in the subsequent subsection how these methods can be used to learn about restrictions imposed by the preservation of supersymmetry on intersecting domain walls. We conclude with a discussion of 3/4 supersymmetry, and a proof that this fraction is not realized in the WZ model.
### 2.1 Supersymmetry fractions
The analysis of the conditions on the invariants $`a,b,c`$ required for the preservation of the various possible fractions of supersymmetry is fairly straightforward when the polynomial $`P(H)`$ has at least two equal roots, and is especially simple when there are three equal roots. We shall therefore begin with the case of three equal roots, followed by the case of two equal roots, arriving finally at the generic case.
The quartic polynomial $`P(H)`$ has three equal roots if
$$c=\frac{a^2}{3},b=(\frac{a}{3})^{3/2},$$
(11)
and the roots are
$$\lambda _1=\lambda _2=\lambda _3\lambda =\pm (\frac{a}{3})^{1/2},\lambda _4=3\lambda .$$
(12)
If $`\lambda `$ is positive then we have the BPS bound $`H\lambda `$, and charge configurations saturating this bound preserve 3/4 supersymmetry. If $`\lambda `$ is negative then we instead find the BPS bound $`H3\lambda `$, with only 1/4 supersymmetry being preserved by charge configurations that saturate it.
Charge configurations preserving 1/2 supersymmetry can occur only when $`P(H)`$ has two equal roots. The conditions for the special case in which $`\lambda _1=\lambda _2`$ and $`\lambda _3=\lambda _4`$ are
$`b=c=0`$
$`\lambda _1=\lambda _3=\pm \sqrt{a}`$ (13)
In the more general case when $`\lambda _1=\lambda _2\lambda `$ and $`\lambda _3\rho `$ we have $`\lambda _4=(2\lambda +\rho )`$. If $`\lambda =0`$ we have $`a^2=4c`$, $`b=0`$ and $`\rho ^2=2a`$, with 1/4 supersymmetry when $`H=|\rho |`$. Otherwise we find the condition
$$4a^3b^2+27b^418ab^2ca^2c^2+4c^3=0$$
(14)
with
$$3\lambda ^2=a\pm 2(a^23c)^{1/2},\rho ^2+2\lambda \rho +3\lambda ^2=2a.$$
(15)
with 1/2 supersymmetry possible when $`\lambda `$ is the largest root.
The general case of four unequal roots is quite complicated, unless $`b=0`$, in which case
$$(\lambda _1,\lambda _2,\lambda _3,\lambda _4)=(\sqrt{a+2\sqrt{c}},\sqrt{a2\sqrt{c}},\sqrt{a2\sqrt{c}},\sqrt{a+2\sqrt{c}}).$$
(16)
One way to achieve $`b=0`$ is to set $`𝐏=\mathrm{𝟎}`$. In this case the bound on $`H`$ becomes
$$H\sqrt{U^2+V^2+2|𝐔\times 𝐕|}.$$
(17)
Note that this becomes $`H|𝐔|+|𝐕|`$ when $`𝐔𝐕=\mathrm{𝟎}`$, which is typical of 1/4 supersymmetric orthogonal intersections of branes. The four eigenvalues of $`\{Q,Q\}`$ are, in order of increasing magnitude,
$$H\sqrt{a+2\sqrt{c}},H\sqrt{a2\sqrt{c}},H+\sqrt{a2\sqrt{c}},H+\sqrt{a+2\sqrt{c}}.$$
(18)
The first of these vanishes when the bound is saturated. The last two are never zero unless all four vanish, which is the vacuum charge sector. The second eigenvalue equals the first only when $`c=0`$, so in this case there are two zero eigenvalues when the bound is saturated and we have 1/2 supersymmetry. Otherwise we have 1/4 supersymmetry.
As emphasized earlier, static configurations need not have $`𝐏=\mathrm{𝟎}`$ because $`𝐏`$ may have an interpretation as a domain-wall junction charge, rather than 3-momentum (in general it must be interpreted as a sum of the 3-momentum and a string junction charge). Nevertheless, one may still have $`b=0`$ if $`𝐔\times 𝐕`$ vanishes, which it will do if, say, $`𝐕=0`$. In this case, the results are exactly as in the $`𝐏=\mathrm{𝟎}`$ case just analysed but with $`𝐕`$ replaced by $`𝐏`$. In particular, if $`𝐏𝐔=0`$ we then have $`H|𝐏|+|𝐔|`$, and static 1/4 supersymmetric configurations have $`H|𝐏|+|𝐔|`$. For this case, we can bring the charges to the form
$$𝐏=(0,0,Q),𝐔=(u_1,u_2,0),𝐕=(0,0,0).$$
(19)
where $`Q`$ is a junction charge. This case is the one analysed in , with $`T=u_1+iu_2`$ being the complex scalar charge in the D=3 supersymmetry algebra obtained by dimensional reduction on the 3-direction. In agreement with we find that $`H=|T|+|Q|`$, so the junction charge contributes positively to the energy of the whole configuration.
More generally, we might have
$$𝐏=(0,0,Q),𝐔=(u_1,u_2,0),𝐕=(v_1,v_2,0).$$
(20)
This case was analysed in , and an explicit realization of it was found in a model with several chiral superfields; in this model the charge Q is again associated with a domain wall junction. In agreement with we find the four roots to be
$`\lambda _1`$ $`=`$ $`Q+\sqrt{(u_2+v_1)^2+(u_1v_2)^2}`$
$`\lambda _2`$ $`=`$ $`Q\sqrt{(u_2+v_1)^2+(u_1v_2)^2}`$
$`\lambda _3`$ $`=`$ $`Q\sqrt{(u_2v_1)^2+(u_1+v_2)^2}`$
$`\lambda _4`$ $`=`$ $`Q+\sqrt{(u_2v_1)^2+(u_1+v_2)^2}.`$ (21)
Note that the four roots are distinct, in general, and (in contrast to the previous case) $`b0`$. If $`Q`$ is positive and $`\lambda _1`$ is the largest root, the junction charge $`Q`$ contributes negatively to the total energy as in .
The case just considered is a special case of the larger class of configurations with $`b0`$ for which $`P(H)`$ has four distinct roots. At this point the analysis becomes quite complicated, and we shall not pursue it further.
### 2.2 Stability of BPS states
Our aim in this subsection is to prove the stability of BPS states. We begin by considering the possible decay of a general state, not necessarily BPS, with energy $`H_3`$ into two other states, not necessarily BPS, with energies $`H_1`$ and $`H_2`$. This can be represented schematically as
$$(state)_3(state)_1+(state)_2.$$
(22)
Let us write
$$\{Q,Q\}=H+K(a,b,c),$$
(23)
where $`K`$ is a traceless symmetric matrix, and $`(a,b,c)`$ are the three $`SU(2)\times SU(2)_R`$ invariants introduced previously. Conservation of charges and energy requires that
$`H_3`$ $`=`$ $`H_1+H_2`$ (24)
$`K_3`$ $`=`$ $`K_1+K_2`$ (25)
where $`K_i=K(a_i,b_i,c_i)`$, with $`(a_i,b_i,c_i)`$ being the values of the invariants $`(a,b,c)`$ for the $`i`$th state. Since the matrices $`H_i+K_i`$ are positive definite they are subject to the Minkowski reverse triangle inequality (see e.g. )
$$[det(H_3+K_3)]^{\frac{1}{4}}[det(H_1+K_1)]^{\frac{1}{4}}+[det(H_2+K_2)]^{\frac{1}{4}}.$$
(26)
We now want to see the consequences of supposing state 3 to be BPS. We observe that the left hand side of (26) vanishes if state 3 is BPS, but the right hand side can vanish only if both states 1 and 2 are also BPS. The extension to more than two decay products is immediate so we conclude that any unstable BPS state would have to decay into other BPS states.
To complete the proof of stability we now show that a BPS state cannot decay into other BPS states. A BPS state has an energy $`H=E(K)E(a,b,c)`$ where $`E(K)`$ is the largest value of $`H`$ for which $`det(H+K)=0`$. An equivalent characterization of $`E(K)`$ is as the smallest eigenvalue of $`K`$. It follows that $`E(K)=min\left(\zeta ^TK\zeta \right)`$, where $`\zeta `$ is a commuting spinor normalized such that $`\zeta ^T\zeta =1`$ but otherwise arbitrary. From this and the fact that $`min(a+b)min(a)+min(b)`$, we deduce the triangle inequality
$$E(K_1+K_2)E(K_1)+E(K_2).$$
(27)
Generic models will have a spectrum of BPS states for which this inequality is never saturated. In such cases BPS states are absolutely stable. In those cases for which there are BPS energies saturating the inequality (27) there may be states of marginal stability<sup>3</sup><sup>3</sup>3It is well known that marginal stability is the mechanism by which BPS states ‘decay’ as one moves in the space of parameters defining certain theories, but this is a discontinuity of the BPS spectrum as a function of parameters and not a process within a given theory.. The inequality (27) is saturated when $`K_1`$ and $`K_2`$ are proportional, with positive constant of proportionality, but this is only a sufficient condition for equality. Another sufficient condition, which we believe to be necessary, is the coincidence, up to normalization, of the eigenvectors of $`K_1`$ and $`K_2`$ with lowest eigenvalue.
It is instructive to see how the above comments apply to the special case in which $`H+K=C\gamma ^\mu P_\mu `$. The Minkowski inequality becomes
$$\sqrt{(P_1+P_2)^2}\sqrt{P_1^2}+\sqrt{P_2^2}.$$
(28)
Since $`\sqrt{P^2}`$ is the rest mass $`m`$ of a particle with 4-momentum $`P`$, we learn that
$$m_3m_1+m_2.$$
(29)
This is the familiar rule that the sum of the masses of the decay products cannot exceed the mass of the particle undergoing decay. Given that $`m_3=0`$ we deduce that $`m_1=m_2=0`$, so if a massless particle decays into two other particles those two particles must also be massless. For this special case the triangle inequality (27) reduces to
$$|𝐏_1+𝐏_2||𝐏_1|+|𝐏_2|,$$
(30)
which is saturated if and only if $`𝐏_1`$ and $`𝐏_2`$ are parallel, and in this case there is no phase space for the decay.
### 2.3 Domain Walls at Angles
Each supersymmetric configuration is associated with a set of Killing spinors $`ϵ`$ which span the kernel of $`\{Q,Q\}`$. With the exception of the vacuum configuration, these spinors are subject to constraints that reduce the dimension of the space that they span. Some properties of supersymmetric configurations follow directly from the nature of these constraints. In particular, intersecting brane configurations can be considered as configurations obtained from parallel branes by rotation of one or more of them. The constraints can be similarly obtained, and then analysed to determine the dimension of the space of Killing spinors they allow . We shall apply this analysis here to intersecting domain walls of N=1 D=4 theories.
We begin with two coincident domain walls, corresponding to the constraint
$$\gamma _{013}ϵ=ϵ.$$
(31)
We then rotate one of them around the 3-axis until it makes an angle $`\beta `$ in the 12-plane, and simultaneously rotate by some angle $`\alpha `$ in the electric-magnetic charge space. This operation is represented by the matrix
$$R=e^{\frac{1}{2}\alpha \gamma _5}e^{\frac{1}{2}\beta \gamma _{12}},$$
(32)
which satisfies
$$\gamma _{013}R^1=R\gamma _{013}.$$
(33)
The constraint on the Killing spinor $`ϵ`$ imposed by the rotated brane is
$$R\gamma _{013}R^1ϵ=ϵ.$$
(34)
Using (33) and (31), one easily verifies that this second constraint is equivalent to
$$\left(R^21\right)ϵ=0.$$
(35)
It is not difficult to show that this equation has no non-zero solutions for $`ϵ`$ unless $`\alpha \pm \beta =0`$. We thus have
$$R=e^{\alpha \mathrm{\Sigma }},\mathrm{\Sigma }=\frac{1}{2}\left(\gamma _5\pm \gamma _{12}\right).$$
(36)
Using the identity $`\mathrm{\Sigma }^3=\mathrm{\Sigma }`$ one can establish that
$$R^21=(2R)(\mathrm{sin}\alpha \mathrm{\Sigma }).$$
(37)
Since $`2R`$ is invertible, it follows that (35) is equivalent to
$$\mathrm{sin}\alpha \mathrm{\Sigma }ϵ=0.$$
(38)
This is trivially satisfied if $`\mathrm{sin}\alpha =0`$. Otherwise it reduces to $`\mathrm{\Sigma }ϵ=0`$, which is equivalent to
$$\gamma _{03}ϵ=\pm ϵ.$$
(39)
If this is combined with (31) we deduce that
$$\gamma _5\gamma _{023}ϵ=ϵ,$$
(40)
which is the constraint associated with a purely magnetic domain wall in the 23-plane. We may take any two of these three constraints as the independent ones; the choice (31) and (40) have an obvious interpretation as the constraints associated with the orthogonal intersection of an electric wall with a magnetic one. This constitutes the special $`\alpha =\pi /2`$ case of the more general configuration of rotated intersecting branes that we have been studying. But we have now derived these constraints for any angle $`\alpha 0,\pi `$. The fraction of supersymmetry preserved by the general rotated brane configuration is therefore the same as the fraction preserved in the special case of orthogonal intersection. Standard arguments can now be used to show that this fraction is 1/4.
We have thus shown that starting from a 1/2 supersymmetric configuration of two parallel coincident domain walls with normal $`𝐧`$, one of them may be rotated relative to the other by an arbitrary angle in a plane containing $`𝐧`$, preserving 1/4 supersymmetry, provided that the charge of the rotated wall is simultaneously rotated by the same angle in the ‘electric-magnetic’ charge space. In practice it may not be possible for the domain walls to intersect at arbitrary angles (preserving supersymmetry). For example, in the $`\text{}_3`$-invariant model discussed in , supersymmetric intersections are necessarily at $`2\pi /3`$ angles. But such additional restrictions are model-dependent. What we learn from the supersymmetry algebra is the model-independent result that the angle separating 1/4 supersymmetric intersecting domain walls must equal the angle between them in the ‘electric/magnetic’ charge space.
Since the constraint (39) is associated with non-zero $`P_3`$ we also learn from the above analysis that we can include this charge, provided it has the appropriate sign, which is determined by the sign in (36), without affecting the constraints imposed by 1/4 supersymmetry, although we then leave the class of configurations for which $`b=0`$. Setting $`P_30`$ might be considered as performing a boost along the 3-direction except for the previously noted fact that $`P_3`$ is not necessarily to be interpreted as momentum. Nevertheless, as a terminological convenience we shall call $`𝐏`$ the ‘3-momentum’ in what follows. Consider the charge configuration obtained by adding the charges of an electric brane in the 13-plane with a brane rotated in the 12-plane, preserving 1/4 supersymmetry, and then adding momentum in the 3 direction:
$`𝐔`$ $`=`$ $`v\mathrm{cos}\alpha (\mathrm{sin}\alpha ,\mathrm{cos}\alpha ,0)+(0,u,0)`$
$`𝐕`$ $`=`$ $`v\mathrm{sin}\alpha (\mathrm{sin}\alpha ,\mathrm{cos}\alpha ,0)`$
$`𝐏`$ $`=`$ $`(0,0,p)`$ (41)
We now have
$`a`$ $`=`$ $`u^2+v^2+2uv\mathrm{cos}^2\alpha +p^2`$
$`b`$ $`=`$ $`puv\mathrm{sin}^2\alpha `$
$`c`$ $`=`$ $`u^2v^2\mathrm{sin}^4\alpha +p^2(u^2+v^2+2uv\mathrm{cos}^2\alpha )`$ (42)
One can show that the eigenvalues of $`\{Q,Q\}`$ are
$$H+p\pm \sqrt{u^2+v^2+2uv\mathrm{cos}2\alpha },Hp\pm (u+v)$$
(43)
For $`u,v,p0`$, we conclude that $`Hp+u+v`$ and that 1/4 supersymmetry is preserved when the bound is saturated. Note that in this case
$$\{Q,Q\}=u\left(1\gamma _{013}\right)+v\left(1\gamma _{013}R^2\right)+p\left(1\gamma _{03}\right),$$
(44)
for the upper sign in (36), confirming that the projections remain unchanged by the inclusion of momentum.
### 2.4 3/4 Supersymmetry
Continuing the above analysis, we now turn to the case in which $`u,v,p`$ are not necessarily all positive because this case includes the possibility of domain wall configurations preserving 3/4 supersymmetry . Consider the case $`\alpha =\pi /2`$ for an electric wall and a magnetic wall intersecting at right angles, so that the eigenvalues (43) are
$$H+p\pm (uv),Hp\pm (u+v).$$
(45)
It follows that $`H`$ is bounded below by each of the eigenvalues
$`\lambda _1`$ $`=`$ $`puv`$
$`\lambda _2`$ $`=`$ $`vup`$
$`\lambda _3`$ $`=`$ $`uvp`$
$`\lambda _4`$ $`=`$ $`u+v+p.`$ (46)
If only one of the charges is non-zero, $`u`$ say, then we obtain the standard BPS bound, $`H|u|`$, which is saturated by the electrically charged BPS domain wall. With two charges, $`u`$ and $`v`$ say, we obtain $`H|u+v|`$ and $`H|uv|`$, and when the stronger of these is saturated we have the intersecting domain wall configuration preserving 1/4 supersymmetry. With all three charges, there are four bounds corresponding to the four eigenvalues and 1/4 supersymmetry is preserved, generically, when the strongest bound is saturated. There are then two subcases to consider according to whether or not $`\lambda _4`$ is the largest eigenvalue. If $`\lambda _4`$ is the largest eigenvalue, as happens, for example, when $`u,v,p`$ are all positive, then we recover the standard 1/4 supersymmetric case considered above, unless two of the three charges $`u,v,p`$ vanish in which case 1/2 supersymmetry is preserved. If $`\lambda _4`$ is not the largest eigenvalue then one of the others is, and we may choose it to be $`\lambda _1`$ because the other possibilities are related to this one by $`SU(2)_R`$ transformations. Given this, $`H`$ is bounded below by $`puv`$ and if there is a state saturating this bound with $`H=puv`$ then the eigenvalues of $`\{Q,Q\}`$ are
$$0,2(pv),2(pu),2(u+v).$$
(47)
It follows that 1/4 supersymmetry is preserved generically but more supersymmetry is preserved for special values of the charges. The possibility of this kind of enhancement of supersymmetry, including the possibility of 3/4 supersymmetry, was recently discussed in and the case under consideration here is very similar. If $`p=v`$ or $`p=u`$ or $`u=v`$, then a charge configuration saturating the BPS bound will preserve 1/2 supersymmetry and if $`p=u=v`$ or $`u=v=\pm p`$ then 3/4 supersymmetry will be preserved. Thus, a charge configuration saturating the bound $`H\lambda _1`$ will preserve 1/4 supersymmetry for generic values of the charges, but 1/2 or 3/4 supersymmetry for certain special values.
We should stress that the above analysis is purely algebraic and it is an open question whether there exists a physical model with domain wall configurations preserving 3/4 supersymmetry. As we now show, this possibility is not realized by the WZ model.
### 2.5 BPS Solutions of the Wess-Zumino Model
The WZ model is known to admit both 1/4 and 1/2 supersymmetric classical solutions, which (at least potentially) correspond to states in the quantum theory. We shall show here that there are no classical solutions preserving 3/4 supersymmetry. We shall begin by considering purely bosonic field configurations and then extend the result to arbitrary classical configurations.
The fields of the WZ model belong to a single chiral superfield, the components of which are a complex physical scalar $`\varphi =A+iB`$, a complex two-component spinor, which is equivalent to a 4-component Majorana spinor $`\lambda `$, and a complex auxiliary field $`F=f+ig`$. We will continue to use a real representation of the four Dirac matrices $`\gamma ^\mu `$. For purely bosonic field configurations we need only consider fermion supersymmetry transformations. Our starting point will therefore be the (off-shell) supersymmetry transformation of $`\lambda `$, which takes the form $`\delta \lambda =Mϵ`$, where $`ϵ`$ is a real constant spinor parameter and $`M`$ is the real $`4\times 4`$ matrix
$$M=\gamma ^\mu (_\mu A+\gamma _5_\mu B)+f+\gamma _5g.$$
(48)
This transformation is valid for the spinor component of any chiral superfield. The Wess-Zumino model is characterised by the fact that the auxiliary field equation is
$$Ff+ig=W^{}(\varphi ),$$
(49)
where $`W^{}(\varphi )`$ is the derivative with respect to $`\varphi `$ of the holomorphic superpotential $`W(\varphi )`$.
A bosonic field configuration of the WZ model will be supersymmetric if there is a spinor field $`ϵ`$ that is both annihilated by $`M(x)`$, for all $`x`$, and covariantly constant with respect to a metric connection on $`\text{𝔼}^{(1,3)}`$. Thus, for there to be $`n`$ preserved supersymmetries it is a necessary condition that $`M(x)`$ has an $`n`$ dimensional kernel for each $`x`$. Our strategy for showing that there are no 3/4 supersymmetric field configurations will be to analyse necessary conditions for the matrix $`M_0M(x_0)`$ at a fixed point $`x_0`$ to have an $`n`$-dimensional kernel.
We begin by noting that a WZ field configuration can preserve 1/4 supersymmetry only if $`detM_0`$ vanishes, which is equivalent to
$$\left[(A)^2+(B)^2f^2g^2\right]^2=4\left[(A)^2(B)^2(AB)^2\right].$$
(50)
This condition is necessary for the preservation of at least 1/4 supersymmetry in any model with a single chiral superfield, and in particular in the WZ model. Configurations preserving more than 1/4 supersymmetry are characterized by additional constraints on the fields. Necessary constraints can be found very easily by making use of the fact that $`M_0`$ can be brought to (real) upper-triangular form by a similarity transformation. We may therefore assume that $`M_0`$ is upper triangular. If, in addition, it has a 2-dimensional kernel then it may be brought to the form
$$\left(\begin{array}{cccc}0& 0& & \\ & 0& & \\ & & & \\ & & & \end{array}\right)$$
(51)
where $``$ indicates an entry that is not zero (or not necessarily zero in the case of the off-diagonal entries). This matrix has the property that
$$2\mathrm{t}\mathrm{r}M_0^33\mathrm{t}\mathrm{r}M_0\mathrm{tr}M_0^2+(\mathrm{tr}M_0)^3=0,$$
(52)
and substituting (48) we learn that
$$f\left[f^2+g^2(A)^2(B)^2\right]=0.$$
(53)
This condition is therefore necessary for a field configuration to preserve 1/2 supersymmetry.
Similarly, any upper-triangular matrix with a 3-dimensional kernel can be brought to the form
$$\left(\begin{array}{cccc}0& 0& 0& \\ & 0& 0& \\ & & 0& \\ & & & \end{array}\right)$$
(54)
This matrix satisfies both (52) and $`\mathrm{tr}M_0^2=(\mathrm{tr}M_0)^2`$, in addition to (50). These conditions, which are therefore necessary for 3/4 supersymmetry, are equivalent to the joint conditions
$`f`$ $`=`$ $`0`$
$`g^2`$ $`=`$ $`(A)^2+(B)^2`$
$`(A)^2(B)^2`$ $`=`$ $`(AB)^2.`$ (55)
We are now in a position to show that there are no 3/4 supersymmetric WZ field configurations (other than the vacuum which has $`4/4`$ supersymmetry). The conditions (2.5) must be satisfied by such a field configuration. We shall analyse these conditions at a fixed point $`x=x_0`$ and consider separately the cases in which $`g=0`$ and $`g0`$ at that point. If $`g=0`$ then the second condition in (2.5) implies that at $`x_0`$ either the 4-vectors $`A`$ and $`B`$ are both null or one is spacelike and the other is timelike. The latter option contradicts the third of eqs (2.5) so both are null. It then follows from (2.5) that $`A`$ and $`B`$ are parallel, so that
$$f=g=0,A=\alpha v,B=\beta v$$
(56)
where $`\alpha `$ and $`\beta `$ are constants and $`v`$ is a null 4-vector. This field configuration is therefore a candidate for 3/4 supersymmetry, but because the conditions leading to it were not sufficient for 3/4 supersymmetry this must be checked. In fact, it is readily shown that the matrix $`M`$ corresponding to the configuration (56) has only a two-dimensional kernel so that at most 1/2 supersymmetry can be preserved.
The remaining candidates for 3/4 supersymmetry in the WZ model arise from field configurations in which $`f`$ vanishes but $`g`$ is non-zero. Then (2.5) implies that at $`x_0`$ either $`A`$ and $`B`$ are both spacelike, or one is spacelike and the other is null. Suppose first that either $`A`$ or $`B`$ is null. In the case in which $`B`$ is null we have
$$f=0A=gs,B=v,$$
(57)
where $`v`$ is a null vector orthogonal to a spacelike vector $`s`$ normalized such that $`s^2=1`$. For this configuration one can check that the matrix $`M`$ generically has a one dimensional kernel, and has a two dimensional kernel when either $`g=0`$ or $`\beta =0`$. The case in which $`A`$ is null is similar, with the same result that at most 1/2 of the supersymmetry is preserved.
If neither $`A`$ nor $`B`$ is null then they are both spacelike and we can arrange for them to take the form
$`B`$ $`=`$ $`\beta (0,1,0,0)`$
$`A`$ $`=`$ $`\alpha (\mathrm{sin}\theta ,\mathrm{cos}\theta ,0,\mathrm{sin}\theta )`$ (58)
with $`g^2=\alpha ^2\mathrm{cos}^2\theta +\beta ^2`$. One then finds that the kernel of $`M(x_0)`$ is 2-dimensional if $`\alpha \beta \mathrm{sin}\theta =0`$ and otherwise 1-dimensional. Configurations of the form (2.5) can therefore preserve at most 1/2 supersymmetry.
We have now shown that there are no non-vacuum bosonic WZ field configurations that preserve 3/4 supersymmetry. We now wish to consider whether this remains true when we consider general configurations that are not necessarily bosonic. This question is perhaps best posed in the context of the quantum theory, which we will not consider here, but it can also be posed classically by taking all fields to be supernumbers with a ‘body’ and a nilpotent ‘soul’. Any general field configuration of this kind preserving 3/4 supersymmetry must have a body preserving at least 3/4 supersymmetry and, as we have just seen, the vacuum configuration is the only candidate. It follows that the only remaining way in which a classical field configuration could be 3/4 supersymmetric is if the 4/4 supersymmetry of the bosonic vacuum configuration is broken to 3/4 by fermions. Preservation of any fraction of supersymmetry in a fermionic background requires the vanishing of the supersymmetry transformations of the bosons. For the WZ model this implies ($`\overline{\lambda }\lambda ^TC`$)
$$\overline{\lambda }ϵ=0,\overline{\lambda }\gamma _5ϵ=0,$$
(59)
and for 3/4 supersymmetry there must be a three-dimensional space of parameters $`ϵ`$ for which this condition holds. At a given point in space we may choose, without loss of generality, a basis in spinor space such that $`Cϵ=(0,,,)^T`$, where an asterisk indicates an entry that may be non-zero. The first equation then implies that $`\lambda ^T=(,0,0,0)`$ and the second that $`\lambda ^T\gamma _5=(,0,0,0)`$. But since $`\gamma _5`$ is both real and satisfies $`\gamma _5^2=1`$ these conditions are not mutually compatible. This concludes our proof that the WZ model has no non-vacuum classical configurations, bosonic or otherwise, that preserve 3/4 supersymmetry
## 3 The geometry of supersymmetry
We now turn to a discussion of the geometry associated with BPS representations of the algebra (2), which we may re-write in terms of a positive semi-definite symmetric bispinor $`Z`$ as $`\{Q,Q\}=Z`$. The positivity of $`\{Q,Q\}`$ implies that $`Z`$ is a vector in a convex cone, with the boundary of the cone corresponding to the BPS condition $`detZ=0`$. We shall first explain some of the geometry associated with convex cones, and how it relates to BPS states. We will then explain how this ties in with the theory of Jordan algebras.
### 3.1 Convex cones
Let us begin with the standard D=4 N=1 supersymmetry algebra, in which case $`Z=\gamma P`$ and the positivity of $`\{Q,Q\}`$ implies that $`P`$ lies either in the forward lightcone of D=4 Minkowski momentum-spacetime or on its boundary, the lightfront. In the latter case, $`P^2=0`$ and any states with this 4-momentum are BPS, preserving 1/2 supersymmetry. The forward lightcone in momentum space and the forward lightcone in position space are both examples of convex cones. An $`n`$-dimensional cone $`𝒞`$ is a subspace of an $`n`$-dimensional vector space $`V`$ with the property that $`\lambda x𝒞`$ for all $`x𝒞`$ and all real positive $`\lambda `$. The cone is convex if the sum of any two vectors in the cone is also in it. The dual cone is then defined as follows. Let $`y`$ be a vector in the dual vector space $`V^{}`$ and let $`yx`$ be a bilinear map from $`V\times V^{}`$ to . The dual cone $`𝒞^{}`$ is the subspace of $`V^{}`$ for which $`yx>0`$ for all $`x𝒞`$.
Given a translation-invariant measure on $`V`$ we can associate with each convex cone in $`V`$ a characteristic function $`\omega `$ defined by
$$\omega ^1(x)=_𝒞^{}e^{yx}d^ny.$$
(60)
As all translation-invariant measures are multiples of any given translation-invariant measure, this formula defines $`\omega `$ up to a scale factor, but this ambiguity will not affect the statements to follow. The cone is foliated by hypersurfaces of constant $`\omega `$, with the limiting hypersurface $`\omega =0`$ being the boundary of the cone. In the case of the forward light cone in D=4 Minkowski spacetime the vector space $`V`$ is $`\text{}^4`$ and $`\omega =𝒩^2`$, where $`𝒩(x)=\eta _{\mu \nu }x^\mu x^\nu `$ is the quadratic form defined by the Minkowski metric $`\eta `$ (we adopt a ‘mostly plus’ metric convention). The hypersurfaces of constant $`\omega `$ are therefore hyperboloids homothetic to $`SO(1,3)/SO(3)`$. Note that this is a symmetric space; this is a general feature of self-dual homogeneous convex cones, of which the forward lightcone in Minkowski space is an example. Homogeneous convex cones that are not self-dual are foliated by homogeneous spaces that are not symmetric spaces.
Because, in this example, $`\omega `$ is determined by a quadratic function $`𝒩`$, the vector space $`V=\text{}^4`$ can be viewed as a metric space, with Minkowski metric $`\eta `$. More generally, $`\omega `$ is not quadratic and hence does not furnish $`V`$ with a metric. Nevertheless, $`\omega `$ does provide a positive definite metric for $`𝒞`$ (obviously, this differs from the Minkowski metric of the ‘quadratic’ case discussed above). Let us first note that, by the definition of a cone, the map $`D:x\lambda x`$ is an automorphism, in that $`Dx𝒞`$ if $`x𝒞`$. It follows immediately that $`\omega (x)`$ is a homogeneous function of degree $`n`$. A corollary of this is that $`\pi (x)x=1`$ where
$$\pi (x)=\frac{1}{n}\frac{\mathrm{log}\omega }{x}.$$
(61)
Thus, $`\pi 𝒞^{}`$, and as $`x`$ ranges over all vectors in $`𝒞`$ so $`\pi `$ ranges over all vectors in $`𝒞^{}`$. One can now introduce a metric $`g`$ on $`𝒞`$ with components<sup>4</sup><sup>4</sup>4For the forward light-cone in Minkowski spacetime with Minkowski metric $`\eta `$, we have $`g_{ij}=(x^2)^2(2x_ix_jx^2\eta _{ij})`$ where $`x^2=\eta _{ij}x^ix^j`$ and $`x_i=\eta _{ij}x^j`$, so that $`\pi _i=(x^2)^1x_i`$.
$$g_{ij}=\frac{1}{n}_i_j\mathrm{log}\omega (x).$$
(62)
One may verify that
$$\pi _j=x^ig_{ij}.$$
(63)
The map from $`𝒞`$ to $`𝒞^{}`$ provided by the metric (62) has a natural interpretation in terms of Hamilton-Jacobi theory: if $`\mathrm{log}\omega `$ is interpreted as a characteristic function in the sense of Hamilton, then $`\pi `$ as defined by (61) is the conjugate momentum.
A feature of the metric $`g`$ is that it is invariant under automorphisms of $`𝒞`$. For example it follows from the homogeneity of $`\omega `$ that the linear map $`D`$ is an isometry of $`g`$. The group of automorphisms will generally be a semi-direct product of $`D`$ with some group $`G`$ that acts on the leaves of the foliation. The cone is homogeneous if $`G`$ acts transitively. A homogeneous cone is foliated by homogeneous hypersurfaces of the form $`G/H`$ for some isotropy group $`H`$. For a self-dual cone this homogeneous space is also a symmetric space. As already mentioned, the forward light cone in $`\text{𝔼}^{(1,3)}`$ is foliated by hyperboloids homothetic to $`SO(1,3)/SO(3)`$, so $`G`$ is the (proper orthochronous) Lorentz group. The metric induced on each leaf of the foliation by the metric $`g_{ij}`$ of the cone is the positive-definite $`SO(1,3)`$-invariant metric on $`SO(1,3)/SO(3)`$.
Let us now turn to the general D=4 N=1 supersymmetry algebra $`\{Q,Q\}=Z`$. The bispinor charge $`Z`$ can be interpreted as a vector in the convex cone of positive-definite real $`4\times 4`$ symmetric matrices. This is a cone in $`\text{}^{10}`$ which, since $`Z`$ includes the 4-momentum, we may consider as a ‘momentum-space’ cone $`𝒞^{}`$. We set aside to the following subsection consideration of the corresponding ‘position space’ cone $`𝒞`$. The characteristic function of $`𝒞^{}`$ is<sup>5</sup><sup>5</sup>5Note that $`\omega ^2`$ is a polynomial. A theorem of Koecher states that $`\omega ^2`$ is a polynomial for all self-dual homogeneous convex cones.
$$\omega (Z)=(detZ)^{\frac{5}{2}}.$$
(64)
The cone is again a self-dual homogeneous one, and is foliated by symmetric spaces that are homothetic to $`SL(4;\text{})/[SO(4)]`$. Of principal interest here is the boundary of $`𝒞^{}`$, defined by $`detZ=0`$, because this is the condition for preservation of supersymmetry. The geometry of this boundary is now rather more complicated than it was before.
The basic observation required to understand this geometry is that the cone is a stratified space with strata $`𝒮_n`$, $`n=0,1,2,3,4`$, where $`𝒮_n`$ is the subspace in which at least $`n`$ of the four eigenvalues vanish, corresponding to at least $`n`$ supersymmetries being preserved, and $`𝒮_{n+1}`$ is the boundary of $`𝒮_n`$. The boundary of the cone is the space $`𝒮_1`$, which is the 9-dimensional space of matrices of rank 3 or less. The boundary of this is the space $`𝒮_2`$ of matrices of rank 2 or less which make up a 7 dimensional space. To see why it is 7 dimensional recall that to specify a matrix of rank 2 it suffices to give the normalised eigenvectors with non-vanishing eigenvalues together with their eigenvalues. The two eigenvectors define a 2-plane in $`\text{}^4`$, corresponding to an element of the 4-dimensional Grassmannian $`SO(4)/(SO(2)\times SO(2))`$. Giving the orientation of the eigenvectors within the 2-plane means specifying one of the $`SO(2)`$ factors. In other words the basis of 2 eigenvectors corresponds to the 5 dimensional Stiefel manifold $`SO(4)/SO(2)`$. Taking into account the two eigenvalues we have a 7-dimensional space, as claimed. The boundary of this stratum is the set $`𝒮_3`$ of matrices of rank 1 or less. These span a 4-dimensional space, since a rank 1 matrix is specified by the direction, up to a sign, of its eigenvector with non-zero eigenvalue together with the eigenvalue. This is a point in $`\text{}P^3\times \text{}^+`$. Finally, the boundary of $`𝒮_3`$ is the stratum $`𝒮_4`$ consisting of the zero matrix, which is the vertex of the cone.
### 3.2 Reverse triangle inequalities
The Minkowski inequality that we used previously to establish the stability of BPS states is a special case of a reverse-triangle inequality valid for all convex cones. Let us define the ‘length’ of a vector in an n-dimensional convex cone with characteristic function $`\omega `$ as
$$L(x)=\omega ^{1/n}(x).$$
(65)
This is a homogeneous function of degree 1. Because the hypersurfaces of constant $`\omega `$ are concave, this ‘length’ satisfies the reverse triangle inequality
$$L(x+x^{})L(x)+L(x^{}).$$
(66)
with equality if and only if $`x`$ and $`x^{}`$ are proportional. In the case of the cone of $`m\times m`$ positive definite hermitian matrices we have $`L(x)=(detx)^{1/m}`$ and the reverse triangle inequality is the Minkowski inequality
$$[det(x+y)]^{\frac{1}{m}}[detx]^{\frac{1}{m}}+[dety]^{\frac{1}{m}},$$
(67)
with equality if the two matrices are proportional. In the special case of diagonal matrices, the cone becomes the positive orthant $`R_+^m`$ in $`\text{𝔼}^m`$. The length of a vector $`x=diag(x_1,\mathrm{},x_m)`$ in $`R_+^m`$ is $`L(x)=(x_1\mathrm{}x_n)^{1/n}`$, and Minkowski’s inequality for positive definite matrices reduces to a form of Holder’s inequality (see e.g. ). The metric $`g`$ on $`R_+^m`$ is the flat metric $`dl^2=(1/n)(d\mathrm{log}x^i)^2`$. The automorphism group is the permutation group $`S_m`$, which is clearly an invariance of the length.
### 3.3 Conformal Invariance
For the standard D=4 N=1 supersymmetry algebra without central charges all BPS states have $`P^2=0`$. This is the momentum space version of the massless wave-equation, which is invariant under the action of the conformal group $`SU(2,2)`$ on compactified Minkowski spacetime. Our aim here is to show how this generalizes when the domain wall charges are included. This will turn out to be a straightforward extension of the standard case, appropriately formulated, so we consider that first.
It is convenient to identify a point in Minkowski spacetime with a matrix $`X=X^\mu \sigma _\mu `$, where $`\sigma _\mu =(1,\sigma _1,\sigma _2,\sigma _3)`$ are the $`2\times 2`$ Hermitian sigma-matrices. The conjugate momentum $`P`$ is then similarly a $`2\times 2`$ Hermitian matrix and $`P^2`$ becomes $`detP`$. (The momentum $`P`$ should not be confused with the dual variable $`\pi `$ introduced in the previous subsection.) Let us now consider the massless particle action
$$I=[\mathrm{tr}PdXedetP]$$
(68)
where $`e`$ (the einbein) is a Lagrange multiplier for the mass-shell constraint $`detP=0`$. The conformal group $`SU(2,2)`$ acts on the compactification of Minkowski space via the fractional linear transformation
$$XX^{}=(AX+B)(CX+D)^1$$
(69)
where the hermiticity of $`X^{}`$ requires that
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)SU(2,2).$$
(70)
This implies that
$$dX^{}(CX+D)=(AX^{}C)dX.$$
(71)
We deduce from this that the $`PdX`$ part of the action $`I`$ is invariant (up to a surface term) if
$$PP^{}=(CX+D)P(AX^{}C)^1.$$
(72)
This transformation implies
$$detPdetP^{}=\mathrm{\Omega }^1detP$$
(73)
where
$$\mathrm{\Omega }=\frac{det(AX^{}C)}{det(CX+D)}.$$
(74)
The action $`I`$ is therefore invariant if we assign to the einbein the transformation $`ee^{}=\mathrm{\Omega }e`$.
We now wish to determine the analogous symmetry group of the more general BPS condition $`detZ=0`$. The matrix $`Z`$ can be viewed as a vector in a 10-dimensional vector space. Let $`X`$ be coordinates of the dual space and consider the particle action
$$I=[\mathrm{tr}ZdXedetZ].$$
(75)
Special cases of actions of this type were considered previously by Cederwall , with a motivation derived from Jordan algebra considerations that we shall explain in the following subsection (see also ). Now consider the fractional linear transformation
$$XX^{}=(AX+B)(CX+D)^1,$$
(76)
which acts on the compactification of the space of symmetric matrices . The matrix $`X^{}`$ will also be real and symmetric provided that
$$\left(\begin{array}{cc}A& B\\ C& D\end{array}\right)Sp(8;\text{}).$$
(77)
That is,
$$A^TDC^TB=1,A^TC=C^TA,B^TD=D^TB.$$
(78)
As before, we deduce (71) and from this that the $`ZdX`$ term is invariant up to a surface term if
$$ZZ^{}=(CX+D)Z(AX^{}C)^1.$$
(79)
This implies
$$detZdetZ^{}=\mathrm{\Omega }^1detZ$$
(80)
where $`\mathrm{\Omega }`$ has form of (74). We may again take $`ee^{}=\mathrm{\Omega }e`$ to achieve an invariance of the action $`I`$. In this case, the invariance group is $`Sp(8;\text{})`$.
Note that this conclusion rests on an interpretation of the 4-dimensional compactified Minkowski spacetime as a subspace of a ten-dimensional vector space of the $`4\times 4`$ real symmetric matrices $`X`$. A field theory realization of Sp(8;R) would require fields defined on this larger space. For example, the analogue of the massless wave equation on Minkowski space is the fourth-order equation
$$det(i/X)\mathrm{\Psi }=0$$
(81)
The symmetry group of this equation is $`Sp(8;\text{})`$. By analogy with the Minkowski case, we expect this to be the maximal symmetry group of this equation.
### 3.4 Jordan algebras
The results of the previous subsections have an interpretation in terms of Jordan algebras. A Jordan algebra $`J`$ of dimension $`n`$ and degree $`\nu `$ is an $`n`$-dimensional real vector space with a commutative, power associative, bilinear product, and a norm $`𝒩`$ that is a homogeneous polynomial of degree $`\nu `$ (see e.g. ). There are four infinite series of simple Jordan algebras, realizable as matrices with the Jordan product being the anticommutator: the degree 2 algebras $`\mathrm{\Sigma }(n)`$ to be discussed below, and the series $`J_k^{\text{}}`$, $`J_k^{\text{}}`$, $`J_k^{\text{}}`$, which are realized by $`k\times k`$ hermitian matrices over , or , with norm given by the determinant, $`𝒩(x)=det(x)`$. In addition, there is one ‘exceptional’ Jordan algebra $`J_3^\text{𝕆}`$ realizable by $`3\times 3`$ hermitian matrices over the octonions.
Associated with any Jordan algebra $`J`$ with product $`xy`$ is a self-dual homogeneous convex cone $`𝒞(J)`$. This is the subspace of $`J`$ consisting of elements $`e^x`$ with $`xJ`$ (where $`e^x`$ is defined by the usual power series with $`x^{n+1}=x^nx`$). The characteristic function is
$$\omega =𝒩^{n/\nu },$$
(82)
so the boundary of the cone corresponds to $`𝒩=0`$. The cone is foliated by copies of the homogeneous space $`Str(J)/Aut(J)`$, where $`Str(J)`$ is the invariance group of $`𝒩`$ (the ‘structure group’ of the algebra) and $`Aut(J)`$ is the automorphism group of the algebra (the subgroup of $`Str(J)`$ that fixes the identity element in $`J`$).
The relation of self-dual homogeneous convex cones to Jordan algebras has similarities to the relation between Lie groups and Lie algebras. Recall that a Lie group is parallelizable but has a non-zero torsion given by the structure constants of its Lie algebra. A self-dual homogeneous convex cone $`𝒞`$, on the other hand, is not parallelizable (in general) but its torsion-free affine connection is determined by the structure constants of a Jordan algebra. Because of homogeneity it suffices to know the connection at the ‘base’ point $`c𝒞`$ defined by<sup>6</sup><sup>6</sup>6There is only one such point, even in those cases for which $`𝒞`$ is flat. It corresponds to the identity element in the algebra. We use the notation $`c`$ to indicate both the identity element of $`J`$ and the base point of the cone $`𝒞(J)`$.
$$g_{ij}|_c=\delta _{ij}.$$
(83)
Let $`f_{ij}^k`$ be the structure constants of $`J`$ in a basis $`e^i=(c,e_a)`$. Then
$$\mathrm{\Gamma }_{ij}{}_{}{}^{k}|_{c}^{}=f_{ij}{}_{}{}^{k}.$$
(84)
Although Jordan algebras are commutative they are nonassociative. Define the associator
$$\{a,b,c\}(ab)ca(bc).$$
(85)
The curvature tensor of the cone at the base point is then given by the relation
$$\{e_i,e_j,e_k\}=R_{ijk}{}_{}{}^{l}|_{c}^{}e_l.$$
(86)
In addition to the automorphism and structure groups, there is a larger ‘Möbius group’ associated with any Jordan algebra $`J`$, acting on elements of $`J`$ by fractional linear transformations. We therefore have the sequence of groups
$$Aut(J)Str(J)Mo(J),$$
(87)
associated with any Jordan algebra $`J`$. These can be interpreted as generalized, rotation, Lorentz and conformal groups, respectively . To motivate this interpretation, we return to the representation of a Minkowski 4-vector as the $`2\times 2`$ Hermitian matrix $`X`$. This is an element in the degree 2 Jordan algebra $`J_2^{\text{}}`$. The dimension is 4 and the norm is the determinant, which is the $`SL(2;\text{})`$ invariant Minkowski norm $`𝒩`$ on $`\text{}^4`$. The group $`SL(2;\text{})`$ acts on $`2\times 2`$ matrices by conjugation so the subgroup leaving invariant the identity matrix is its maximal compact $`SU(2)`$ subgroup. The convex cone associated with this Jordan algebra is the forward light-cone of D=4 Minkowski spacetime. As we saw in the previous subsection, the group of fractional linear transformations of $`X`$ is $`SU(2,2)`$, so the sequence (87) is, in this case,
$$SU(2)SL(2;\text{})SU(2,2).$$
(88)
These are the standard rotation, Lorentz and conformal groups.
The inclusion of domain wall charges means that we should replace $`J_2^{\text{}}`$ by $`J_4^{\text{}}`$, the algebra of $`4\times 4`$ symmetric real matrices. One can see that $`J_2^{\text{}}`$ is a subalgebra of $`J_4^{\text{}}`$ from the fact that $`J_2^{\text{}}\mathrm{\Sigma }(4)`$, where $`\mathrm{\Sigma }(n)`$ is the n-dimensional Jordan algebra with basis $`(1,\sigma _1,\mathrm{}\sigma _{n1})`$ and Jordan product $`\sigma _a\sigma _b=2\delta _{ab}`$; this has a realization in which $`\sigma _a`$ are sigma-matrices of an $`n`$-dimensional Minkowski spacetime, with the Jordan product being the anticommutator; it follows that the standard supersymmetry algebra in $`D`$ dimensions is naturally associated with $`\mathrm{\Sigma }(D)`$. For $`D=4`$ one can choose the $`\sigma _a`$ to be the three $`2\times 2`$ hermitian Pauli matrices, hence the isomorphism $`J_2^{\text{}}\mathrm{\Sigma }(4)`$. All simple Jordan algebras of degree 2 are isomorphic to $`\mathrm{\Sigma }(n)`$ for some $`n`$. Having replaced $`J_2^{\text{}}`$ by $`J=J_4^{\text{}}`$ we find that the sequence (88) is generalized to
$$SU(2)\times SU(2)SL(4;\text{})Sp(8;\text{})$$
(89)
We now turn to the Jordan algebraic interpretation of the boundary of the convex cone $`𝒞(J)`$. This consists of elements $`\lambda PJ`$ where $`\lambda `$ is a positive real number and $`P`$ is an idempotent of $`J`$ with less than maximal rank, i.e. its trace, defined by $`\mathrm{tr}X=\mathrm{log}𝒩(e^X)`$, is less than $`\nu `$. An idempotent is a non-zero element $`PJ`$ satisfying $`PP=P`$, and two idempotents $`P`$ and $`P^{}`$ are said to be orthogonal if $`PP^{}=0`$. The idempotents with unit trace are called the primitive idempotents, and the number of mutually orthogonal primitive idempotents equals the degree $`\nu `$ of the algebra. For a Jordan algebra of degree 2 all idempotents of less than maximal rank have unit trace and are therefore primitive. This is true of $`J_2^{\text{}}`$, in particular, corresponding to the fact that the only supersymmetric states other than the vacuum permitted by the standard D=4 N=1 supersymmetry algebra are 1/2 supersymmetric states associated with massless particles (for which the 4-momentum lies on the positive light-front). Note that although at most two primitive idempotents of a degree 2 Jordan algebra can be orthogonal in the above sense, the space of primitive idempotents of $`\mathrm{\Sigma }(D)`$ is $`(D1)`$-dimensional. The boundary of the associated convex cone is therefore $`(D1)`$-dimensional. For $`\mathrm{\Sigma }(4)J_2^{\text{}}`$, in particular, this boundary is the three-dimensional forward light-front of the origin of 4-dimensional Minkowski momentum space.
For a Jordan algebra $`J`$ of degree $`\nu >2`$, there are idempotents of less than maximal rank that are not primitive. For an algebra of degree 3, these non-primitive idempotents generate faces of the boundary of $`e^J`$ which themselves have a boundary generated by the primitive idempotents. An example is the (non-simple) Jordan algebra $`J=\text{}\text{}\text{}`$ for which $`e^J`$ is the positive octant in $`\text{𝔼}^3`$; its boundary consists of three faces that meet on the three axes generated by the three primitive idempotents (in this case there are only three primitive idempotents, which are therefore orthogonal; details can be found in ). More generally, for Jordan algebras of higher degree, the boundary of the associated convex cone is a stratified set of faces. In particular, $`J_4^{\text{}}`$ has degree 4 so the faces of the boundary of its associated convex cone are generated by idempotents of trace 1,2 and 3, corresponding to 3/4,1/2 and 1/4 supersymmetry respectively. The primitive idempotents, of unit trace, correspond to 3/4 supersymmetry.
### 3.5 Entropy of BPS fusion
In a quantum field theory realization of the D=4 N=1 supersymmetry algebra the central charges $`Z`$ are labels of quantum states. We have now seen that the set of these charges naturally carries the structure of a Jordan algebra. This algebra may itself be regarded as a finite-dimensional state space (not to be confused with infinite-dimensional space of states of the field theory that carry these charges). This interpretation is of course how Jordan algebras originally arose (see for a review). The exceptional Jordan algebra provides a state space more general than conventional quantum mechanics but for all other Jordan algebras the formalism is equivalent to one in which a state is represented by a density matrix. The general state is therefore a mixed state. The pure states correspond to the primitive idempotents; these lie on the boundary of the convex cone $`𝒞(J)`$ but do not in general exhaust it. Rather, the boundary is stratified by sets of states of successively less purity, corresponding in our application to states with successively less supersymmetry. Thus, the pure states in this sense are the charge configurations that preserve 3/4 supersymmetry, the remaining supersymmetric configurations corresponding to states on the boundary of the cone that are not pure.
We previously showed that a BPS state is stable against decay into any other pair of states; in particular it cannot decay into two BPS states. Consider now the reverse process, i.e. fusion of two BPS states to form a third via the inverse of the reaction (22), i.e.
$$(BPS)_1+(BPS)_2(BPS)_3,$$
(90)
If the first two states preserve 3/4 supersymmetry then the third one will generally preserve less supersymmetry. This is like passing from a pure to a mixed state. There is also a formal resemblance here to classical thermodynamics. The Jordan algebra $`J`$, now viewed as vector space $`V`$ containing the convex cone $`𝒞(J)`$, is spanned by the extensive quantities while the dual vector space $`V^{}`$ is spanned by the intensive variables. The function
$$S(x)=\mathrm{log}\omega (x)$$
(91)
of the extensive variables may be interpreted as entropy. Because it is convex
$$S(\mu x+(1\mu )x^{})\mu S(x)+(1\mu )S(x^{}),$$
(92)
with equality when $`x`$ is proportional to $`x^{}`$, the entropy can not decrease as a result of a fusion process such as (90). Conversely, the (marginal) stability of a single BPS state against decay into two other BPS states can now be understood as being forbidden by a version of the second law of thermodynamics.
## 4 BPS states for adS
The N=1 D=4 adS anticommutator (4) may be written as
$$\{Q_\alpha ,Q_\beta \}=\frac{1}{2}M_{AB}\left(𝒞\mathrm{\Gamma }^{AB}\right)_{\alpha \beta },$$
(93)
where
$$\mathrm{\Gamma }^A=(\gamma ^\mu ,\gamma _5)$$
(94)
and $`M_{AB}=M_{BA}`$ are the generators of the adS group $`SO(3,2)`$ (and so are no longer central). The matrix $`𝒞`$ is the $`SO(3,2)`$ charge conjugation matrix; we can choose a representation in which
$$𝒞=\gamma _0\gamma _5$$
(95)
and this choice will be implicit in what follows. Note that
$$\{\mathrm{\Gamma }^A,\mathrm{\Gamma }^B\}=2\eta ^{AB},$$
(96)
where $`\eta `$ is a flat metric on $`\text{𝔼}^{(2,4)}`$, such that $`\eta =\mathrm{diag}(1,1,1,1,1)`$ in cartesian coordinates. Although (4) is preserved by $`GL(4;\text{})`$, the automorphism group of the adS supergroup $`OSp(1|4;\text{})`$ is $`Sp(4;\text{})SL(4;\text{})GL(4;\text{})`$.
The anticommutator (4) can also be written in the form (2), with
$$M_{04}=H,M_{i4}=P_i,M_{0i}=U_i,J_i\frac{1}{2}ϵ_{ijk}M^{jk}=V_i$$
(97)
where $`H`$ is the hamiltonian, $`𝐏`$ the 3-momentum, $`𝐉`$ the angular momentum while the 3-vector $`𝐔`$ generates boosts. The analysis of supersymmetric charge configurations is then exactly the same as in the super-Poincaré case considered earlier, and in particular requiring $`\frac{1}{4},\frac{1}{2}`$ or $`\frac{3}{4}`$ supersymmetry gives exactly the same conditions on the charges $`H,𝐔,𝐕,𝐏`$ as were found earlier.
The condition for preservation of supersymmetry can be expressed in terms of the $`SO(3,2)`$ Casimirs. We will first show how the values of these Casimirs are constrained by the physical state condition, and then turn to the supersymmetric states.
### 4.1 Physical States in adS
Physical states lie either in the convex cone for which $`Z=\frac{1}{2}M_{AB}C\mathrm{\Gamma }_{\alpha \beta }^{AB}`$ is positive, or on its boundary, for which $`\mathrm{det}Z=0`$. This cone is a subspace of the 10-dimensional vector space spanned by $`5\times 5`$ skew-symmetric matrices $`M`$ with entries $`M_{AB}`$. The matrix commutator turns this space into the Lie algebra $`so(3,2)`$. This algebra has rank 2, with quadratic Casimir<sup>7</sup><sup>7</sup>7The quadratic Casimir provides a metric of signature $`(4,6)`$ on the 10-dimensional vector space, but this metric (which is inherited from the metric $`\eta `$ on $`\text{𝔼}^{(3,2)}`$) does not play a crucial role in the following analysis.
$$c_2=\frac{1}{2}M_{AB}M^{AB},$$
(98)
and quartic Casimir
$$c_4=M_B^AM_C^BM_D^CM_A^D.$$
(99)
Since $`\mathrm{det}Z`$ is both a quartic polynomial of the charges and $`SO(3,2)`$ invariant it must be a linear combination of $`c_4`$ and $`c_2^2`$. In fact
$$\mathrm{det}Z=c_4c_2^2,$$
(100)
and hence
$$c_4c_2^2$$
(101)
for physical states.
There is a further constraint on the Casimirs required by physical states. To see this, we begin by noting that the vacuum is the only physical state for which the energy $`M_{04}`$ vanishes. This follows from the fact that $`\{Q,Q\}`$ is positive semi-definite, with a trace equal to $`4M_{04}`$. We next prove that $`M_{04}`$ must vanish if the kernel of $`M`$ contains a timelike 5-vector. Suppose that such a 5-vector exists. By an $`SO(3,2)`$ transformation, we can arrange for it to have only one non-vanishing component, in the $`4`$-direction. It then follows that the only non-vanishing components of $`M`$ are $`M_{\mu \nu }`$. In particular, the energy $`M_{04}`$ vanishes. Thus, for any non-vacuum physical state the kernel of $`M`$ contains no timelike vectors. Note that the kernel of $`M`$ has dimension $`1`$, $`3`$ or $`5`$, according to whether $`M`$ has rank $`4`$, $`2`$ or $`0`$, respectively. The vacuum is the only physical state for which $`M`$ has rank $`0`$.
Now consider the Pauli-Lubanski 5-vector
$$s^A=\frac{1}{8}ϵ^{ABCDE}M_{BC}M_{DE}.$$
(102)
This satisfies the identity
$$M_{AB}s^B0,$$
(103)
which shows that, unless it vanishes, $`s`$ is in the kernel of $`M`$. A timelike $`s`$ would therefore be in the kernel of $`M`$ but, as we have just seen, the kernel of $`M`$ cannot contain timelike vectors unless $`M`$ vanishes, but in that case $`s`$ also vanishes. Thus, $`s`$ cannot be timelike. Now,
$$s^2\eta ^{AB}s_As_B=\frac{1}{4}(2c_2^2c_4),$$
(104)
so $`s`$ will be non-timelike if and only if
$$c_42c_2^2.$$
(105)
This bound implies (for physical states) that $`c_4=0`$ when $`c_2=0`$ .
### 4.2 Supersymmetric States
Our main interest is in BPS states, i.e. the subset of physical states that are supersymmetric. These must saturate the bound (101), so BPS states are those for which
$$c_4=c_2^2.$$
(106)
Using this in (104) we see that
$$s^2=\frac{1}{4}c_2^2$$
(107)
for supersymmetric states. We will organise our discussion of the supersymmetric states according to whether $`s`$ is zero, spacelike or non-vanishing null.
If $`s`$ vanishes then $`M`$ has either a 3-dimensional or a 5-dimensional kernel. $`M`$ will have a 5-dimensional kernel only if it vanishes. If the kernel is 3-dimensional then, as we have seen, it cannot contain timelike vectors. It may contain null vectors but any such null vector must be orthogonal to all other vectors in the kernel, spacelike or null, because we could otherwise find a timelike linear combination. Since the maximum number of mutually orthogonal null 5-vectors is $`2`$, a 3-dimensional kernel must contain at least one spacelike vector. There are three possible choices for the other two linearly independent 5-vectors: (i) both spacelike, (ii) one spacelike and one null, or (iii) both null. In all cases $`M`$ can be brought to a form in which $`M_{04}=E0`$ is its only independent entry. In case (i) $`M_{04}`$ and $`M_{40}`$ are the only entries, and the only supersymmetric state with this property is the vacuum, with $`E=0`$. In case (ii) $`M`$ can be brought to a form for which the only non-zero upper-triangular entries are $`M_{04}=M_{02}=E`$. It then follows from the discussion of section 2.4, on which we will elaborate below, that all such states are 1/2 supersymmetric. In case (iii) $`M`$ can brought to a form for which the only non-zero upper-triangular entries are $`M_{04}=M_{02}=M_{23}=M_{34}`$; all such states are 3/4 supersymmetric.
Consider now spacelike $`s`$. In this case we may choose the only non-vanishing component of $`s`$ to be its $`1`$-component. Since $`s`$ now spans the kernel of $`M`$, this $`5\times 5`$ matrix $`M`$ then reduces to a $`4\times 4`$ matrix $`F`$ acting on the 4-dimensional $`(0234)`$ subspace orthogonal to $`s`$, on which $`\eta `$ restricts to a metric $`\stackrel{~}{\eta }`$ of signature $`(2,2)`$. The matrix $`F`$ is equivalent to a second-rank antisymmetric tensor in $`\text{𝔼}^{(2,2)}`$ that can be written uniquely as $`F=F^++F^{}`$ where $`F^+`$ is real and self-dual while $`F^{}`$ is real and anti-self-dual matrix. Now
$$c_4c_2^2=\left[\mathrm{tr}(\stackrel{~}{\eta }F^+)^2\right]\left[\mathrm{tr}(\stackrel{~}{\eta }F^{})^2\right].$$
(108)
We can write $`F`$ as
$$F=\left(\begin{array}{cccc}0& u& b& E\\ u& 0& v& c\\ b& v& 0& p\\ E& c& p& 0\end{array}\right)$$
(109)
provided that
$$vE+bc+up0,$$
(110)
since $`s`$ would otherwise vanish. Now
$$\mathrm{tr}(\stackrel{~}{\eta }F^\pm )^2=(Ev)^2(u\pm p)^2(b\pm c)^2.$$
(111)
Configurations with self-dual or anti-self-dual $`F`$, for which $`E=v`$, $`u=\pm p`$ and $`b=\pm c`$, are 1/2 supersymmetric. However, any configuration for which
$$(Ev)^2=(u\pm p)^2+(b\pm c)^2$$
(112)
is also supersymmetric. In fact
$$\{Q,Q\}=\left[(Ev)(b\pm c)\gamma ^{012}+(u\pm p)\gamma ^{013}\right]+\left(vc\gamma ^{02}+p\gamma ^{03}\right)\left(1\pm \gamma ^1\right).$$
(113)
Given (112), the term in square brackets is proportional to a 1/2 supersymmetry projector that commutes with the 1/2 supersymmetry projector $`(1/2)\left(1\pm \gamma ^1\right)`$ which leads generically to 1/4 supersymmetry.
The final case to consider is $`s`$ null but non-zero. By means of an $`SO(3,2)`$ transformation we may choose
$$s(1,0,1,0,0)$$
(114)
This choice is preserved by an $`SO(1,2)`$ ‘stability’ subgroup, and by a transformation in the $`SO(2)`$ subgroup of this group we can bring $`M`$ to the standard form
$$M=E\left(\begin{array}{ccccc}0& 0& a& 0& 1\\ 0& 0& a& 0& 1\\ a& a& 0& t& q\\ 0& 0& t& 0& r\\ 1& 1& q& r& 0\end{array}\right)$$
(115)
One then finds that
$$c_2=E^2(t^2q^2r^2),$$
(116)
so that supersymmetric states are those with
$$t=\pm \sqrt{q^2+r^2}.$$
(117)
Actually, in arriving at the above form of $`M`$ we have used only that the null 5-vector $`(1,1,0,0,0)`$ is in the kernel of $`M`$. To ensure that this 5-vector is proportional to $`s`$ (with non-zero constant of proportionality) we require that
$$t+ra0.$$
(118)
This condition also ensures (as it must) that $`M`$ has rank 4. When combined with (117) it implies that
$$t0.$$
(119)
For $`M`$ of the form (115) we have
$$\{Q,Q\}=E\left\{\left(1a\gamma _3\right)\left(1\gamma _{01}\right)t\gamma _1\left[1(q/t)\gamma _{012}(r/t)\gamma _{013}\right]\right\}$$
(120)
A spinor $`ϵ`$ is in the kernel of $`\{Q,Q\}`$ if
$$\left[(q/t)\gamma _{012}+(r/t)\gamma _{013}\right]ϵ=ϵ,$$
(121)
and
$$\gamma _{01}ϵ=ϵ,$$
(122)
and these two constraints imply 1/4 supersymmetry. Note that when $`a=\pm 1`$ and $`q=0`$ and hence $`t=\pm r`$, the latter constraint can be replaced by $`\gamma _3ϵ=\pm ϵ`$, which again yields 1/4 supersymmetry.
### 4.3 Examples
Many of the possibilities for BPS configurations just noted are illustrated by the class of examples considered in section 2.4. This means, in the language of this section, that the non-zero upper-triangular components of $`M_{AB}`$ are taken to be $`M_{04}=E`$, $`M_{34}=p`$, $`M_{02}=u`$ and $`M_{23}=v`$. The Pauli-Lubanski 5-vector is then
$$s=(0,Ev+up,0,0,0),$$
(123)
so $`s`$ is spacelike unless it vanishes. The Casimirs for this class are given by
$`c_2`$ $`=`$ $`E^2+v^2p^2u^2`$ (124)
$`c_4`$ $`=`$ $`2[E^4+u^4+v^4+p^42(v^2+E^2)(u^2+p^2)4Euvp]`$ (125)
The BPS condition $`c_4=c_2^2`$ becomes
$$(Euvp)(Euv+p)(Eu+vp)(E+uvp)=0,$$
(126)
in agreement with (45).
Let us first consider vanishing $`s`$. We have seen above that $`M`$ can be brought to a standard form in which all charges are determined in terms of $`M_{04}=E`$. The non-vacuum BPS states occured for cases (ii) and (iii) discussed above. An example of case (ii) within the class of configurations now under discussion is found by setting $`v=p=0`$ and $`E=|u|`$<sup>8</sup><sup>8</sup>8The charge $`u`$ can be interpreted as a membrane charge. To see this note that there is a static planar solution of the equations of motion of a test membrane in $`adS_4`$ at a fixed radial distance, in horospherical coordinates, from the Killing horizon . This solution must preserve 1/2 supersymmetry of the $`adS_4`$ supersymmetry because $`adS_4`$ can itself be interpreted as a membrane, at the horizon, to which the test membrane is parallel. Because this test membrane remains at a fixed distance from the horizon, the worldline of a point on it is uniformly accelerated, and therefore naturally associated with a non-zero boost $`u`$.. Finally, an example of case (iii), with 3/4 supersymmetry, is obtained by setting $`u=v=p=E<0`$, although there is no known field theoretic realization of this case.
We next to turn to examples with $`s`$ spacelike. Let us first consider $`u=p=0`$ and set $`v=J`$, where $`J`$ is the spin about the 1-axis. We then have
$$c_2=E^2+J^2,c_4=2E^4+2J^4,$$
(127)
which is equivalent to
$`E`$ $`=`$ $`\sqrt{{\displaystyle \frac{c_2}{2}}+{\displaystyle \frac{1}{2}}\sqrt{c_4c_2^2}},`$
$`J`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{c_2}{2}}{\displaystyle \frac{1}{2}}\sqrt{c_4c_2^2}}.`$ (128)
The physical states satisfy $`E|J|`$ and states that saturate this bound preserve 1/2 supersymmetry. For these configurations the matrix $`F`$ of (109) is either self-dual or anti-self-dual. An example of states with $`s`$ spacelike and $`F`$ neither self-dual or anti-self-dual can be obtained by taking $`u,v,p`$ to be positive and solving (112) via $`E=u+v+p`$. We then have
$$\{Q,Q\}=u(1+\gamma ^{013})+p(1+\gamma ^{03})+v(1+\gamma ^1)$$
(129)
and 1/4 of the supersymmetry is preserved.
## 5 Comments
We have seen that a full analysis of the D=4 N=1 supersymmetry algebra not only confirms the existence of 1/2 and 1/4 supersymmetric states, realizable within the WZ model, and determines some of their properties, but it also permits states with 3/4 supersymmetry which, as we have shown, cannot be realized by solutions of the WZ model. However, it has been argued that such ‘exotic’ fractions might play a role in other contexts, and with this in mind we have provided a detailed analysis of the BPS states of D=4 N=1 supersymmetry. We have also seen that these states can be understood in terms of the geometry associated with the convex cone of the Jordan algebra $`J_4^{\text{}}`$, and that this leads to a natural generalization of the rotation, Lorentz and conformal groups.
In general, the $`U(1)_R`$ symmetry will be broken to at most a discrete subgroup. For theories with domain walls (e.g. the WZ model), the R-symmetry will be explicitly broken by the scalar potential. In theories with only massless particles, and no domain walls, the $`U(1)_R`$ symmetry will be generically broken to a discrete subgroup by chiral anomalies. For theories in which the domain wall charges are quantized, the $`U(1)_R`$ symmetry will be broken to the discrete subgroup preserving the quantization condition. An example of this is given by M-theory compactified on a 7-manifold of $`G_2`$ holonomy, yielding a D=4 N=1 theory in which the domain walls are M2-branes and wrapped M5-branes, with the M2-brane and M5-brane charges quantized. Given that only a discrete subgroup of $`U(1)_R`$ survives the same is true of the larger group $`SU(2)_R`$.
We noted that, in the classical theory, the automorphism group of the full supertranslation algebra is $`GL(4,\text{})`$, but it seems that any realization of this on fields, and any realisation of the generalized conformal group $`Sp(8,\text{})`$, requires an enlargement of 3-space to include coordinates conjugate to the ‘domain-wall’ charges $`𝐔`$ and $`𝐕`$. Of course, the domain wall interpretation is probably no longer appropriate in this case. Other interpretations are certainly possible in the context of particle mechanics . In such one-dimensional field theories it is possible to realize the $`SU(2)_R`$ symmetry between the three 3-vector ‘charges’ $`𝐏,𝐔,𝐕`$ as an internal symmetry. For such models that arise from the toroidal compactification of some D=4 theory with quantized $`𝐔`$ and $`𝐕`$, the 3-momentum will also be quantized and the classical $`GL(4;\text{})`$ symmetry will be broken to the discrete $`GL(4;\text{})`$ subgroup preserving the 9-dimensional charge lattice.
Many of the observations made here for $`N=1`$ $`D=4`$ can of course be generalized to $`N>1`$ or to $`D>4`$. For example the general $`N`$ extended $`D=4`$ supersymmetry algebra has automorphism group $`GL(4N;\text{})`$ and $`det\{Q,Q\}`$ is preserved by the subgroup $`SL(4N,\text{})`$. This leads to the sequence
$$SO(4N)SL(4N;\text{})Sp(8N;\text{}).$$
(130)
for the Jordan algebra $`J_{4N}^{\text{}}`$ of $`4N\times 4N`$ symmetric matrices over the reals. The generalised conformal symmetry of the BPS condition is then $`Sp(8N;\text{})`$, as deduced from a different analysis in .
A $`D>4`$ case of particular interest is the D=11 ‘M-theory algebra’ $`\{Q,Q\}=Z`$ where $`Q`$ is now a 32 component real spinor of the D=11 Lorentz group and $`Z`$ is a $`32\times 32`$ real symmetric matrix containing the Hamiltonian and 527 central charges carried by M-branes . This supersymmetry algebra has automorphism group $`GL(32;\text{})`$, as noted independently in , and $`Z`$ takes values in the convex cone associated with the Jordan algebra $`J_{32}^{\text{}}`$. The sequence (87) of groups associated with this algebra is
$$SO(32)SL(32;\text{})Sp(64;\text{}),$$
(131)
so that $`Sp(64;\text{})`$ is the M-theoretic generalisation of the D=11 conformal group. As in the D=4 case, the realization of any of these larger ‘spacetime’ symmetry groups, or discrete subgroups such as $`GL(32;\text{})`$, would seem to require consideration of an enlarged space of 527 coordinates, as considered for other reasons in .
Finally, we have found many possibilities for new BPS states in anti de Sitter space. It seems likely that some of these, in particular those with 1/4 supersymmetry, will have a realization in the context of N=1 D=4 supersymmetric field theories in an adS spacetime.
## Acknowledgments
We would like to thank C. Gui for bringing ref. to our attention. We also thank M. Günaydin and J. Lukierski for helpful correspondence. JPG thanks the EPSRC for partial support. The work of CMH was supported in part by the National Science Foundation under Grant No. PHY94-07194. All authors are supported in part by PPARC through their SPG $`\mathrm{\#}`$613. |
warning/0001/hep-th0001110.html | ar5iv | text | # Effective Low Energy Theories and QCD Dirac Spectra
## 1 Introduction
There are two essentially different approaches to many-body problems. The first approach is to study them by means of Monte-Carlo simulations of the microscopic theory. The second approach is to isolate the relevant degrees of freedom and to describe them by an effective theory. Both approaches have their merits and generally their complementarity leads to a deeper understanding of the underlying phenomena. In this lecture we will focus on Quantum Chromo-Dynamics (QCD) which is the theory of the strong interactions. A great deal of effort has been devoted to Monte Carlo simulations of lattice QCD . They provide a firm footing for our understanding of nonperturbative phenomena such as confinement and chiral symmetry breaking. For QCD at low energy an alternative approach is possible. Because of spontaneous breaking of chiral symmetry, QCD at low energy reduces to a theory of weakly interacting Goldstone bosons. Although this theory cannot be derived from QCD by means of an ab-initio calculation, its Lagrangian is determined uniquely by chiral symmetry and Lorentz invariance . The validity of this low-energy theory is based on the presence of a mass-gap which is a highly nontrivial and nonperturbative feature of QCD.
Chiral symmetry is spontaneously broken at low temperatures. Mainly through lattice simulations, it is now widely accepted that a chiral restoration transition takes place at a temperature of about 140 $`MeV`$. The order parameter of this transition is the density per unit volume of Dirac eigenvalues near zero. This is the reason why we are interested in the properties of the smallest eigenvalues of the Dirac operator . The situation at finite baryon density and zero temperature is much less well understood. Monte-Carlo simulations are only possible if the fermion determinant is ignored. It has been shown that the so called quenched approximation fails spectacularly at nonzero chemical potential . On general grounds it is expected that a chiral phase transition occurs at a value of the chemical potential where it becomes advantageous to create the lightest particles with nonzero baryon number. In QCD these are the nucleons, but in the quenched approximation, they are Goldstone modes made from quarks and conjugate anti-quarks . At nonzero chemical potential, the relation between the chiral order parameter and the QCD Dirac spectrum, which is now scattered in the complex plane, is much less transparent. The failure of quenching can be understood as the absence of spectral ergodicity; at finite density the ensemble averaged Dirac spectrum is completely different form the spectral averaged Dirac spectrum. We will investigate the Dirac spectrum by means of an effective theory. Our approach is similar to the one used for QCD with two colors and for QCD with adjoint quarks . By investigating the properties of the smallest Dirac eigenvalues we hope to obtain a better understanding of this problem and other related issues.
In this lecture we wish to discuss to what extent spectra of the QCD Dirac operator can be derived from an effective low energy partition function. In the first half of this talk we formulate an effective theory for the Dirac spectrum based on a spectrum generating function that in addition to the usual quarks contains bosonic ghost quarks to properly normalize the spectral density. This trick has been widely used in the super-symmetric formulation of Random Matrix Theory . We identify a domain where the results for the resolvent are given by a chiral Random Matrix Theory with the symmetries of the QCD partition function. We find that our results are in agreement with recent lattice simulations. In the second half of this lecture we discuss quenched QCD Dirac spectra at nonzero chemical potential in terms of an effective theory. This lecture is based on several recent articles and additional background material can be found in several recent reviews .
## 2 The QCD partition function
The QCD partition function for $`N_f`$ quarks with mass $`m`$, temperature $`T`$ and quark chemical potential $`\mu `$ is given by
$`Z(m,\mu ,T)=\mathrm{Tr}e^{\frac{H_{\mathrm{QCD}}\mu N}{T}},`$ (1)
where $`H_{\mathrm{QCD}}`$ is the Hamiltonian of QCD and $`N`$ is the quark number operator. The trace is over all states of the theory. This partition function can be rewritten as a Euclidean Feynman path integral
$`Z(m,\mu ,T)=\stackrel{N_f}{det}(D+m)_{\mathrm{YM}},`$ (2)
where the average of the fermion determinant is over the Euclidean Yang-Mills action and the Dirac operator is denoted by $`D`$. In this lecture we focus on the basic symmetry properties of $`D`$ and its explicit representation is not necessary. For simplicity we take all $`N_f`$ quark masses equal to $`m`$.
### 2.1 The eigenvalues of the Dirac operator
The main subject of this talk are the properties of the smallest eigenvalues of the QCD Dirac operator. In a chiral representation of the gamma-matrices, the Euclidean Dirac operator has the block structure
$`D=\left(\begin{array}{cc}0& id+\mu \\ id^{}+\mu & 0\end{array}\right),`$ (5)
where $`d`$ is the covariant derivative of the color group. For $`\mu 0`$ the Dirac operator does not have any hermiticity properties and the eigenvalues will be scattered in the complex plane. We will return to this case in the second half of this talk. For $`\mu =0`$ the Dirac operator is anti-Hermitian, $`D^{}=D`$ and has purely imaginary eigenvalues. The nonzero eigenvalues occur in pairs $`\pm i\lambda _k`$, whereas zero eigenvalues related to the topology of the gauge fields remain unpaired. For simplicity we will restrict the discussion in this lecture to gauge fields with a trivial topology where all eigenvalues occur as pairs. For example, our gauge fields could be composed of fields of an equal number of instantons and anti-instantons . In terms of its eigenvalues, the Euclidean partition function can be written as
$`Z(m,\mu ,T)={\displaystyle \underset{k}{}}(i\lambda _k+m)^{N_f}_{\mathrm{YM}}.`$ (6)
The order parameter of the chiral phase transition is the chiral condensate
$`\mathrm{\Sigma }`$ $`=`$ $`\underset{m0}{lim}\underset{V\mathrm{}}{lim}{\displaystyle \frac{1}{VN_f}}_m\mathrm{log}Z`$ (7)
$`=`$ $`\underset{m0}{lim}\underset{V\mathrm{}}{lim}{\displaystyle \frac{1}{V}}{\displaystyle \underset{k}{}}{\displaystyle \frac{1}{m+i\lambda _k}}_{\mathrm{QCD}}=\underset{m0}{lim}\underset{V\mathrm{}}{lim}{\displaystyle \frac{1}{V}}{\displaystyle \underset{\lambda _k>0}{}}{\displaystyle \frac{2m}{m^2+\lambda _k^2}}_{\mathrm{QCD}},`$
where the average with label QCD includes both the fermion determinant and the Yang-Mills action. At finite volume, the chiral condensate is zero for $`m=0`$. Only if the thermodynamic limit is taken before the chiral limit ($`m0`$) can the chiral condensate become nonzero. If this is the case, chiral symmetry is broken spontaneously. The situation is analogous to the magnetization in a Heisenberg model. At finite volume and zero magnetic field the magnetization is zero. Spontaneous magnetization arises if the thermodynamic limit is taken before putting the external field to zero. According to Goldstone’s theorem spontaneous breaking of a continuous symmetry implies the existence of massless Goldstone bosons. Because of confinement QCD has a mass gap and, at low energy, the partition function is dominated by the Goldstone modes. The effective theory describing their interactions then follows from chiral symmetry and Lorentz invariance.
One of the questions we would like to address in this lecture is whether there is a relation between the existence of a well-defined low-energy theory and the spectrum of the Dirac operator. A second, seemingly unrelated question is the connection of this low-energy with Random Matrix Theory.
## 3 Symmetries of the Spectrum Generating Function at $`\mu =0`$
We will study the Dirac spectrum by means of the resolvent defined by
$`\mathrm{\Sigma }(z)={\displaystyle \frac{1}{V}}\mathrm{Tr}{\displaystyle \frac{1}{D+z}}_{\mathrm{QCD}}.`$ (8)
The mass $`z`$ in the resolvent is not related to the mass $`m`$ in the fermion determinant. Such mass is known in the literature as a valence quark mass. However, we will use the more appropriate name of spectral mass.
In terms of the spectral density
$`\rho (\lambda )={\displaystyle \underset{k}{}}\delta (\lambda \lambda _k)_{\mathrm{QCD}},`$ (9)
the resolvent can be written as
$`\mathrm{\Sigma }(z)={\displaystyle \frac{1}{V}}{\displaystyle \rho (\lambda )𝑑\lambda \frac{1}{i\lambda +z}}.`$ (10)
Therefore, $`\mathrm{\Sigma }(z)`$ is an analytic function with a cut on the imaginary axis. This identity can be inverted by taking the discontinuity across the cut
$`{\displaystyle \frac{\rho (\lambda )}{V}}=\underset{ϵ0}{lim}{\displaystyle \frac{1}{2\pi }}(\mathrm{\Sigma }(i\lambda +ϵ)\mathrm{\Sigma }(i\lambda ϵ)).`$ (11)
In order to obtain a generating function for $`\mathrm{\Sigma }(z)`$ we use a method that has been widely used in the theory of disordered systems , namely
$`Z^{\mathrm{sp}}(z,J,m)={\displaystyle \frac{det(D+z+J)}{det(D+z)}}\stackrel{N_f}{det}(D+m)_{\mathrm{YM}}.`$ (12)
The resolvent is then given by
$`\mathrm{\Sigma }(z)={\displaystyle \frac{1}{V}}_JZ(z,J,m)|_{J=0}.`$ (13)
The spectrum generating function (12) contains $`N_f+1`$ fermionic quarks and one bosonic quark. In addition to the chiral symmetry, this partition function also contains a super-symmetry that mixes fermionic and bosonic quarks. This can be seen by rearranging the fermion determinant as
$`det\left(\begin{array}{cc}m& id\\ id^{}& m\end{array}\right)=det\left(\begin{array}{cc}id& m\\ m& id^{}\end{array}\right).`$ (18)
We observe that for $`m=z=J=0`$ the partition function is invariant under
$`\left(\begin{array}{cccccc}id& & & & & \\ & \mathrm{}& & & & \\ & & id& & & \\ & & & id^{}& & \\ & & & & \mathrm{}& \\ & & & & & id^{}\end{array}\right)\left(\begin{array}{cc}U& \\ & V\end{array}\right)\left(\begin{array}{cccccc}id& & & & & \\ & \mathrm{}& & & & \\ & & id& & & \\ & & & id^{}& & \\ & & & & \mathrm{}& \\ & & & & & id^{}\end{array}\right)\left(\begin{array}{cc}U^1& \\ & V^1\end{array}\right),`$ (35)
where $`U`$ an $`V`$ are $`(N_f+1|1)\times (N_f+1|1)`$ super-matrices. Mathematically, this symmetry group is known as $`Gl(N_f+1|1)\times Gl(N_f+1|1)`$. A $`Gl(1)`$ subgroup is broken by the anomaly. The chiral condensate $`\mathrm{\Sigma }\overline{\psi }\psi `$ is only invariant under the diagonal subgroup with $`U=V`$. Therefore the symmetry is broken spontaneously to $`Gl(N_f+1|1)`$. As is the case in QCD, the masses of the Goldstone modes are given by the Gell-Mann-Oakes-Renner relation
$`{\displaystyle \frac{\sqrt{2m\mathrm{\Sigma }}}{F}},{\displaystyle \frac{\sqrt{(m+z)\mathrm{\Sigma }}}{F}},{\displaystyle \frac{\sqrt{2z\mathrm{\Sigma }}}{F}},`$ (36)
where $`F`$ is the pion decay constant. The Goldstone modes corresponding to a fermionic and a bosonic quark are fermionic whereas all other Goldstone modes are bosonic.
## 4 Effective Low-Energy Theory
The manifold $`Gl(N_f+1|1)`$ is not Riemannian, and is therefore not suitable as a Goldstone manifold. The Goldstone manifold is given by the maximum Riemannian submanifold of the symmetric superspace $`Gl(N_f+1|1)`$ which will be denoted by $`\widehat{Gl}(N_f+1|1)`$ . If we ignore certain complications related to the topological structure of the QCD vacuum, this partition function is given by
$`Z(m,J,z)={\displaystyle _{\widehat{Gl}(N+f+1|1)}}𝑑Ue^{{\scriptscriptstyle d^4x\left[{\scriptscriptstyle \frac{F^2}{4}}\mathrm{Str}_\mu U_\mu U^1{\scriptscriptstyle \frac{1}{2}}\mathrm{\Sigma }\mathrm{Str}M(U+U^1)\right]}},`$ (37)
where the mass matrix $`M=\mathrm{diag}(m,\mathrm{},m,z+J,z)`$.The super-matrix $`U`$ is parameterized as
$`U=\left(\begin{array}{cc}V& \alpha \\ \beta & e^s\end{array}\right)e^{i\sqrt{2}\mathrm{\Phi }/F}.`$ (40)
Here, $`V`$ is a $`U(N_f+1)`$-matrix, $`\alpha `$ and $`\beta `$ are Grassmann valued vectors of length $`N_f+1`$ and $`s`$ is a real number.
In order to estimate the relative importance of the two terms in the effective Lagrangian we expand the fields to second order in the pion fields $`\mathrm{\Phi }=\pi ^at_a`$,
$`^{\mathrm{eff}}={\displaystyle \frac{1}{2}}_\mu \pi ^a_\mu \pi ^a+{\displaystyle \frac{1}{2}}M_a^2\pi _a^2,`$ (41)
where $`M_a`$ is one of the Goldstone masses given in (36). In a box of volume $`L^4`$, the smallest nonzero momenta are of the order $`p_\mu 1/L`$. Let us consider QCD in the chiral limit (with $`m=0`$). Then the regular mesons are massless whereas mesons containing one or two spectral quarks have a nonzero mass given by (36). Therefore, if
$`{\displaystyle \frac{z\mathrm{\Sigma }}{F^2}}{\displaystyle \frac{1}{L^2}}`$ (42)
the correlation functions involving spectral quarks are dominated by contributions from the zero momentum Goldstone modes. Nonperturbatively, the partition function reduces to a group integral in this limit. This integral has been calculated analytically resulting in the following dependence of the condensate on the spectral mass
$`{\displaystyle \frac{\mathrm{\Sigma }(z)}{\mathrm{\Sigma }}}=\mu _z\left[I_{N_f}(\mu _z)K_{N_f}(\mu _z)+I_{N_f+1}(\mu _z)K_{N_f1}(\mu _z)\right].`$ (43)
As already could be observed from the generating function, it depends on $`z`$ only through the combination $`\mu _zzV\mathrm{\Sigma }`$. This result was first obtained by means of chiral Random Matrix Theory to be discussed in the next section.
QCD is not the only theory that reduces to this effective partition function. In fact, any partition function with the same pattern of chiral symmetry breaking and a mass gap reduces to the same low-energy theory. Explicit examples are given by the random flux model and other models with a hopping term, disorder and chiral symmetry . A natural question is what is the simplest theory with the same zero momentum sector as QCD. This theory is chiral Random Matrix Theory.
## 5 Chiral Random Matrix Theory
In the sector of topological charge $`\nu `$ and for $`N_f`$ quarks with mass $`m`$, the chiral random matrix partition function with the global symmetries of the QCD partition function is defined by
$`Z_\beta ^\nu ()={\displaystyle DW\underset{f=1}{\overset{N_f}{}}det\left(\begin{array}{cc}m& iW\\ iW^{}& m\end{array}\right)e^{\frac{N\beta }{4}\mathrm{\Sigma }\mathrm{Tr}W^{}W}},`$ (46)
where $`W`$ is a $`n\times (n+|\nu |)`$ matrix and $`N=2n+|\nu |`$. As is the case in QCD, we assume that the equivalent of the topological charge $`\nu `$ does not exceed $`\sqrt{N}`$, so that, to a good approximation, $`n=N/2`$. Then the parameter $`\mathrm{\Sigma }`$ can be identified as the chiral condensate and $`N`$ as the dimensionless volume of space time (Our units are defined such that the density of the modes $`N/V=1`$). The matrix elements of $`W`$ are either real ($`\beta =1`$, chiral Gaussian Orthogonal Ensemble (chGOE)), complex ($`\beta =2`$, chiral Gaussian Unitary Ensemble (chGUE)), or quaternion real ($`\beta =4`$, chiral Gaussian Symplectic Ensemble (chGSE)). For QCD with three or more colors and quarks in the fundamental representation the matrix elements of the Dirac operator are complex and we have $`\beta =2`$. The ensembles with $`\beta =1`$ and $`\beta =4`$ are relevant in the case of two colors and adjoint fermions, respectively. The reason for choosing a Gaussian distribution of the matrix elements is its mathematical simplicity. It can be shown that the correlations of the eigenvalues on the scale of the level spacing do not depend on the details of the probability distribution .
Together with the Wigner-Dyson ensembles and the superconducting random matrix ensembles the chiral ensembles can be classified according to the Cartan classification of large symmetric spaces .
## 6 Scales in the Dirac Spectrum
For a nonzero value of the chiral condensate $`\mathrm{\Sigma }`$ we can identify three important scales in the Dirac spectrum. The first scale is the smallest nonzero eigenvalue of the Dirac operator given by $`\lambda _{\mathrm{min}}=1/\rho (0)=\pi /\mathrm{\Sigma }V`$. The second scale is the spectral mass for which the Compton wavelength of the associated Goldstone bosons is equal to the size of the box. As discussed above this scale is given by
$`E_c={\displaystyle \frac{F^2}{\mathrm{\Sigma }L^2}}.`$ (47)
In mesoscopic physics this scale is known as the Thouless energy. It is given by the inverse diffusion time of an electron through a sample of length $`L`$. In Euclidean QCD, this scale can be interpreted in terms of diffusive motion of quark in 4 Euclidean dimensions and one artificial time dimension . A third scale is given by a typical hadronic mass scale. The three scales are ordered as $`\lambda _{\mathrm{min}}E_c\mathrm{\Lambda }`$.
For spectral masses $`zE_c`$ the kinetic term in the effective action can be neglected and the low-energy partition function can be reduced to a zero dimensional integral. As already mentioned, the results for $`\mathrm{\Sigma }(z)`$ obtained from the spectral partition function and from chRMT coincide in this domain.
In the domain $`E_cz\mathrm{\Lambda }`$ the kinetic term has to be taken into account. The slope of the Dirac spectrum at $`\lambda =0`$ can be obtained from a one-loop calculation. For $`N_f`$ massless flavors the result is given by
$`{\displaystyle \frac{\rho ^{}(0)}{\rho (0)}}={\displaystyle \frac{(N_f2)(N_f+\beta )}{16\pi \beta }}{\displaystyle \frac{\mathrm{\Sigma }_0}{F^4}}.`$ (48)
Here, $`\beta `$ denotes the Dyson index of the Dirac operator defined before.
The domain below $`E_c`$ has been investigated extensively by means of lattice QCD simulations and complete agreement with the chRMT results has been found <sup>,</sup> . In Fig. 1 we show a comparison of the ratio $`\mathrm{\Sigma }(z)/\mathrm{\Sigma }`$ versus $`zV\mathrm{\Sigma }`$ of lattice results obtained by the Columbia group and eq. (43) for $`N_f=0`$ and 2. The point where the lattice data depart from the chRMT result roughly coincides with the scale $`E_c`$ defined in eq. (47).
## 7 QCD Dirac spectra at $`\mu 0`$
The QCD partition function at $`\mu 0`$ is given by
$`Z={\displaystyle \underset{\alpha }{}}e^{\frac{E_\alpha \mu N_\alpha }{T}}.`$ (49)
Below we will derive the Dirac spectrum from the property that for $`T0`$ only states with $`E_\alpha /N_\alpha <\mu `$ contribute to the partition function.
For $`\mu 0`$ the eigenvalues are scattered in the complex plane. The spectral density is then given by
$`\rho (z)={\displaystyle \frac{1}{\pi }}_z^{}G(z),`$ (50)
with the resolvent $`G(z)`$ defined as usual by
$`G(z)={\displaystyle \frac{1}{V}}\mathrm{Tr}{\displaystyle \frac{1}{D+z}}\stackrel{N_f}{det}(D+m)_{\mathrm{YM}}`$ (51)
Writing out explicitly the real and imaginary parts, we observe that $`G(z)`$ is the electric field in the plane of charges located at the positions of the eigenvalues. The spectrum of the Dirac operator has been studied extensively in the context of Random Matrix Theory <sup>,</sup> . Below, we will obtain the quenched Dirac spectrum using an effective low energy. To lowest order in $`\mu ^2`$ our results agree with quenched Random Matrix Theory.
At $`\mu 0`$ a new ingredient enters in the definition of the the spectrum generating function. In order to assure convergence of the bosonic integral we need an additional factor $`det(D^{}+z^{})`$ in the denominator, and a corresponding factor in the numerator,
$`Z(z,m)={\displaystyle \frac{det(D+z+J)det(D^{}+z^{}+J^{})det^{N_f}(D+m)}{det(D+z)det(D^{}+z^{})}}_{\mathrm{YM}}.`$ (52)
The additional factor can be interpreted in terms of conjugate quarks with a baryon number opposite to that of ordinary quarks . This opens the possibility of baryonic Goldstone modes consisting of a quark and a conjugate anti-quark. For simplicity we only discuss the quenched case where $`N_f=0`$.
In order to construct the effective partition function we have to identify the symmetries of the spectrum generating function. To this end we rewrite the product of the determinant and the conjugate determinant as follows
$`det(D+z)det(D+z^{})`$ $`=`$ $`det\left(\begin{array}{cc}id+\mu & z\\ z& id^{}+\mu \end{array}\right)det\left(\begin{array}{cc}id+\mu & z^{}\\ z^{}& id^{}+\mu \end{array}\right)`$ (57)
$`=`$ $`det\left(\begin{array}{cccc}id+\mu & 0& z& 0\\ 0& id\mu & 0& z^{}\\ z& 0& id^{}+\mu & 0\\ 0& z^{}& 0& id^{}\mu \end{array}\right)`$ (62)
As can be observed from the sign of the chemical potential, the baryon number of the conjugate quarks is opposite to that of the usual quarks.
In this case the chiral symmetry is broken spontaneously according to
$`Gl(2|2)\times Gl(2|2)Gl(2|2).`$ (63)
This symmetry is broken explicitly by the term proportional to $`\mu `$ as
$`Gl(1,1|1,1)\times Gl(1,1|1,1).`$ (64)
This gives rise to an additional term in the effective action. In order to maintain invariance under both the left and right symmetry groups we need at least a term of second order in $`\mu `$ and $`U`$. One easily verifies that to this order, the only static term with the correct symmetry properties is given by
$`{\displaystyle d^4x\mu ^2\mathrm{Str}UT_3U^1T_3},`$ (65)
where the symmetry breaking matrix $`T_3`$ is a diagonal matrix with nonzero matrix elements equal to 1 except on the positions corresponding to the conjugate quarks, where they are -1. For example, for $`N_f=0`$ we have
$`T_3=\left(\begin{array}{cccc}1& & & \\ & 1& & \\ & & 1& \\ & & & 1\end{array}\right).`$ (70)
If we ignore complications related to the topology of the gauge field configurations, the static part of the effective partition function is given by
$`{\displaystyle _{U\widehat{Gl}(2|2)}}e^{\frac{1}{2}\mathrm{\Sigma }V\mathrm{Str}M(U+U^1)\frac{1}{2}V\mu ^2F^2\mathrm{Str}UT_3U^1T_3},`$ (71)
where the mass matrix is given by $`M=\mathrm{diag}(z+J,z^{}+J^{},z,z^{})`$. The integral is over the maximum Riemannian submanifold of $`Gl(2|2)`$. The relative coefficient of the two symmetry breaking terms is determined by the condition that the partition function should have a singularity when $`2\mu `$ becomes equal to the mass of the lightest meson. The coefficient of the term $`\mu ^2`$ has to be chosen such that the effective meson mass vanishes at
$`\mu ^2={\displaystyle \frac{\mathrm{Re}z\mathrm{\Sigma }}{2F^2}},`$ (72)
guaranteeing a nonanlyticity of the $`\mu `$-dependence at this point. The value of this coefficient can be obtained more elegantly by means of a gauge principle. This construction has been performed for QCD with two colors and for QCD with adjoint fermions . In addition to the term proportional to $`\mu ^2`$ we then obtain the coefficients of the terms in the effective Lagrangian that are linear in $`\mu `$ and the momentum.
Baryonic Goldstone modes contain one ordinary quark and one conjugate quark both with mass $`z`$ resulting in a square mass of $`2\mathrm{R}\mathrm{e}z\mathrm{\Sigma }/F^2.`$ For $`\mu ^2<\mathrm{Re}z\mathrm{\Sigma }/2F^2`$ only the vacuum state contributes to the QCD partition function, and thus
$`Z(z,J)=e^{2\mathrm{R}\mathrm{e}J\mathrm{\Sigma }V}.`$ (73)
In the effective theory this corresponds to the saddle-point $`U=1`$. For $`\mu ^2>\mathrm{Re}z\mathrm{\Sigma }/2F^2`$, baryonic Goldstone modes contribute to the partition function resulting in a nonzero baryon density. The density stays finite because of the repulsive interaction between the Goldstone modes. In terms of the effective partition function this happens by the rotation of the saddle point from $`U=1`$ to a nontrivial value. By making the Ansatz for the boson-boson and the fermion-fermion blocks
$`U_{\mathrm{BB}}=U_{FF}=\left(\begin{array}{cc}\mathrm{cos}\theta & \mathrm{sin}\theta \\ \mathrm{sin}\theta & \mathrm{cos}\theta \end{array}\right),`$ (76)
one finds that $`\mathrm{cos}\theta =\mathrm{\Sigma }\mathrm{Re}z/2\mu ^2F^2`$. A similar rotation of the saddle point has been found in the analysis of nonhermitian random matrix models . We then find the following $`J`$-dependence of the partition function
$`Z(z,J)=\mathrm{exp}[{\displaystyle \frac{V\mathrm{\Sigma }^2\mathrm{Re}J\mathrm{Re}z}{\mu ^2F^2}}].`$ (77)
The resolvent in both domains follows by differentiation with respect to the source $`J`$. For $`\mathrm{Re}z>2\mu ^2F^2/\mathrm{\Sigma }`$ we find from (73)
$`G(z)=\mathrm{\Sigma }\mathrm{and}\rho (z)=0,`$ (78)
and for $`\mathrm{Re}z<2\mu ^2F^2/\mathrm{\Sigma }`$ the result for the resolvent following from (77) is given by
$`G(z)={\displaystyle \frac{\mathrm{\Sigma }^2\mathrm{Re}z}{2\mu ^2F^2}}\mathrm{and}\rho (z)={\displaystyle \frac{\mathrm{\Sigma }^2}{4\mu ^2F^2}}.`$ (79)
We observe that the resolvent is continuous at the transition point. The eigenvalues are thus located inside a strip of width $`4F^2\mu ^2/\mathrm{\Sigma }`$. In terms of the interpretation of the resolvent as an electric field, an eigenvalue density that does not depend on the imaginary part of $`z`$ results in a constant electric field outside of the strip of eigenvalues. Ignoring finite size effects this is indeed what is observed in quenched lattice QCD simulations . In the quenched approximation these results are in agreement with an explicit Random Matrix Model calculation. Therefore, to order $`\mu ^2`$ the Random Matrix partition function reproduces exact QCD results.
## 8 Conclusions
We have shown that the correlations of QCD Dirac eigenvalues at a scale well below the hadronic mass scale can be obtained analytically. They are given by an effective theory for the generating function of the resolvent for the QCD Dirac spectrum. For mass scales below $`F^2/\mathrm{\Sigma }L^2`$ the kinetic term in the effective theory can be ignored resulting in eigenvalue correlations given by chRMT. Our results have been confirmed by numerous lattice QCD simulations. The same procedure can be followed for QCD at nonzero chemical potential. Our results indicate that in the quenched limit the eigenvalues are scattered in a strip with a width determined by the lightest meson of our theory. This shows that, to lowest nontrivial order in the chemical potential, chiral Random Matrix Theory provides an exact description of the QCD Dirac spectrum.
Acknowledgments. We gratefully acknowledge all our collaborators in this project. A. Altland, B. Simons, M. Stephanov and M. Zirnbauer are thanked for useful discussions. This work was partially supported by the US DOE grant DE-FG-88ER40388. One of us (D.T.) was supported in part by “Holderbank”-Stiftung and by Janggen-Pöhn-Stiftung. |
warning/0001/math0001171.html | ar5iv | text | # Wavelet filters and infinite-dimensional unitary groups
## 1. Introduction
Digital filters are defined by a choice of coefficients which weight digital time signals from signal processing. It has been known since the eighties that these coefficients are also those which relate the operations of scaling and translation on the real line $``$ in the construction of multiresolution wavelets; see (2.11) and (2.18). There the digital filters are also known as quadrature mirror filters and they are fundamental in generating wavelet bases in $`L^2()`$. They can be determined by various rules which however at times appear somewhat *ad hoc* in the literature \[JiMi91\], \[Dau92\], \[BrJo97b\]. Here we will describe a pairwise bijective correspondence between three sets:
1. the wavelet filters,
2. a certain infinite-dimensional unitary group (also called a loop group),
and finally
1. a specific class of representations on $`L^2(𝕋)`$ of some relations from $`C^{}`$-algebra theory. These last mentioned relations go under the name of the Cuntz relations and are widely studied in operator theory.
The representations in (iii) are not representations of groups, but rather representations of relations. These latter relations have been noted independently in the operator algebra community and in the subband filtering community; see especially \[Cun77\] and \[CMW92a\]. In fact, the present authors stressed in \[BrJo97b\] and \[BrJo98\] that the digital filters which go into the analysis and the synthesis of time signals with subbands are special cases of these representations. In the particular case of quadrature mirror filters (QMF) there are two bands, which are called high-pass/low-pass. In Section 5, we work out the general case of $`N`$ bands, and make explicit the interplay between the separate viewpoints (i)–(iii). This also has the advantage of explaining how wavelet packets (see \[CMW92b\]) arise very naturally from the representation-theoretic method in (iii). Band filtering is described from an engineering viewpoint in \[Vai93\].
The question of reducibility versus irreducibility of the representations in (iii) will be central, as will be unitary equivalence of pairs of representations. Our study of the equivalence question for representations is motivated by a desire for classifying wavelets and systematically constructing interesting examples. The question of reducibility of the representations in turn is motivated by the need for isolating when wavelet data is minimal in a suitable sense. Wavelets come with some specified numerical scaling which may be a natural number, $`N=2,3,\mathrm{}`$, or it may be an expansive integral matrix of some specified dimension $`d`$. (Here we shall restrict to $`d=1`$.) It turns out that the reducibility question for the representations is closely tied to whether or not the scaling data (in this case $`N`$) is minimal in a suitable sense; see Section 6 below.
We also mention that the algebra $`𝒪_N`$ is used in the theory of superselection sectors in quantum field theory \[DoRo89\]. While $`𝒪_N`$ is used there in connection with representation theory, our viewpoint here is completely different: Doplicher and Roberts use in \[DoRo87\] a category of endomorphisms of $`𝒪_N`$ in deriving a noncommutative Pontryagin duality theory, also called Tannaka-Krein theory, in giving an algebraic description of $`\widehat{G}`$ when $`G`$ is a given compact (non-abelian) Lie group. The intertwiners for systems of representations of $`G`$ in \[DoRo89\] and \[DoRo87\] induce the endomorphisms of $`𝒪_N`$. Our viewpoint is in a sense the opposite: We identify a class of representations of $`𝒪_N`$ which has the structure of a compact loop group, i.e., the group of all maps from $`𝕋`$ into $`\mathrm{U}\left(N\right)`$ where $`\mathrm{U}\left(N\right)`$ denotes the group of unitary $`N\times N`$ matrices.
## 2. Background on wavelet filters and extensions
To fix notation, let us give a short rundown of the standard multiresolution wavelet analysis of scale $`N`$. We follow Section 10 in \[BrJo97b\], but see also \[GrMa92\], \[Mey87\], \[MRV96\] and \[Hör95\]. Define scaling by $`N`$ on $`L^2()`$ by
(2.1)
$$\left(U\xi \right)(x)=N^{\frac{1}{2}}\xi \left(N^1x\right),$$
and translation by $`1`$ on $`L^2()`$ by
(2.2)
$$\left(T\xi \right)(x)=\xi (x1).$$
A *scaling function* is a function $`\phi L^2()`$, such that if $`𝒱_0`$ is the closed linear span of all translates $`T^k\phi `$, $`k`$, then $`\phi `$ has the following four properties:
(2.3) $`\{T^k\phi ;k\}`$ is an orthonormal set in $`L^2()`$;
(2.4) $`U\phi 𝒱_0`$;
(2.5) $`_nU^n𝒱_0=\{0\}`$;
(2.6) $`_nU^n𝒱_0=L^2().`$
The simplest example of a scaling function is the Haar function $`\phi `$, which is the characteristic function of the interval $`[0,1]`$.
By (2.3) we may define an isometry
(2.7)
$$_\phi :𝒱_0L^2(𝕋):\xi m$$
as follows:
(2.8)
$$\begin{array}{cc}\xi ()& =_nb_n\phi (n)\hfill \\ _\phi & \\ m(t)& =m(e^{it})=_nb_ne^{int}.\hfill \end{array}$$
Then
(2.9)
$$\widehat{\xi }(t)=m(t)\widehat{\phi }(t)$$
where $`\xi \widehat{\xi }`$ is the Fourier transform. If $`\xi 𝒱_1=U^1𝒱_0`$, then $`U\xi 𝒱_0`$, so we may define
(2.10)
$$m_\xi =_\phi \left(U\xi \right)L^2(𝕋),$$
and then
(2.11)
$$\sqrt{N}\widehat{\xi }(Nt)=m_\xi (t)\widehat{\phi }(t)\text{.}$$
In particular, defining
(2.12)
$$m_0(t)=m_\phi (t)$$
we note that condition (2.3) is equivalent to
(2.13)
$$\mathrm{PER}\left(\left|\widehat{\phi }\right|^2\right)(t):=\underset{k}{}\left|\widehat{\phi }(t+2\pi k)\right|^2=(2\pi )^1\text{,\hspace{1em}\hspace{1em}constancy }\mathrm{a}.\mathrm{e}.t\text{,}$$
which in turn implies
(2.14)
$$\underset{k=0}{\overset{N1}{}}\left|m_0(t+2\pi k/N)\right|^2=N.$$
(Note that (2.14) does *not* conversely imply (2.13). Condition (2.14) merely implies that $`\phi `$ defines a tight frame. In the representation theory in Section 5 this distinction does not play a role. Only the unitarity of $`M(z)`$ in (2.22) is relevant. See \[Hör95, Theorem 3.3.6\]. Similarly, the representation theory in Section 5 applies in cases more general than the ones tied to the $`L^2()`$ wavelets via (2.12) and (2.18). In the latter case one also has the low-pass condition $`m_0(1)=\sqrt{N}`$. In the loop representation (2.23)–(2.24), this is equivalent to $`A_{0,j}(1)=1/\sqrt{N}`$ for $`j=0,\mathrm{},N1`$. These conditions involve the evaluation of the functions $`m_0()`$ or $`A_{0,j}()`$ at $`z=1`$. But $`z=e^{i\omega }`$, where $`\omega `$ is the frequency variable, so the conditions are the low-pass conditions for $`\omega 0.`$)
If $`\xi ,\eta U^1𝒱_0`$, then we have the equivalence:
(2.15)
$$\begin{array}{c}\xi T^k\eta \text{ for all }k\\ \\ _{k=0}^{N1}\overline{m_\xi \left(t+2\pi k/N\right)}m_\eta \left(t+2\pi k/N\right)=0\text{ for almost all }t\text{.}\end{array}$$
If $`\xi U^1𝒱_0`$, then
(2.16)
$$\begin{array}{c}\xi (k),k\text{, is an orthonormal set}\\ \\ _{k=0}^{N1}\left|m_\xi (t+2\pi k/N)\right|^2=N.\end{array}$$
With the *low-pass filter* $`m_0`$ already given, now choose *high-pass filters* $`m_1,\mathrm{},m_{N1}`$ in $`L^2(𝕋)`$ such that
(2.17)
$$\underset{k=0}{\overset{N1}{}}\overline{m_i\left(t+2\pi k/N\right)}m_j(t+2\pi k/N)=\delta _{i,j}N$$
for almost all $`t`$, and define *wavelet functions* $`\psi _1,\mathrm{},\psi _{N1}`$ by
(2.18)
$$\sqrt{N}\widehat{\psi }_i(Nt)=m_i(t)\widehat{\phi }(t).$$
Using (2.15)–(2.18) one then shows that the functions
(2.19)
$$\{T^k\psi _i;k,i=1,\mathrm{},N1\}$$
form an orthonormal basis for $`𝒱_1𝒱_0^{}`$, and thus, using (2.5)–(2.6), the set
(2.20)
$$\left\{U^nT^k\psi _i;n,k,i=1,\mathrm{},N1\right\}$$
forms an orthonormal basis for $`L^2()`$. Concretely, the functions in (2.20) are
(2.21)
$$U^nT^k\psi _i(x)=N^{\frac{n}{2}}\psi _i(N^nxk).$$
Thus, reformulating (2.17), orthonormality of $`\{U^nT^k\psi _i\}`$ is equivalent to unitarity of the matrix
(2.22)
$$M(z)=\frac{1}{\sqrt{N}}\left(\begin{array}{cccc}m_0(z)& m_0(\rho z)& \mathrm{}& m_0\left(\rho ^{N1}z\right)\\ m_1(z)& m_1(\rho z)& \mathrm{}& m_1\left(\rho ^{N1}z\right)\\ \mathrm{}& & & \\ m_{N1}(z)& m_{N1}(\rho z)& \mathrm{}& m_{N1}\left(\rho ^{N1}z\right)\end{array}\right)$$
for almost all $`z𝕋`$, where $`\rho =e^{\frac{2\pi i}{N}}`$ (so it is enough to consider $`z=e^{ix}`$, $`0x<\frac{2\pi }{N}`$). Here $`M`$ is a function from $`𝕋`$ into $`\mathrm{U}(N)`$ of a special kind. It will be convenient for our analysis to consider general functions from $`𝕋`$ into $`\mathrm{U}(N)`$. To this end, we do a Fourier analysis over the cyclic group $`/N`$, and we introduce
(2.23)
$$A_{i,j}(z)=\frac{1}{N}\underset{w^N=z}{}w^jm_i(w),$$
and the inverse transform
(2.24)
$$m_i(z)=\underset{j=0}{\overset{N1}{}}z^jA_{i,j}(z^N).$$
These transforms are also the key to the correspondences (i$``$ (ii$``$ (iii) mentioned in the Introduction. Then (2.24) may be summarized as
(2.25)
$$M(z)=A(z^N)\frac{1}{\sqrt{N}}\left(\begin{array}{cccc}1& 1& \mathrm{}& 1\\ z& \rho z& \mathrm{}& \rho ^{N1}z\\ \mathrm{}& \mathrm{}& & \mathrm{}\\ z^{N1}& \rho ^{N1}z^{N1}& \mathrm{}& \rho ^{(N1)^2}z^{N1}\end{array}\right)\text{;}$$
and the requirement that $`M(z)`$ is unitary for all $`z𝕋`$ is now equivalent to $`A(z)`$ being unitary for all $`z𝕋`$. But $`A`$ is an arbitrary loop. (A *loop* is by definition a function from $`𝕋`$ into $`\mathrm{U}\left(N\right)`$. See the introduction to Section 3.) Hence the formulas (2.23)–(2.24) represent the bijection (i)–(ii) alluded to in the Introduction above.
One main problem considered in this paper is the problem whether the high-pass filters $`m_1,\mathrm{},m_{N1}`$ selected such that (2.17) is valid can be chosen to be “nice” functions when the low-pass filter $`m_0`$ is “nice”, where “nice” may mean continuous, differentiable, polynomial, etc.
Using (2.23), this problem amounts to choosing a unitary matrix $`A(z)`$ once the first row of $`A(z)`$ is given as a unit row vector $`a(z)`$. Ideally, one would like to use a selection map $`F`$ from the unit sphere $`S^{2N1}`$ of $`^N`$ into $`N`$-dimensional orthogonal frames in $`^N`$, i.e.,
(2.26)
$$𝐅(𝐱)=\left(\begin{array}{c}𝐅_0(𝐱)\\ \mathrm{}\\ 𝐅_{N1}(𝐱)\end{array}\right),$$
where the vectors $`𝐅_0(𝐱),\mathrm{},𝐅_{N1}(𝐱)`$ form an orthonormal basis for $`^N`$ for each $`𝐱S^{2N1}`$ and such that $`𝐅_0(𝐱)=𝐱`$ for each $`𝐱`$, and define
(2.27)
$$A(z)=\left(\begin{array}{c}a(z)\\ 𝐅_1(a(z))\\ \mathrm{}\\ 𝐅_{N1}(a(z))\end{array}\right).$$
For example, when $`N=2`$ one may use Daubechies’s choice
(2.28)
$$𝐅_1(x_1,x_2)=(\overline{x}_2,\overline{x}_1).$$
In general one can of course find measurable functions with these properties, but it is remarkable that one cannot choose them to be continuous when $`N>2`$. Following Theorem 4.1 in \[JiSh94\] and Remark 10.2 in \[BrJo97b\], if $`S^{n1}`$ is the unit sphere in $`^n`$, then a theorem of Adams \[Ada62\] says that the highest number of pointwise linearly independent vector fields that can be defined on $`S^{n1}`$ is $`\rho (n)1`$, where the function $`\rho (n)`$ is defined as follows: Let $`b`$ be the multiplicity of $`2`$ in the prime decomposition of $`n`$; write $`b=c+4d`$ where $`c\{0,1,2,3\}`$ and $`d\{0,1,2,\mathrm{}\}`$; and put $`\rho (n)=2^c+8d`$. Thus one verifies $`\rho (2)=2`$, $`\rho (4)=4`$, $`\rho (8)=8`$, $`\rho (16)=9`$, and in general $`\rho (n)<n`$ if $`n\{2,4,8\}`$. Hence the only possibilities of finding a continuous selection function $`F`$ are when $`N=2`$ or $`N=4`$. When $`N=2`$ one can use Daubechies’s selection function (2.28). When $`N=4`$ it is tempting to use the Cayley numbers to construct a selection function, but we will see that this is impossible. Note that if $`\mathrm{U}(N)`$ is replaced by the real orthogonal group $`\mathrm{O}(N)`$, a selection function can be found if and only if $`N\{2,4,8\}`$ by using the complex numbers, the quaternions and the Cayley numbers, respectively. For example, for $`N=4`$, the matrix
(2.29)
$$\left(\begin{array}{cccc}\hfill z_1& \hfill z_2& \hfill z_3& \hfill z_4\\ \hfill z_2& \hfill z_1& \hfill z_4& \hfill z_3\\ \hfill z_3& \hfill z_4& \hfill z_1& \hfill z_2\\ \hfill z_4& \hfill z_3& \hfill z_2& \hfill z_1\end{array}\right)$$
is in $`\mathrm{O}(4)`$ whenever the first row is a unit vector in $`^4`$. Using the identification of $``$ as $`2\times 2`$ real matrices
(2.30)
$$z_1+iz_2\left(\begin{array}{cc}\hfill z_1& \hfill z_2\\ \hfill z_2& \hfill z_1\end{array}\right),$$
this corresponds to Daubechies’s selection function (2.28). If $`N=4`$, a selection function does not exist in the complex case, by Proposition 3.6 in \[Jam76\]. Specializing to the case $`n=k=N`$ in this proposition, it says that, if the space of orthonormal $`N`$-frames over $`^N`$ admits a cross-section over the unit sphere in $`^N`$, and if $`p`$ is any prime, then $`N`$ is divisible by the smallest power of $`p`$ exceeding $`\left(N1\right)/\left(p1\right)`$. If $`N=4`$, this condition fails for $`p=3`$ (while it is fulfilled for all $`p`$ if $`N=2`$). (Note that a direct application of Adams’s theorem only gives nonexistence of the cross-section for $`N>4`$, as is also indicated in the proof of Theorem 4.1 in \[JiSh94\].) (Note also that the discussion above settles the question implied by footnote 3 to the remark before Lemma 3.2.3 in \[Hör95\].) The claim made in \[BrJo97b, Remark 10.2\] that there is such a selection function for $`N=8`$ is definitely erroneous. What is true is that the matrix
(2.31)
$$\left(\begin{array}{cccccccc}\hfill z_1& \hfill z_2& \hfill z_3& \hfill z_4& \hfill z_5& \hfill z_6& \hfill z_7& \hfill z_8\\ \hfill z_2& \hfill z_1& \hfill z_4& \hfill z_3& \hfill z_6& \hfill z_5& \hfill z_8& \hfill z_7\\ \hfill z_3& \hfill z_4& \hfill z_1& \hfill z_2& \hfill z_7& \hfill z_8& \hfill z_5& \hfill z_6\\ \hfill z_4& \hfill z_3& \hfill z_2& \hfill z_1& \hfill z_8& \hfill z_7& \hfill z_6& \hfill z_5\\ \hfill z_5& \hfill z_6& \hfill z_7& \hfill z_8& \hfill z_1& \hfill z_2& \hfill z_3& \hfill z_4\\ \hfill z_6& \hfill z_5& \hfill z_8& \hfill z_7& \hfill z_2& \hfill z_1& \hfill z_4& \hfill z_3\\ \hfill z_7& \hfill z_8& \hfill z_5& \hfill z_6& \hfill z_3& \hfill z_4& \hfill z_1& \hfill z_2\\ \hfill z_8& \hfill z_7& \hfill z_6& \hfill z_5& \hfill z_4& \hfill z_3& \hfill z_2& \hfill z_1\end{array}\right)$$
is in $`\mathrm{O}\left(8\right)`$ whenever the first row is a unit vector in $`^8`$. However, trying to convert this to an element of $`\mathrm{U}\left(4\right)`$ by the same trick as above leads to
(2.32)
$$\left(\begin{array}{cccc}\hfill z_1& \hfill z_2& \hfill z_3& \hfill z_4\\ \hfill \overline{z_2}& \hfill \overline{z_1}& \hfill z_4& \hfill z_3\\ \hfill \overline{z_3}& \hfill z_4& \hfill \overline{z_1}& \hfill z_2\\ \hfill z_4& \hfill \overline{z_3}& \hfill \overline{z_2}& \hfill z_1\end{array}\right),$$
but this is *not* in general unitary when the first row is a unit vector in $`^4`$. For an interesting short account of the matrices (2.29)–(2.31), see \[Tau71\]. The similar selection problem in the context of splines and higer-dimensional spaces $`^s`$ has been treated in \[RiSh91\] and \[RiSh92\]; see in particular pp. 142–145 in \[RiSh91\].
However, there is a completely different method of selecting $`A(z)`$ using more global properties of the first row $`a(z)`$ which applies in the case that $`a`$ is a polynomial in $`z`$. This special case is interesting because it corresponds to compactly supported scaling functions $`\phi `$: the scaling function $`\phi `$ constructed from $`m_0`$ by cascade approximation has support in $`[0,Ng1]`$ if and only if the low-pass filter $`m_0(z)`$ is a polynomial of degree at most $`g1`$, and satisfies $`m_0(1)=\sqrt{N}`$. See \[BrJo97b\], \[BrJo98\], \[BrJo99\], \[BEJ00\], \[ReWe98\], and \[Dau92\]. In this context there is a method of extending $`a(z)`$ to $`A(z)`$ which is independent of the selection theory, and $`A(z)`$ may even be taken to be a polynomial in $`z`$ (with coefficients in $`(^N)`$) of degree equal to the degree of $`a`$. The embryonic form of this method is due to Pollen \[Pol90\], \[Pol92\], but has been developed further: see, e.g., \[ReWe98\], \[Vai93\]. We will give an account of the method in operator-theoretic form, and give the results in Proposition 3.2 and Theorem 4.1. Combining Theorem 4.1 with (2.25) we therefore obtain
###### Theorem 2.1.
Let $`N\{2,3,\mathrm{}\}`$ and let $`\phi `$ be a scaling function with support in $`[0,Ng1]`$ where $`g\{1,2,3,\mathrm{}\}`$. Then one can use multiresolution wavelet analysis to find associated wavelet functions $`\psi _1,\mathrm{},\psi _{N1}`$ also having support in $`[0,Ng1]`$.
###### Proof.
The conditions on $`\phi `$, $`\psi _1,\mathrm{}`$, $`\psi _{N1}`$ correspond to the functions $`m_0`$, $`m_1,\mathrm{}`$, $`m_{N1}`$ being polynomials of degree $`Ng1`$, \[Hör95\], hence Theorem 2.1 follows from Theorem 4.1 and (2.25). More details are given in Corollary 4.2 and its proof. ∎
In Sections 5 and 6 we will consider the representations of the Cuntz algebra $`𝒪_N`$ defined by $`m_0,m_1,\mathrm{},m_{N1}`$ when these are polynomials, and thus extend results from \[BEJ00\].
## 3. Decomposition of polynomial loops into linear factors
We define a *loop* to be a continuous function from the circle $`𝕋`$ into the compact group $`\mathrm{U}(N)`$ of unitary $`N\times N`$ matrices, where $`N=1,2,\mathrm{}`$. The set $`C(𝕋,\mathrm{U}(N))`$ of loops has a natural group structure under pointwise multiplication, and this group is called the *loop group* \[PrSe86\]. We say that a loop is *polynomial* if its matrix elements are polynomials in the basic variable $`z𝕋=\{w;\left|w\right|=1\}`$. The set $`𝒫(𝕋,\mathrm{U}(N))`$ of polynomial loops then forms a semigroup in $`C(𝕋,\mathrm{U}(N))`$ called the *polynomial loop* *semigroup*. (Note that the formulas (2.23)–(2.24) for the bijection between the set of wavelet filters and the loop group are valid also in the wider category of $`L^{\mathrm{}}`$-functions, i.e., when each $`m_i^{(A)}L^{\mathrm{}}(𝕋)`$ and each $`A_{i,j}L^{\mathrm{}}(𝕋)`$. This is the generality of \[BrJo97b\].) Note that if $`U`$ is a given loop, then $`U`$ has a class in the $`K_1`$-group of $`C(𝕋,M_N)`$, i.e., $`K_1\left(C(𝕋,M_N)\right)`$, see \[Bla86\]. This class is called the McMillan degree in the wavelet literature \[ReWe98\], and we will denote it by $`K_1\left(U\right)`$. It can be computed as the winding number of the map
(3.1)
$$𝕋zdet\left(U(z)\right)𝕋\text{.}$$
If $`U`$ is in the loop semigroup, the integer $`K_1(U)`$ is necessarily nonnegative. This follows from Lemma 3.1, below, applied to $`u\left(z\right)=det\left(U\left(z\right)\right)`$. The map $`UK_1(U)`$ is a group homomorphism from $`C(𝕋,\mathrm{U}(N))`$ onto $``$. We need an elementary lemma \[ReWe98\].
###### Lemma 3.1.
If $`u:𝕋𝕋`$ is a polynomial, then $`u`$ is a monomial.
###### Proof.
We have $`u=_{k=0}^nc_kz^k`$ for suitable coefficients $`c_k`$. By multiplying $`u`$ by a nonpositive power of $`z`$ we may assume $`c_00`$. But then $`1=u(z)\overline{u(z)}=c_n\overline{c_0}z^n+`$ powers of $`z`$ between $`(n1)`$ and $`(n1)`$ $`+\overline{c_n}c_0z^n`$ and hence $`c_n=0`$. By induction, $`u(z)=c_0`$. ∎
We now use this to show that any polynomial loop $`U`$ has a decomposition into linear factors \[ReWe98, pp. 60–61\]. The idea in the proof of Lemma 3.3 goes at least back to \[Pol90\] and \[LLS96, Lemma 3.3\].
###### Proposition 3.2.
Let $`U`$ be a polynomial loop in $`𝒫(𝕋,\mathrm{U}(N))`$. The following conditions are equivalent.
(3.2) $`K_1(U)=d`$.
(3.3) There exist one-dimensional projections $`P_1,\mathrm{},P_d`$ on $`^N`$ and a $`V\mathrm{U}(N)`$ such that
$$U(z)=U_1(z)U_2(z)\mathrm{}U_d(z)V,$$
where
$$U_i(z)=zP_i+(1P_i)$$
for $`i=1,\mathrm{},d`$ and $`z𝕋`$. (If $`d=0`$, then $`U(z)=V`$.)
###### Proof.
Since $`K_1(U_i)=1`$ and $`K_1(V)=0`$, the implication (3.3) $``$ (3.2) is obvious. We prove the other implication by using induction with respect to $`d`$. To this end, we will use the following reduction lemma.
###### Lemma 3.3.
Assume that
(3.4)
$$U(z)=A_0+A_1z+\mathrm{}+A_kz^k$$
defines a $`U𝒫(𝕋,\mathrm{U}(N))`$, where $`A_k0`$ and $`k1`$. Then $`U`$ has a unique decomposition
(3.5)
$$U(z)=((1Q)+zQ)(B_0+B_1z+\mathrm{}+B_{k1}z^{k1}),$$
where $`Q`$ is the projection onto the range $`A_k^N`$ of $`A_k`$ (and then $`(1Q)+zQ`$ and $`B_0+B_1z+\mathrm{}+B_{k1}z^{k1}`$ are in $`𝒫(𝕋,\mathrm{U}(N))`$).
###### Proof.
Let $`Q`$ be the projection onto $`A_k^N`$. We have
(3.6) $`1`$ $`=U(z)^{}U(z)`$
$`=(A_0^{}+A_1^{}z^1+\mathrm{}+A_k^{}z^k)(A_0+A_1z+\mathrm{}+A_kz^k)`$
$`=A_k^{}A_0z^k+\text{ terms in higher powers of }z\text{.}`$
Hence
(3.7)
$$A_k^{}A_0=0,$$
and thus, since $`Q`$ is the projection onto the range of $`A_k`$, we obtain
(3.8)
$$QA_0=0\text{ and }(1Q)A_k=0\text{.}$$
Now put
(3.9) $`W(z)`$ $`=((1Q)+zQ)^1U(z)`$
$`=((1Q)+z^1Q)(A_0+\mathrm{}+z^kA_k)`$
$`=z^1QA_0`$
$`+((1Q)A_0+QA_1)`$
$`+\mathrm{}`$
$`+z^k(1Q)A_k\text{ (ordered by increasing powers).}`$
It follows from (3.8) that $`W(z)`$ has the form
(3.10)
$$W(z)=B_0+B_1z+\mathrm{}+B_{k1}z^{k1},$$
and (3.5) follows. (In fact, $`B_0=\left(1Q\right)A_0+QA_1=A_0+QA_1`$, $`\mathrm{}`$, and $`B_{k1}=\left(1Q\right)A_{k1}+A_k`$.) By our choice of $`Q`$, the decomposition (3.5) is unique by (3.8). ∎
###### End of the proof of Proposition 3.2.
We will prove the proposition by induction with respect to the McMillan index $`d`$. Assume first that $`d1`$ and that the proposition holds for all indices $`d1`$, and suppose that
(3.11)
$$U(z)=A_0+A_1z+\mathrm{}+A_kz^k$$
has index $`d`$. But by Lemma 3.3, $`U`$ has a decomposition
(3.12)
$$U(z)=((1Q)+zQ)W(z),$$
where $`Q`$ is a nonzero projection. But then
(3.13) $`d`$ $`=K_1(U)=K_1((1Q)+zQ)+K_1(W(z))`$
$`=dim(Q)+K_1(W(z)),`$
so $`K_1(W(z))d1`$. (We define $`dimQ=rankQ`$ when $`Q`$ is a projection.) We may then apply the induction hypothesis to $`W(z)`$. Finally, there exist $`dimQ`$ one-dimensional projections $`P_1,\mathrm{},P_{dim(Q)}`$ such that $`Q=P_1+P_2+\mathrm{}+P_{dim(Q)}`$, and then
(3.14)
$$((1Q)+zQ)=\underset{n=1}{\overset{dimQ}{}}((1P_n)+zP_n).$$
It follows that $`U`$ has the form (3.3).
Finally, if $`K_1(U)=0`$ and $`U`$ has the form (3.11), there are two possibilities:
1. $`k=0`$ and $`U`$ is a constant unitary matrix. Then $`U`$ already has the form in (3.3) with $`d=0`$.
2. $`k>0`$. Then one may apply Lemma 3.3 to write $`U(z)=((1Q)+zQ)W(z)`$ with $`Q`$ a nonzero projection, but as
$$0=K_1(U)=dimQ+K_1(W)$$
and $`K_1(W)0`$, this is impossible.
This ends the proof of Proposition 3.2. The last argument in the induction chain is simpler if we use induction with respect to $`k`$ rather than induction with respect to $`d`$. ∎
## 4. Extensions of low-pass polynomial filters to high-pass filters
In this section we will prove Theorem 2.1. This theorem follows from Corollary 4.2, below. Theorem 4.1 states that every vector $`\alpha `$ in $`^{Ng}`$ which satisfies a certain orthogonality condition is the first row of coefficients of some element of the loop group.
###### Theorem 4.1.
Let $`\alpha =(\alpha _0,\alpha _1,\mathrm{},\alpha _{g1})`$ be $`g`$ row vectors in $`^N`$. The following conditions are equivalent.
(4.1) The vectors satisfy the relations
$$\underset{i}{}\alpha _i\text{ }\alpha _{i+j}=0\text{ for }j0,$$
and
$$\underset{i}{}\alpha _i\text{ }\alpha _i=1,$$
where we use the convention that $`\alpha _i=0`$ if $`i<0`$ or $`ig`$.
(4.2) There exists a polynomial loop
$$A\left(z\right)𝒫(𝕋,\mathrm{U}\left(N\right))$$
of degree $`g1`$ such that the first row of $`A\left(z\right)`$ is
$$\underset{i=0}{\overset{g1}{}}z^i\alpha _i.$$
###### Proof.
The implication (4.2) $``$ (4.1) follows from considering the $`(0,0)`$ matrix entry of
(4.3)
$$A\left(z\right)A\left(z\right)^{}=1$$
in $`M_N`$ for all $`z`$, i.e., $`A\left(z\right)A\left(z\right)^{}`$ is the constant Laurent polynomial $`1`$.
The other implication is proved by a very similar method as in the proof of Proposition 3.2. Again we use induction with respect to $`g`$.
If $`g=1`$, condition (4.1) just says that $`\alpha _0`$ is a unit row vector, and we can find a constant function $`A\left(z\right)=A`$ just by Gram-Schmidt orthogonalization.
Assume that $`g>1`$, and that the result has been proved for all smaller $`g`$. We may assume that $`\alpha _{g1}0`$ (otherwise we are already through by the induction hypothesis). Let $`P`$ be the one-dimensional projection onto $`\alpha _{g1}`$, i.e.,
(4.4)
$$\alpha _{g1}P=\alpha _{g1}.$$
Now define
(4.5) $`\beta \left(z\right)`$ $`=\alpha \left(z\right)\left(1P+z^1P\right)`$
$`=\left(\alpha _0+\alpha _1z+\mathrm{}+\alpha _{g1}z^{g1}\right)\left(1P+z^1P\right)`$
$`=z^1\alpha _0P`$
$`+\alpha _0\left(1P\right)+\alpha _1P`$
$`+z\left(\mathrm{}\right)`$
$`+\mathrm{}`$
$`+z^{g2}\left(\alpha _{g1}P+\alpha _{g2}\left(1P\right)\right)`$
$`+z^{g1}\alpha _{g1}\left(1P\right).`$
But since $`\alpha _0\text{ }\alpha _{g1}=0`$ by (4.1), the $`z^1`$ term disappears, and (4.4) implies that the $`z^{g1}`$ term disappears. Therefore $`\beta \left(z\right)`$ is a polynomial in $`z`$ of degree $`g2`$. One now verifies from the unitarity of $`\left(1P+z^1P\right)`$ that the coefficient vectors of $`\beta `$ satisfy the same relations (4.1) as those of $`\alpha `$. Hence, applying the induction hypothesis, there exists a polynomial loop $`B\left(z\right)`$ of degree $`g2`$ such that the first row of $`B\left(z\right)`$ is $`\beta \left(z\right)`$. Thus, putting
(4.6)
$$A\left(z\right)=B\left(z\right)\left(1P+zP\right),$$
it follows from (4.5) that the first row of $`A\left(z\right)`$ is $`\alpha \left(z\right)`$. This completes the induction, and thus the proof of Theorem 4.1. ∎
Let us restate Theorem 2.1 and its proof in terms of the low-pass filter $`m_0(z)`$:
###### Corollary 4.2.
Let $`m_0`$ be a polynomial, and let $`N\{2,3,\mathrm{}\}`$. Suppose
(4.7)
$$\left|m_0\left(z\right)\right|^2+\left|m_0\left(z\rho \right)\right|^2+\mathrm{}+\left|m_0\left(z\rho ^{N1}\right)\right|^2=N$$
for all $`z𝕋`$, where $`\rho =e^{\frac{2\pi i}{N}}`$. Then there are polynomials $`\left\{m_i;i=1,\mathrm{},N1\right\}`$ such that the combined system $`\left\{m_i;0i<N\right\}`$ satisfies the unitarity property of (2.22), or equivalently (2.17) with the convention $`ze^{it}`$. In other words, every $`m_0`$ may be completed to a quadrature mirror filter system.
Furthermore, when $`g`$ is chosen such that the degree of the polynomial $`m_0`$ is at most $`Ng1`$, then the polynomials $`m_1,\mathrm{},m_{N1}`$ can be taken to have degree at most $`Ng1`$.
###### Proof.
With $`m_0`$ and $`N`$ given as stated in the corollary, set
(4.8)
$$A_{0,j}\left(z\right):=\frac{1}{N}\underset{w^N=z}{}w^jm_0\left(w\right),z𝕋.$$
Then it follows from (4.7) that the coefficients $`\alpha _j`$, as a set of row vectors in $`^N`$, satisfy (4.1) in Theorem 4.1. Picking then $`A\left(z\right)=\left(A_{i,j}\left(z\right)\right)𝒫(𝕋,\mathrm{U}\left(N\right))`$ as in (4.2), using the theorem, and setting
(4.9)
$$m_i\left(z\right)=\underset{j=0}{\overset{N1}{}}z^jA_{i,j}\left(z^N\right),$$
see (2.24), we note that these functions satisfy the conclusion in the corollary.
Finally, we note that, as in Remark 5.6, if $`m_0`$ has degree at most $`Ng1`$, it follows from (4.8) that $`A_{0,j}\left(z\right)`$ has degree at most $`g1`$; thus $`A_{i,j}\left(z\right)`$ can be taken to have degree at most $`g1`$ by Theorem 4.1, and then all $`m_i\left(z\right)`$ have degree at most $`Ng1`$ by Remark 5.6 or (4.9). ∎
## 5. Representations associated with polynomial filters
If $`m_0,m_1,\mathrm{},m_{N1}`$ are complex functions on $`𝕋`$ such that the matrix $`M\left(z\right)`$ in (2.22) is unitary, one may define isometries $`S_0,\mathrm{},S_{N1}`$ on $`L^2\left(𝕋\right)`$ by
(5.1)
$$\left(S_i\xi \right)\left(z\right)=m_i\left(z\right)\xi \left(z^N\right)$$
and then
(5.2)
$$\left(S_i^{}\xi \right)\left(z\right)=\frac{1}{N}\underset{\begin{array}{c}w\\ w^N=z\end{array}}{}\overline{m}_i\left(w\right)\xi \left(w\right),$$
see \[BrJo97b, (1.16)–(1.17)\] and, for $`N=2`$, \[CMW92a, (1.1)\]. The isometries $`S_0,\mathrm{},S_{N1}`$ then have mutually orthogonal ranges,
(5.3)
$$S_i^{}S_j=\delta _{i,j}\text{1}\text{1},$$
and these ranges add up to 11,
(5.4)
$$\underset{i=0}{\overset{N1}{}}S_iS_i^{}=\text{1}\text{1}.$$
The relations (5.3)–(5.4) are called the *Cuntz relations* (of order $`N`$), and the universal $`C^{}`$-algebra generated by operators $`s_0,\mathrm{},s_{N1}`$ satisfying these relations is called the *Cuntz algebra* (of order $`N`$). It is a simple antiliminal $`C^{}`$-algebra, which means that its space of Hilbert-space realizations cannot be equipped with a standard Borel structure (see \[Cun77\]). In \[BrJo97b\] and Section 2 above, we identify intertwining operators $`W=_\phi ^{}`$ from $`L^2\left(𝕋\right)`$ onto resolution subspaces in $`L^2\left(\right)`$ which intertwine the isometry $`S_0`$ with the scaling operator $`fN^{1/2}f\left(x/N\right)`$ on the appropriate resolution subspace in $`L^2\left(\right)`$. The family of subspaces
(5.5)
$$S_{i_1}S_{i_2}\mathrm{}S_{i_k}L^2\left(𝕋\right),k=0,1,\mathrm{},i_j\{0,1,\mathrm{},N1\},$$
corresponds under $`W`$ to wavelet packets (see \[CMW92b\]) in $`L^2\left(\right)`$. Relations (5.3)–(5.4) make it immediately clear that the respective projections onto the subspaces (5.5) are $`S_{i_1}\mathrm{}S_{i_k}S_{i_k}^{}\mathrm{}S_{i_1}^{}`$. In particular, these projections are clearly mutually orthogonal when $`k`$ is fixed and multi-indices are varied.
Let us recall some results from \[BJKW00\] and \[BEJ00\]. Let $``$ be a Hilbert space, and $`s_i\pi \left(s_i\right)=S_i\left(\right)`$ a representation of the Cuntz relations on $``$. Assume that there exists a finite-dimensional subspace $`𝒦`$ with the properties
(5.6)
$$S_i^{}𝒦𝒦,i=0,1,\mathrm{},N1,$$
and
(5.7)
$$𝒦\text{ is cyclic, i.e., }\left[\pi \left(𝒪_N\right)𝒦\right]=$$
(here $`\left[\pi \left(𝒪_N\right)𝒦\right]`$ denotes the closure of the linear span of all polynomials in$`S_0,\mathrm{},S_{N1}`$ applied to vectors in $`𝒦`$). Define operators $`V_i\left(𝒦\right)`$ by
(5.8)
$$V_i^{}=S_i^{}|_𝒦.$$
Then
(5.9)
$$\underset{i=0}{\overset{N1}{}}V_iV_i^{}=\text{1}\text{1}_𝒦.$$
Define a map $`\sigma `$ on $`\left(𝒦\right)`$ by
(5.10)
$$\sigma \left(A\right)=\underset{i=0}{\overset{N1}{}}V_iAV_i^{}.$$
Then $`\sigma `$ is a unital completely positive map, and
###### Theorem 5.1.
(\[BJKW00\]) There is a positive norm-preserving linear isomorphism between the commutant algebra
(5.11)
$$\pi \left(𝒪_N\right)^{}=\left\{A\left(\right);A\pi \left(x\right)=\pi \left(x\right)A\text{ for all }x𝒪_N\right\}$$
and the fixed-point set
(5.12)
$$\left(𝒦\right)^\sigma =\left\{A\left(𝒦\right);\sigma \left(A\right)=A\right\}$$
given by
(5.13)
$$\pi \left(𝒪_N\right)^{}APAP,$$
where $`P`$ is the projection of $``$ onto $`𝒦`$. In particular, $`\pi `$ is irreducible if and only if $`\sigma `$ is ergodic.
More generally, if $`𝒦_1`$, $`𝒦_2`$ (with corresponding projections $`P^{\left(1\right)}`$ and $`P^{\left(2\right)}`$) are $`S^{}`$-invariant cyclic subspaces for two representations $`\pi _1`$, $`\pi _2`$ of $`𝒪_N`$ on $`_1`$, $`_2`$, and
(5.14)
$$V_i^{\left(j\right)}=P^{\left(j\right)}\pi _j\left(s_i\right)|_{𝒦_j}$$
for $`j=1,2`$, $`i=0,\mathrm{},N1`$, define $`\rho `$ on $`(𝒦_1,𝒦_2)`$ by
(5.15)
$$\rho \left(A\right)=\underset{i}{}V_i^{\left(2\right)}AV_i^{\left(1\right)}.$$
Then there is an isometric linear isomorphism between the set of intertwiners
(5.16)
$$\left\{A(_1,_2);A\pi _1\left(x\right)=\pi _2\left(x\right)A\text{ for all }x𝒪_N\right\}$$
and the fixed-point set
(5.17)
$$\left\{B(𝒦_1,𝒦_2);\rho \left(B\right)=B\right\}$$
given by
(5.18)
$$AB=P^{\left(2\right)}AP^{\left(1\right)}.$$
We argued in \[BEJ00, (4.14)–(4.18)\] that if $`m_0,\mathrm{},m_{N1}`$ are polynomials in the circle variable $`z`$, $`N=2`$, and $`=L^2\left(𝕋\right)`$, then such a finite-dimensional subspace $`𝒦`$ exists, having dimension $`Ng`$ and spanned by $`1,z^1,z^2,\mathrm{},z^{Ng+1}`$. In this section we will extend this result to $`N>2`$, and we will see in Proposition 5.5 that we can do slightly better than what the $`N=2`$ result indicates. Let us assume from now on that the loop group element $`A\left(z\right)`$ in (2.23)–(2.25) is a polynomial of degree $`g1`$, so that $`m_0\left(z\right),\mathrm{},m_{N1}\left(z\right)`$ are polynomials of degree at most $`N\left(g1\right)+N1=Ng1`$.
We will use the notation $`A\left(z\right)=\left(A_{i,j}\left(z\right)\right)_{i,j=0}^{N1}`$ for the loop-group element $`A:𝕋\mathrm{U}\left(N\right)`$. Since the Fourier expansion is finite, $`A\left(z\right)`$ has the form
(5.19)
$$A\left(z\right)=\underset{k=0}{\overset{g1}{}}z^kA^{\left(k\right)},$$
where $`A^{\left(k\right)}\left(^N\right)`$ for $`k=0,\mathrm{},g1`$. The factorization in Lemma 3.3 motivates the name *genus* for $`g`$.
We have
###### Corollary 5.2.
If $`A\left(z\right)`$ is a general polynomial of $`z`$ with values in $`\left(^N\right)`$ of the form (5.19), the following four conditions (5.20)–(5.23) are equivalent:
(5.20)
(5.21)
(5.22)
and
(5.23) there are projections $`Q_1,Q_2,\mathrm{},Q_{g1}`$ and a unitary $`V\mathrm{U}\left(N\right)`$ such that
$`A^{\left(0\right)}`$ $`=V{\displaystyle \underset{j=1}{\overset{g1}{}}}\left(\text{1}\text{1}_NQ_j\right),`$
$`A^{\left(1\right)}`$ $`=V{\displaystyle \underset{j=1}{\overset{g1}{}}}\left(\text{1}\text{1}_NQ_1\right)\mathrm{}`$
$`\mathrm{}\left(\text{1}\text{1}_NQ_{j1}\right)Q_j\left(\text{1}\text{1}_NQ_{j+1}\right)\mathrm{}`$
$`\mathrm{}\left(\text{1}\text{1}_NQ_{g1}\right),`$
$`\mathrm{}`$ $`\mathrm{}`$
$`A^{\left(g1\right)}`$ $`=V{\displaystyle \underset{j=1}{\overset{g1}{}}}Q_j.`$
###### Proof.
This follows from Proposition 3.2 and explicit computations. ∎
###### Remark 5.3.
The case $`g=2=N`$ includes the family of wavelets introduced by Daubechies \[Dau92\] and studied further in \[BEJ00\]. Note that $`g=2`$ yields the representation
(5.24)
$$A^{\left(0\right)}=V\left(\text{1}\text{1}_NQ\right),A^{\left(1\right)}=VQ,$$
by (5.23). But then (5.21) takes the form
(5.25)
$$A^{\left(0\right)}A^{\left(0\right)}=\text{1}\text{1}_NQ,A^{\left(1\right)}A^{\left(1\right)}=Q,$$
which will be used in the sequel.
Let us return to the connection between loop-group elements and representations of $`𝒪_N`$. The algebra $`𝒪_N`$ has the following basic representation $`s_iS_i`$ on $`L^2\left(𝕋\right)`$:
(5.26)
$$S_j\xi \left(z\right)=z^j\xi \left(z^N\right),0j<N,\xi L^2\left(𝕋\right),z𝕋\text{.}$$
If $`s_iT_i`$ is any other representation, then as noted in \[Cun77\], \[Cun80\], and \[BJP96\], the system $`\left(S_j^{}T_i\right)_{i,j}`$ realizes a unitary $`N\times N`$ matrix with entries in $`\left(L^2\left(𝕋\right)\right)`$.
###### Proposition 5.4.
Let $`\left(S_i\right)`$ be the basic representation of $`𝒪_N`$, and let $`\left(T_i\right)`$ be an arbitrary representation. Then the following two conditions are equivalent:
(5.27) Each operator $`S_j^{}T_i`$ on $`L^2\left(𝕋\right)`$ is a multiplication operator;
and
(5.28) the $`T_i`$-representation has the form
$$T_i\xi \left(z\right)=m_i\left(z\right)\xi \left(z^N\right),0i<N,\xi L^2\left(𝕋\right),z𝕋,$$
where $`m_0,\mathrm{},m_{N1}L^{\mathrm{}}\left(𝕋\right)`$.
###### Proof.
(5.27) $``$ (5.28): If (5.27) holds, then there are functions $`A_{i,j}L^{\mathrm{}}\left(𝕋\right)`$ such that $`S_j^{}T_i=A_{i,j}`$ where the right-hand side also denotes the multiplication operator determined by the function in question. Then
(5.29)
$$T_i=\underset{j=0}{\overset{N1}{}}S_jS_j^{}T_i=\underset{j=0}{\overset{N1}{}}S_jA_{i,j}=\underset{j=0}{\overset{N1}{}}A_{i,j}\left(z^N\right)S_j.$$
Setting
(5.30)
$$m_i\left(z\right)=\underset{j=0}{\overset{N1}{}}A_{i,j}\left(z^N\right)z^j$$
and using the Cuntz relations, we conclude that $`T_i`$ has the form in (5.28).
(5.28) $``$ (5.27): If, conversely, the representation $`\left(T_i\right)`$ is given to have the form in (5.28), then one checks that the filter functions $`m_i`$ must satisfy the unitarity property (2.22) above. If then $`A_{i,j}\left(z\right)`$ are given by (2.23), then
(5.31)
$$S_j^{}T_i\xi \left(z\right)=\frac{1}{N}\underset{w^N=z}{}w^jT_i\xi \left(w\right)=\frac{1}{N}\underset{w^N=z}{}w^jm_i\left(w\right)\xi \left(z\right),$$
and we conclude that (5.27) holds, i.e., $`S_j^{}T_i`$ is multiplication by the function $`A_{i,j}()`$, where $`A_{i,j}`$ is derived from $`m_i`$ via (2.23) ∎
We will henceforth denote the representation $`T_i`$ defined by (5.31) by $`T_i^{\left(A\right)}`$, so in particular,
(5.32)
$$T_i^{\left(\text{1}\text{1}\right)}=S_i,$$
where the notation $`S_i`$ is reserved for the representation defined by (5.26).
We will now apply the results in (5.6)–(5.18) to analyze these representations further. As mentioned after (5.18), the linear span $`𝒦`$ of $`1,z^1,z^2,\mathrm{},z^r`$ is cyclic and $`T^{}`$-invariant for a suitable $`r`$, i.e., satisfies (5.6)–(5.7) with $`T`$ in lieu of $`S`$, see \[BEJ00, Proposition 3.1 and Corollary 3.3\]. If $`m_0,\mathrm{},m_{N1}`$ can be derived from a polynomial loop $`A\left(z\right)`$ of degree $`g1`$ by (2.24), we may explicitly estimate $`r`$. To this end, let us look at the action of $`T_i^{\left(A\right)}`$ on $`e_n\left(z\right)=z^n`$. Put $`n=j+Nl`$ uniquely, where $`j\{0,1,\mathrm{},N1\}`$ and $`l`$. Then, using (5.29),
(5.33) $`T_i^{\left(A\right)}e_n`$ $`={\displaystyle \underset{j^{}}{}}S_j^{}^{}\overline{A_{i,j^{}}\left(z^N\right)}e_n`$
$`={\displaystyle \underset{k}{}}{\displaystyle \underset{j^{}}{}}S_j^{}^{}\overline{A_{i,j^{}}^{\left(k\right)}}z^{Nk}e_n={\displaystyle \underset{k}{}}{\displaystyle \underset{j^{}}{}}\overline{A_{i,j^{}}^{\left(k\right)}}S_j^{}^{}e_{nNk}.`$
But $`nNk=j+NlNk=j+N\left(lk\right)`$, and since
(5.34)
$$S_j^{}^{}e_{j+N\left(lk\right)}=\{\begin{array}{cc}e_{lk}\hfill & \text{ if }j^{}=j,\hfill \\ 0\hfill & \text{ if }j^{}j,\hfill \end{array}$$
we obtain
(5.35)
$$T_i^{\left(A\right)}e_{j+Nl}=\underset{k=0}{\overset{g1}{}}\overline{A_{i,j}^{\left(k\right)}}e_{lk}.$$
It follows that $`T_i^{\left(A\right)}`$ is represented by a matrix which is a slanted block matrix of the following form:
(5.36)
$$\begin{array}{cc}& \text{}\end{array}$$
The matrix elements outside the shaded blocks are all zero, and each block has $`N`$ columns and $`g`$ rows. Each block is a translate $`N`$ steps to the right and $`1`$ step up compared to the previous one,and the most central block is located with corners at $`(0,0)`$, $`(N1,0)`$, $`(N1,g+1)`$, $`(0,g+1)`$. It follows that if we fix one $`n`$, iterated applications of $`T_i^{\left(A\right)}`$ for various $`i`$’s ultimately will transform this vector into a linear combination of $`e_m`$’s where $`m\{0,1,2,\mathrm{},r\}`$, where $`r`$ is the largest integer such that $`(r,r)`$ is contained in one of the blocks. Thus the linear span $`𝒦`$ of $`\{e_0,e_1,\mathrm{},e_r\}`$ is $`T_i^{\left(A\right)}`$-invariant, and because of the relation
(5.37)
$$\text{1}\text{1}=\underset{\begin{array}{c}I\\ \left|I\right|=n\end{array}}{}T_I^{\left(A\right)}T_I^{\left(A\right)},$$
$`𝒦`$ will also be cyclic (see \[BEJ00, Section 3\] for more details of this argument). Thus, computing $`r`$ more explicitly, we deduce
###### Proposition 5.5.
If $`T^{\left(A\right)}`$ is the representation of the Cuntz algebra $`𝒪_N`$ defined by a polynomial loop $`A\left(z\right)`$ of degree $`g1`$, then the subspace
(5.38)
$$𝒦=linspan\{e_0,e_1,e_2,\mathrm{},e_r\}$$
satisfies (5.6) and (5.7), i.e., $`𝒦`$ is $`T^{\left(A\right)}`$-invariant and cyclic, where
(5.39)
$$r=g+\frac{g1}{N1}=\frac{gN1}{N1}.$$
Here $`x`$ is the largest integer $`x`$.
###### Proof.
From the figure (5.36) it follows that
(5.40)
$$r=\{\begin{array}{ccc}g\hfill & \text{if }g<N,\hfill & \text{i.e., }g1<N1,\hfill \\ g+1\hfill & \text{if }N+1g+1<2N,\hfill & \text{i.e., }N1g1<2N2,\hfill \\ g+2\hfill & \text{if }2N+1g+2<3N,\hfill & \text{i.e., }2N2g1<3N3,\hfill \\ \text{etc.}\hfill & & \end{array}$$
Here the first column of ranges (“if $`\mathrm{}`$”) is taken from the box diagram (5.36), and the second column of ranges (“i.e., $`\mathrm{}`$”) is derived from the first by subtraction of suitable integers. The second column shows how the ratio $`\frac{g1}{N1}`$ arises in (5.39). ∎
###### Remark 5.6.
Proposition 5.5 could alternatively have been proved by exactly the method used to show (4.18) in \[BEJ00\] from Remark 3.2 there, using joint invariant sets for the maps $`n\left(nk\right)/N`$ for $`k=0,1,\mathrm{},gN1`$. Note that a consequence of (2.23) and (2.24) is that the degree of all the polynomials $`A_{i,j}\left(z\right)`$ for $`i,j=0,\mathrm{},N1`$ is at most $`g1`$ if and only if the degree of all the polynomials $`m_i\left(z\right)`$ for $`i=0,\mathrm{},N1`$ is at most $`gN1`$.
The case considered in detail in \[BEJ00\] was $`N=g=2`$, so we regain the result $`r=3`$ from there, i.e., in that case $`𝒦=span\{e_0,e_1,e_2,e_3\}`$. Note that the denominator $`N1`$ in (5.39) shows the dependence of $`r=\left(dim𝒦\right)1`$ on the scaling number $`N`$ for the given wavelet filter. So when $`N>2`$, the denominator $`N1`$ in (5.39) yields a smaller value for $`r`$ than the value $`gN1`$ from \[BEJ00\] in the special case $`N=2`$.
We will now apply Theorem 5.1 to the representations above, when $`𝒦=𝒦^{\left(A\right)}`$ is given by (5.38) and $`r=r^{\left(A\right)}`$ by (5.39). The operators $`V_i\left(𝒦^{\left(A\right)}\right)`$ defined by $`T_i^{\left(A\right)}`$ in lieu of $`S_i`$ in (5.8) will be denoted by $`V_i^{\left(A\right)}`$, and then the $`\sigma `$ defined by (5.10) is denoted by $`\sigma ^{\left(A\right)}=\sigma ^{(A,A)}`$, i.e.,
(5.41)
$$\sigma ^{(A,A)}\left(X\right)=\underset{i=0}{\overset{N1}{}}V_i^{\left(A\right)}XV_i^{\left(A\right)}$$
for $`X\left(𝒦^{\left(A\right)}\right)`$. If $`T^{\left(A\right)}`$, $`T^{\left(B\right)}`$ are two representations of our kind, let $`\sigma ^{(B,A)}`$ denote the $`\rho `$ defined by (5.15), i.e.,
(5.42)
$$\sigma ^{(B,A)}\left(X\right)=\underset{i}{}V_i^{\left(B\right)}XV_i^{\left(A\right)},$$
so $`\sigma ^{(B,A)}`$ is in $`(𝒦^{\left(A\right)},𝒦^{\left(B\right)})`$. We will mostly be interested in the situation that $`A`$ and $`B`$ have the same genus $`g`$, and then $`𝒦^{\left(A\right)}=𝒦^{\left(B\right)}`$ by Proposition 5.5 (of course we may always replace $`g^{\left(A\right)}`$, $`g^{\left(B\right)}`$ by $`\mathrm{max}\{g^{\left(A\right)},g^{\left(B\right)}\}`$ and thus assume $`𝒦^{\left(A\right)}=𝒦^{\left(B\right)}`$). Now $`(𝒦^{\left(A\right)},𝒦^{\left(B\right)})`$ can be made into a Hilbert space $`𝒦^{(B,A)}`$ in a natural fashion by defining the norm of $`\psi (𝒦^{\left(A\right)},𝒦^{\left(B\right)})`$ by
(5.43)
$$\psi _2^{(B,A)}=\stackrel{\left(A\right)}{Tr}\left(\psi ^{}\psi \right)=\stackrel{\left(B\right)}{Tr}\left(\psi \psi ^{}\right),$$
where $`\stackrel{\left(A\right)}{Tr}`$, $`\stackrel{\left(B\right)}{Tr}`$ are the standard unnormalized traces on $`\left(𝒦^{\left(A\right)}\right)`$, $`\left(𝒦^{\left(B\right)}\right)`$, respectively. Then $`\sigma ^{(B,A)}\left(𝒦^{(B,A)}\right)`$. Because of Theorem 5.1, our main concern will be the spectral multiplicity of $`1`$ in the spectrum $`Sp\left(\sigma ^{(B,A)}\right)`$, as well as the associated eigensubspace $`𝒦^{(B,A)}`$. It turns out that these items are much simpler to compute for the adjoint $`\sigma ^{(B,A)}\left(𝒦^{(B,A)}\right)`$ (i.e., the adjoint of $`\sigma ^{(B,A)}`$ as a Hilbert-space operator). Then $`Sp\left(\sigma ^{(B,A)}\right)=\overline{Sp\left(\sigma ^{(B,A)}\right)}`$, and in particular $`1Sp\left(\sigma ^{(B,A)}\right)`$ if and only if $`1Sp\left(\sigma ^{(B,A)}\right)`$. If $`A=B`$ and the completely positive unital map $`\sigma ^{(A,A)}`$ admits a faithful invariant state, it follows from \[BrJo97a, Lemma 6.3\] that the dimensions of the fixed-point sets of $`\sigma ^{(A,A)}`$ and $`\sigma ^{(A,A)}`$ are the same. However, we will see that $`\sigma ^{(A,A)}`$ does not necessarily admit a faithful invariant state. Our strategy will be to compute $`Sp\left(\sigma ^{(B,A)}\right)`$ by evaluating
(5.44)
$$det\left(x\text{1}\text{1}\sigma ^{(A,B)}\right)=\overline{det\left(\overline{x}\text{1}\text{1}\sigma ^{(A,B)}\right)}$$
and then compute the fixed points by hand. Note that from (5.42) we have
(5.45)
$$\sigma ^{(B,A)}\left(X\right)=\underset{i}{}V_i^{\left(B\right)}XV_i^{\left(A\right)}$$
for $`X𝒦^{(B,A)}=(𝒦^{\left(A\right)},𝒦^{\left(B\right)})`$. If $`X=E_{k,l}=`$ the rank-one partial isometry mapping $`e_l`$ into $`e_k`$, we compute, in Dirac’s notation,
(5.46)
$$\sigma ^{(B,A)}\left(E_{k,l}\right)=\underset{i}{}V_i^{\left(B\right)}E_{k,l}V_i^{\left(A\right)}=\underset{i}{}|V_i^{\left(B\right)}e_kV_i^{\left(A\right)}e_l|,$$
which is complicated to handle. On the other hand,
(5.47)
$$\begin{array}{c}\sigma ^{(B,A)}\left(E_{k,l}\right)=\underset{i}{}V_i^{\left(B\right)}E_{k,l}V_i^{\left(A\right)}\hfill \\ \hfill =\underset{i}{}|V_i^{\left(B\right)}e_kV_i^{\left(A\right)}e_l|=\underset{i}{}|T_i^{\left(B\right)}e_kT_i^{\left(A\right)}e_l|,\end{array}$$
where the last step used $`T^{}`$-invariance of $`𝒦`$, see (5.6) and (5.8). But $`T^{}`$ is given by (5.35), and hence the matrix elements of $`\sigma ^{(B,A)}`$ can be computed explicitly. Let us do the calculation in a special case:
### The case $`g=2`$ and $`N>g`$
Then $`r=2`$ by Proposition 5.5, and $`𝒦`$ is three-dimensional. Then the loop $`A\left(z\right)`$ has the form
(5.48)
$$A\left(z\right)=V\left(1Q+zQ\right)=A^{\left(0\right)}+zA^{\left(1\right)},$$
where $`V`$ is a unitary and $`Q`$ a projection in $`\left(^N\right)`$. As noted in (5.25), we have
(5.49)
$$A^{\left(0\right)}A^{\left(0\right)}=\text{1}\text{1}_NQ,A^{\left(1\right)}A^{\left(1\right)}=Q.$$
Now define
(5.50)
$$\lambda _{i,j}=\left(\text{1}\text{1}Q\right)_{i,j}=\delta _{i,j}Q_{i,j},i,j=0,1,2,\mathrm{},N1.$$
Setting $`\lambda _i:=\lambda _{i,i}`$, then (5.47) and (5.35) give
(5.51)
$$\{\begin{array}{c}\sigma ^{(A,A)}\left(E_{0,0}\right)=\lambda _0E_{0,0}+\left(1\lambda _0\right)E_{1,1},\hfill \\ \sigma ^{(A,A)}\left(E_{1,1}\right)=\lambda _{N1}E_{1,1}+\left(1\lambda _{N1}\right)E_{2,2},\hfill \\ \sigma ^{(A,A)}\left(E_{2,2}\right)=\lambda _{N2}E_{1,1}+\left(1\lambda _{N2}\right)E_{2,2},\hfill \\ \sigma ^{(A,A)}\left(E_{0,1}\right)=\lambda _{0,N1}E_{0,1}\lambda _{0,N1}E_{1,2},\hfill \\ \sigma ^{(A,A)}\left(E_{0,2}\right)=\lambda _{0,N2}E_{0,1}\lambda _{0,N2}E_{1,2},\hfill \\ \sigma ^{(A,A)}\left(E_{1,2}\right)=\lambda _{N1,N2}E_{1,1}\lambda _{N1,N2}E_{2,2},\hfill \\ \sigma ^{(A,A)}\left(E_{k,l}\right)=\sigma ^{(A,A)}\left(E_{l,k}\right)^{}\text{ whenever }0l<k2.\hfill \end{array}$$
From here one computes
(5.52)
$$\begin{array}{c}det\left(x\text{1}\text{1}\sigma ^{(A,A)}\right)=det\left(x\text{1}\text{1}\sigma ^{(A,A)}\right)\hfill \\ \hfill =x^4\left(x1\right)\left(x\lambda _0\right)\left(x\left(\lambda _{N1}\lambda _{N2}\right)\right)\left(x\lambda _{0,N1}\right)\left(x\overline{\lambda _{0,N1}}\right).\end{array}$$
Since the $`\lambda _{i,j}`$ are matrix elements of a projection, we see immediately that the eigenvalue $`1`$ has multiplicity one except in special cases as when $`\lambda _0=1`$, $`\lambda _{N1}=1`$ and $`\lambda _{N2}=0`$, or $`\lambda _{0,N1}=1`$. Hence:
###### Theorem 5.7.
If $`g=2`$ and $`N>2`$, the representation $`T^{\left(A\right)}`$ is irreducible for generic loops $`A`$. It is reducible if and only if the projection $`Q\left(^N\right)`$ defining $`A`$ by (5.48) satisfies at least one of the conditions
(5.53)
$$Q_{0,0}=0(\lambda _0=1)$$
or
(5.54)
$$Q_{N1,N1}=0\text{ and }Q_{N2,N2}=1(\lambda _{N1}=1\text{ and }\lambda _{N2}=0).$$
###### Proof.
We will use the fact that the $`\lambda _{i,j}`$ numbers are the matrix entries of a projection, $`\text{1}\text{1}Q`$. Since $`\left(\text{1}\text{1}Q\right)^2=\left(\text{1}\text{1}Q\right)`$, we have
(5.55)
$$\lambda _k=\lambda _k^2+\underset{jk}{}\left|\lambda _{k,j}\right|^2.$$
Thus the case (5.53) may occur with
(5.56)
$$Q=\left(\begin{array}{cc}0& \begin{array}{cccc}0& 0& \mathrm{}& 0\end{array}\\ \begin{array}{c}0\\ 0\\ \mathrm{}\\ 0\end{array}& \text{ }\begin{array}{ccc}\mathrm{}\mathrm{}& \mathrm{}& \mathrm{}\mathrm{}\\ \mathrm{}& Q_{red}& \mathrm{}\\ \mathrm{}\mathrm{}& \mathrm{}& \mathrm{}\mathrm{}\end{array}\text{}\text{ }\end{array}\right)=\left(0\right)Q_{red}.$$
Also (5.54) may occur with
(5.57)
$$Q=\left(\begin{array}{cc}\begin{array}{ccc}\mathrm{}\mathrm{}& \mathrm{}& \mathrm{}\mathrm{}\\ \mathrm{}& Q_{red}& \mathrm{}\\ \mathrm{}\mathrm{}& \mathrm{}& \mathrm{}\mathrm{}\end{array}& \begin{array}{cc}0& 0\\ 0& 0\\ \mathrm{}& \mathrm{}\\ 0& 0\end{array}\\ \begin{array}{cccc}0& 0& \mathrm{}& 0\\ 0& 0& \mathrm{}& 0\end{array}& \begin{array}{cc}1& 0\\ 0& 0\end{array}\end{array}\right)=Q_{red}\left(\begin{array}{cc}1& 0\\ 0& 0\end{array}\right).$$
But the case $`Q_{0,N1}=1`$, i.e., $`\lambda _{0,N1}=\lambda _{N1,0}=1`$, cannot occur. To see this, recall that, by (5.25) and (5.50), the $`\lambda _{i,j}`$ numbers are the matrix entries of $`\text{1}\text{1}_NQ`$, which is a projection in $`^N`$, since $`Q`$ is. Specifically,
(5.58)
$$\lambda _{0,N1}=\epsilon _0\text{ }Q\epsilon _{N1},$$
where $`\left\{\epsilon _i\right\}_{i=0}^{N1}`$ is the canonical basis for $`^N`$. So we are considering the possibility of having an off-diagonal entry $`\lambda _{0,N1}`$ assume the value $`1`$, and that is impossible. In fact, let $`P=P^{}=P^2`$ be any projection, and let $`\epsilon `$, $`\epsilon ^{}`$ be given orthogonal unit vectors. Then
(5.59)
$$\left|\epsilon \text{ }P\epsilon ^{}\right|\frac{1}{2},$$
and so in particular the value $`1`$ is excluded. Note that (5.59) is clearly sharp: take $`\left(\begin{array}{cc}1/2& 1/2\\ 1/2& 1/2\end{array}\right)`$. To prove (5.59), note that
(5.60)
$$\left|\epsilon \text{ }P\epsilon ^{}\right|^2+\left|\epsilon ^{}\text{ }P\epsilon ^{}\right|^2P\epsilon ^{}^2,$$
by Bessel. Since $`\epsilon ^{}\text{ }P\epsilon ^{}=P\epsilon ^{}^2`$, we get
(5.61)
$$\begin{array}{c}\left|\epsilon P\epsilon ^{}\right|\sqrt{P\epsilon ^{}^2P\epsilon ^{}^4}=P\epsilon ^{}\left(1P\epsilon ^{}^2\right)^{1/2}\hfill \\ \hfill =P\epsilon ^{}\left(\text{1}\text{1}P\right)\epsilon ^{}\frac{1}{2}\left(P\epsilon ^{}^2+\left(\text{1}\text{1}P\right)\epsilon ^{}^2\right)=\frac{1}{2}\epsilon ^{}^2=\frac{1}{2}.\end{array}$$
This proves (5.59). ∎
###### Remark 5.8.
So in summary, we may define a transformation $`\sigma `$ for any system of numbers $`\lambda _{i,j}`$ via (5.51), or the matrix (5.63) below, and its spectrum will be given by (5.52). We are then interested in when points from the spectrum of $`\sigma `$ can attain the value $`1`$, and we saw that some can, and others cannot, viz., if the $`\lambda _{i,j}`$ are derived from an $`A`$ with $`g=2`$, then we get the following estimates on the points $`\lambda _0`$, $`\lambda _{N1}\lambda _{N2}`$, $`\lambda _{0,N1}`$, and $`\lambda _{N1,0}`$ in the spectrum of $`\sigma `$:
$$0\lambda _01,1\lambda _{N1}\lambda _{N2}1,$$
and, by the middle step in (5.61),
$$\left|\lambda _{0,N1}\right|=\left|\lambda _{N1,0}\right|\mathrm{min}\{\left(\lambda _0\left(1\lambda _0\right)\right)^{1/2},\left(\lambda _{N1}\left(1\lambda _{N1}\right)\right)^{1/2}\}\frac{1}{2}.$$
In particular, the spectral radius of such a $`\sigma =\sigma ^{\left(A\right)}`$ is at most $`1`$. Nonetheless, we know that $`\sigma ^{\left(A\right)}`$ is generally not contractive in the Hilbert space $`\left(𝒦^{\left(A\right)}\right)`$ of Hilbert-Schmidt operators, although $`\sigma ^{\left(A\right)}`$ obviously is contractive as a completely positive map on the $`C^{}`$-algebra $`\left(𝒦^{\left(A\right)}\right)`$.
###### Remark 5.9.
Hence for a given unitary $`V\mathrm{U}\left(N\right)`$ and projection $`Q\left(^N\right)`$, we see that the eigenvalue $`1`$ may have the following multiplicities in the characteristic polynomial for various values of $`A(V,Q)`$:
$`1`$ has multiplicity $`3`$: This occurs in one case.
###### Case 1.
$`\lambda _0=1`$, $`\lambda _{N2}=0`$, and $`\lambda _{N1}=1`$. The projection $`\text{1}\text{1}Q`$ then has the form
(5.62) 11Q=(
1
00
0000P00000000000001),11𝑄fragments
1
00
0000𝑃00000000000001\mathchoice{\hbox{\small 1\kern-3.3pt\normalsize 1}}{\hbox{\small 1\kern-3.3pt\normalsize 1}}{\hbox{\tiny 1\kern-2.3pt\scriptsize 1}}{\hbox{\Tiny 1\kern-2.0pt\tiny 1}}-Q=\left(\begin{tabular}[c]{c|c|c}$1$&$\begin{array}[c]{ccc}0&\cdots&0\end{array}$&$\begin{array}[c]{cc}0&0\end{array}$\\
\hline\cr$\begin{array}[c]{c}0\\
\vdots\\
0\end{array}$&$\begin{array}[c]{ccc}\hbox to0.0pt{\hss$\,\cdots$\hss}\hbox to0.0pt{\hss\raisebox{-2.0pt}{$\vdots$}\hss}&\cdots&\hbox to0.0pt{\hss$\,\cdots$\hss}\hbox to0.0pt{\hss\raisebox{-2.0pt}{$\vdots$}\hss}\\
\vdots&P&\vdots\\
\hbox to0.0pt{\hss$\,\cdots$\hss}\hbox to0.0pt{\hss\raisebox{-2.0pt}{$\vdots$}\hss}&\cdots&\hbox to0.0pt{\hss$\,\cdots$\hss}\hbox to0.0pt{\hss\raisebox{-2.0pt}{$\vdots$}\hss}\end{array}$&$\begin{array}[c]{cc}0&0\\
\vdots&\vdots\\
0&0\end{array}$\\
\hline\cr$\begin{array}[c]{c}0\\
0\end{array}$&$\begin{array}[c]{ccc}0&\cdots&0\\
0&\cdots&0\end{array}$&$\begin{array}[c]{cc}0&0\\
0&1\end{array}$\end{tabular}\right),
where $`P`$ is a projection in $`\left(^{N3}\right)`$ (so when $`N=3`$ there is no more choice.)
$`1`$ has multiplicity $`2`$: This occurs in two mutually exclusive cases.
###### Case 2.
$`\lambda _0=1`$ and ($`\lambda _{N2}>0`$ or $`\lambda _{N1}<1`$).
###### Case 3.
$`\lambda _0<1`$, $`\lambda _{N2}=0`$, and $`\lambda _{N1}=1`$.
$`1`$ has multiplicity $`1`$: This occurs in all remaining cases.
Since the dimension of the fixed-point set of $`\sigma ^{\left(A\right)}`$ is always at most equal to the multiplicity of $`1`$ in the characteristic polynomial, it follows from Remark 5.9 and Theorem 5.1 that the linear dimension of the commutant of the $`T^{\left(A\right)}`$($`=T^{(V,Q)}`$)-representation is always at most $`3`$. Since the smallest-dimensional nonabelian $`C^{}`$-algebra is $`\left(^2\right)`$, which has dimension $`4`$, it follows that the commutant is always abelian when $`g=2`$, and the representation decomposes into a sum of at most three irreducible mutually disjoint representations. In order to get the exact number, we have to compute the fixed-point space of $`\sigma ^{\left(A\right)}`$ exactly. To this end, we note that the matrix of $`\sigma ^{\left(A\right)}`$, relative to the basis $`\left\{E_{i,k};2i,k0\right\}`$ of $`\left(^3\right)`$ is (we actually compute the matrix of $`\sigma ^{\left(A\right)}`$ using (5.51) and take the adjoint)
(5.63)
| $`\sigma ^{\left(A\right)}`$ | $`(0,0)`$ | $`(1,1)`$ | $`(2,2)`$ | $`(0,1)`$ | $`(0,2)`$ | $`(1,2)`$ | $`(2,1)`$ | $`(1,0)`$ | $`(2,0)`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`(0,0)`$ | $`\lambda _0`$ | $`1\lambda _0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(1,1)`$ | $`0`$ | $`\lambda _{N1}`$ | $`1\lambda _{N1}`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(2,2)`$ | $`0`$ | $`\lambda _{N2}`$ | $`1\lambda _{N2}`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(0,1)`$ | $`0`$ | $`0`$ | $`0`$ | $`\lambda _{N1,0}`$ $`0`$ $`\lambda _{N2,0}`$ $`0`$ $`\lambda _{N1,0}`$ | | $`\lambda _{N1,0}`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(0,2)`$ | $`0`$ | $`0`$ | $`0`$ | $`\lambda _{N2,0}`$ | | $`\lambda _{N2,0}`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(1,2)`$ | $`0`$ | $`\lambda _{N2,N1}`$ | $`\lambda _{N2,N1}`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(2,1)`$ | $`0`$ | $`\lambda _{N1,N2}`$ | $`\lambda _{N1,N2}`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ |
| $`(1,0)`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`\lambda _{0,N1}`$ | $`\lambda _{0,N1}`$ | $`0`$ |
| $`(2,0)`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`0`$ | $`\lambda _{0,N2}`$ | $`\lambda _{0,N2}`$ | $`0`$ |
Thus one may compute the eigenvectors of $`\sigma ^{\left(A\right)}`$. The generic result is as in Table 1.
Note that in non-generic cases, like $`\lambda _0=1`$, and/or ($`\lambda _{N1}=1`$ and $`\lambda _{N2}=0`$), some of the listed eigenvectors are zero. This is not surprising since these are exactly the cases where the root $`1`$ in the characteristic polynomial is multiple. Let us consider these cases separately:
###### Case 2.
$`\lambda _0=1`$ ($`\lambda _{0,N1}=\lambda _{0,N2}=0`$) and $`\lambda _{N1}<1`$ or $`\lambda _{N2}>0`$: In this case the last two eigenvectors in the list in Table 1 are zero, but one may verify that these may be replaced by the vectors $`E_{0,2}`$ and $`E_{2,0}`$. Thus the eigenspace corresponding to $`1`$ is indeed $`2`$-dimensional in this case, and spanned by $`\text{1}\text{1}_𝒦`$ and $`E_{0,0}`$.
###### Case 3.
$`\lambda _0<1`$, $`\lambda _{N1}=1`$, $`\lambda _{N2}=0`$: In this case it follows from (5.55) that $`\lambda _{N1,k}=0`$ for all $`kN1`$ and $`\lambda _{N2,k}=0`$ for all $`kN2`$, so the matrix in (5.63) degenerates into the upper-left-hand $`3\times 3`$ matrix, which is
(5.64)
$$\left(\begin{array}{ccc}\lambda _0& 1\lambda _0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).$$
So we see that a new eigenvector for $`\lambda _{N1}\lambda _{N2}=1`$ is
(5.65)
$$E_{2,2}.$$
Let us now look at the even more degenerate case where the multiplicity of $`1`$ in the characteristic polynomial is $`3`$:
###### Case 1.
$`\lambda _0=1`$, $`\lambda _{N1}=1`$, $`\lambda _{N2}=0`$: From the discussion of Case 3 above, it follows that then the upper-left $`3\times 3`$ matrix of (5.63) is $`\text{1}\text{1}_3`$, while the rest of the matrix is zero. Thus the fixed-point set is the linear span of
(5.66)
$$E_{0,0},E_{1,1}\text{, and }E_{2,2}.$$
Thus the dimension of the eigensubspace is in all cases where the parameters are restricted as in Remark 5.8 equal to the multiplicity of the root $`1`$ in the characteristic polynomial, and the above analysis implies
###### Corollary 5.10.
If $`g=2`$ and $`N>2`$, the representation $`T^{\left(A\right)}`$ decomposes into at most three irreducible representations which are mutually nonequivalent. Generically, $`T^{\left(A\right)}`$ is irreducible, and otherwise the decomposition structure is summarized in Table 2.
###### Remark 5.11.
Table 2 shows a remarkable feature which these representations share with the $`g=2`$, $`N=2`$ representations studied in \[BEJ00\], namely that the fixed-point set of $`\sigma ^{\left(A\right)}`$ is both abelian and an algebra in all cases. Hence it follows from (5.13) that the projection $`P`$ onto $`𝒦`$ commutes with the commutant $`\pi \left(𝒪_N\right)^{}`$, and thus we may find cyclic vectors for the various subrepresentations by picking vectors in the range of the minimal projections in the fixed-point set $`\left(𝒦\right)^{\sigma ^{\left(A\right)}}`$. For example, for the first case in Table 2 (Case 1), the three vectors $`z^n`$, $`n=0,1,2`$, in $`L^2\left(𝕋\right)`$ are cyclic for the three disjoint representations the original representation decomposes into, respectively. For the second case in Table 2 (Case 3), the pair $`z^0`$, $`z^2`$ has the same property, as well as the pair $`z^1`$, $`z^2`$ or any pair of the form $`\lambda z^0+\mu z^1`$, $`z^2`$. So far we do not know a single example (for general $`N`$, $`g`$) where $`\left(𝒦\right)^{\sigma ^{\left(A\right)}}`$ is not abelian, or is not an algebra.
## 6. Reduction of representations as a reduction of $`N`$
Consider an arbitrary element $`A𝒫(𝕋,\mathrm{U}(N))`$. Hence both $`N`$ and the genus $`g`$ are arbitrary, but given. We saw that there is an associated wavelet filter $`m^{(A)}`$, as well as a representation $`T^{(A)}`$ of $`𝒪_N`$ acting on $`L^2(𝕋)`$. (As before we work in the standard Fourier basis for $`L^2(𝕋)`$, i.e., $`e_n(z)=z^n`$, $`n`$.) We identified the subspace $`𝒦`$ spanned by $`\{e_r,\mathrm{},e_1,e_0\}`$ where $`r`$ was picked so that reducibility of $`T^{(A)}`$ is decided by the fixed-point set of an associated completely positive map $`\sigma _𝒦^{(A)}`$ given by
(6.1)
$$\sigma _𝒦^{(A)}()=\underset{i=0}{\overset{N1}{}}V_i^{(A)}()V_i^{(A)}$$
with
(6.2)
$$V_i^{(A)}=𝒫_𝒦T_i^{(A)}\text{.}$$
While we identified the complete spectral picture of $`\sigma _𝒦^{(A)}`$ in the special case $`g=2`$, the higher-genus case is not yet entirely understood. But we note that the identity from the $`g=2`$ case,
(6.3)
$$\sigma _𝒦^{(A)}(E_{0,0})=\lambda _0(A)E_{0,0},$$
which is clear from (5.63), remains true for arbitrary $`g`$. Hence our use of $`\sigma _𝒦^{(A)}`$ yields reducibility of the representation $`T^{(A)}`$ whenever $`\lambda _0(A)=1`$. We recall that if
(6.4)
$$A(z)=A^{(0)}+zA^{(1)}+\mathrm{}+z^{g1}A^{(g1)},$$
then $`\lambda _0(A)`$ is defined as the $`(0,0)`$-matrix entry in $`A^{(0)}A^{(0)}`$.
###### Remark 6.1.
To prove (6.3) for the general case, recall that
(6.5)
$$T_i^{\left(A\right)}e_0=\underset{j=0}{\overset{N1}{}}A_{i,j}\left(z^N\right)e_j=\underset{j=0}{\overset{N1}{}}\underset{k=0}{\overset{g1}{}}A_{i,j}^{\left(k\right)}e_{j+kN}.$$
Using (6.2), we therefore get
(6.6)
$$V_i^{\left(A\right)}e_0=P_𝒦T_i^{\left(A\right)}e_0=\underset{j=0}{\overset{N1}{}}\underset{k=0}{\overset{g1}{}}A_{i,j}^{\left(k\right)}P_𝒦e_{j+kN}=A_{i,0}^{\left(0\right)}e_0,$$
and by (5.46),
(6.7) $`\sigma _𝒦^{\left(A\right)}\left(E_{0,0}\right)`$ $`=\left({\displaystyle \underset{i=0}{\overset{N1}{}}}|A_{i,0}^{\left(0\right)}e_0A_{i,0}^{\left(0\right)}e_0|\right)=\left({\displaystyle \underset{i=0}{\overset{N1}{}}}\left|A_{i,0}^{\left(0\right)}\right|^2\right)|e_0e_0|`$
$`=\left(A^{\left(0\right)}A^{\left(0\right)}\right)_{0,0}E_{0,0}=\lambda _0\left(A\right)E_{0,0},`$
which is (6.3).
But it is worth stressing that the argument (6.6), for the vector $`e_0`$ in $`𝒦`$, does *not* carry over to the other basis vectors $`e_1`$, $`e_2`$, $`\mathrm{}`$. This is clear already in the case $`g=2`$ from (5.63). More generally, for $`e_1`$, for example,
(6.8)
$$V_i^{\left(A\right)}e_1=A_{i,0}^{\left(1\right)}e_0+A_{i,N1}^{\left(0\right)}e_1+A_{i,N2}^{\left(0\right)}e_2+\mathrm{},$$
and therefore
(6.9)
$$\begin{array}{c}\sigma _𝒦^{\left(A\right)}\left(E_{1,1}\right)=\lambda _{0,0}^{\left(11\right)}E_{0,0}+\lambda _{N1,N1}^{\left(00\right)}E_{1,1}+\lambda _{N2,N2}^{\left(00\right)}E_{2,2}+\mathrm{}\hfill \\ \hfill +\lambda _{N1,0}^{\left(01\right)}E_{0,1}+\lambda _{0,N1}^{\left(10\right)}E_{1,0}+\lambda _{N2,N1}^{\left(00\right)}E_{1,2}+\lambda _{N1,N2}^{\left(00\right)}E_{2,1}+\mathrm{},\end{array}$$
where $`\lambda _{i,j}^{\left(kl\right)}:=\left(A^{\left(k\right)}A^{\left(l\right)}\right)_{i,j}`$, i.e., the $`(i,j)`$ entry in the $`N\times N`$ matrix $`A^{\left(k\right)}A^{\left(l\right)}`$. Even if $`\lambda _{N1,N1}^{\left(00\right)}=1`$, that does not imply that all the other coefficients $`\lambda _{i,j}^{\left(kl\right)}`$ from (6.9) will necessarily vanish; see the details below. While
(6.10)
$$\lambda _{N1,N1}^{\left(00\right)}=1\lambda _{N1,j}^{\left(00\right)}=\lambda _{i,N1}^{\left(00\right)}=0\text{ for all }i,j\{0,\mathrm{},N2\},$$
the other coefficients in (6.9) such as $`\lambda _{N2,N2}^{\left(00\right)}`$ or $`\lambda _{0,0}^{\left(11\right)}`$ would typically be nonzero even if $`\lambda _{N1,N1}^{\left(00\right)}=1`$. So even then, $`E_{1,1}`$ will not be fixed by $`\sigma _𝒦^{\left(A\right)}`$, and there is not a natural analogue to (6.3). While if $`g=2`$, then the difference $`\lambda _{N1,N1}^{\left(00\right)}\lambda _{N2,N2}^{\left(00\right)}`$ is in the spectrum of $`\sigma _𝒦^{\left(A\right)}`$, see Table 1, this will generally not be true if $`g>2`$. Then the spectral picture for $`\sigma _𝒦^{\left(A\right)}`$ is not fully understood.
However, the following result shows that a general wavelet representation $`T^{(A)}`$ may be reducible because a special point $`\lambda _0(A)`$ in the spectrum of $`\sigma _𝒦^{(A)}`$ may be one. This is a special reduction, however, as we also showed that whenever any point $`\lambda `$ in spectrum $`\left(\sigma _𝒦^{(A)}\right)`$ satisfies $`\lambda =1`$, we get a possibly different reduction of the representation $`(T^{(A)},L^2(𝕋))`$.
###### Theorem 6.2.
Let $`A𝒫(𝕋,\mathrm{U}(N))`$. Then the following five conditions are equivalent.
1. $`\lambda _0(A)=1`$.
2. There is a $`B=(B_{i,j})_{i,j=1}^{N1}𝒫(𝕋,\mathrm{U}(N1))`$ such that
(6.11) A(1)1A(z)=(
1
0
0
00B1,1(z)B1,N1(z)00BN1,1(z)BN1,N1(z)),𝐴superscript11𝐴𝑧fragments
1
0
0
00fragmentsB11(z)fragmentsB1𝑁1(z)0missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpression0fragmentsB𝑁11(z)fragmentsB𝑁1𝑁1(z)A(1)^{-1}A(z)=\left(\begin{tabular}[c]{c|cclc}$1$&$0$&$0$&$\cdots$&$0$\\
\hline\cr$0$&$B_{1,1}(z)$&$\cdots$&$\cdots$&$B_{1,N-1}(z)$\\
$0$&$\vdots$&&&$\vdots$\\
$\vdots$&$\vdots$&&&$\vdots$\\
$0$&$B_{N-1,1}(z)$&$\cdots$&$\cdots$&$B_{N-1,N-1}(z)$\end{tabular}\right),
i.e.,
(6.12) A(1)1A(z)=(
1
0
00B(z)0)=(1)B(z).𝐴superscript11𝐴𝑧fragments
1
0
00missing-subexpressionmissing-subexpressionmissing-subexpressionmissing-subexpressionfragmentsB(z)missing-subexpression0missing-subexpressionmissing-subexpressionmissing-subexpressiondirect-sum1𝐵𝑧A(1)^{-1}A(z)=\left(\begin{tabular}[c]{c|lll}$1$&$0$&$\cdots$&$0$\\
\hline\cr$0$&&&\\
$\vdots$&&$B(z)$&\\
$0$&&&\end{tabular}\right)=\left(1\right)\oplus B\left(z\right).
3. After modifying with a $`\mathrm{U}(N)`$-automorphism of $`𝒪_N`$, the correspondingwavelet filter $`m^{(A)}`$ satisfies
(6.13)
$$\{\begin{array}{c}m_0^{(A)}(z)1\text{, for all }z𝕋,\hfill \\ m_i^{(A)}(z)=_{j=1}^{N1}B_{i,j}(z^N)z^j.\hfill \end{array}$$
4. $$T_i^{(A)^{}}e_0e_0\text{, for all }i=0,\mathrm{},N1\text{.}$$
5. $$\sigma _𝒦^{\left(A\right)}\left(E_{0,0}\right)E_{0,0}.$$
###### Remark 6.3.
Formula (6.13) explains the assertion made in the Introduction about reduction in the size of the number $`N`$ which serves as the scaling of the wavelet in question. Once $`B`$ is identified then there is an $`(N1)`$-wavelet filter
(6.14)
$$m_{i1}^{(B)}(z)=\underset{j=0}{\overset{N2}{}}B_{i,j+1}\left(z^{N1}\right)z^j,$$
$`i=1,\mathrm{},N1`$.
###### Proof.
(i) $``$ (ii): Using the factorization (5.23) in Corollary 5.2, we note that
(6.15)
$$\begin{array}{c}A^{(0)}A^{(0)}=\left(\text{1}\text{1}Q_{g1}\right)\left(\text{1}\text{1}Q_{g2}\right)\mathrm{}\hfill \\ \hfill \mathrm{}\left(\text{1}\text{1}Q_1\right)\overline{)\left(\text{1}\text{1}Q_0\right)}\left(\text{1}\text{1}Q_1\right)\mathrm{}\\ \hfill \mathrm{}\left(\text{1}\text{1}Q_{g2}\right)\left(\text{1}\text{1}Q_{g1}\right)\end{array}$$
and the $`V`$ ($`\mathrm{U}(N)`$) from (5.23) is $`V=A(1)`$, i.e., evaluation of $`A(z)`$ at $`z=1`$. Let $`\epsilon _0`$ be the first canonical basis vector in $`^N`$. Then (i) states that
(6.16)
$$\left(\text{1}\text{1}Q_0\right)\left(\text{1}\text{1}Q_1\right)\mathrm{}\left(\text{1}\text{1}Q_{g1}\right)\epsilon _0=1.$$
Hence $`\left(\text{1}\text{1}Q_{g1}\right)\epsilon _0H\left(\text{1}\text{1}Q_{g1}\right)\epsilon _0=1`$, where
(6.17)
$$H:=\left(\text{1}\text{1}Q_{g2}\right)\mathrm{}\left(\text{1}\text{1}Q_1\right)\overline{)\left(\text{1}\text{1}Q_0\right)}\left(\text{1}\text{1}Q_1\right)\mathrm{}\left(\text{1}\text{1}Q_{g2}\right).$$
Using Schwarz’s inequality and induction, we conclude that $`\left(\text{1}\text{1}_NQ_j\right)\epsilon _0=\epsilon _0`$, and therefore $`Q_j\epsilon _0=0`$ for all $`j=0,\mathrm{},g1`$. Hence $`\left(\text{1}\text{1}_NQ_j+zQ_j\right)\epsilon _0=\epsilon _0`$, and from (5.23),
(6.18)
$$A(1)^1A(z)\epsilon _0=\epsilon _0\text{.}$$
The same argument yields
(6.19)
$$\epsilon _0\text{ }A(1)^1A(z)e_j=\delta _{0,j},$$
and so (ii) follows.
(ii) $``$ (iii): Substituting (6.11) into (2.24), we get
(6.20)
$$\left(V^1A\right)_{i,0}(z)=\delta _{i,0},$$
(6.21)
$$\left(V^1A\right)_{0,j}(z)=\delta _{j,0},$$
and
(6.22)
$$\left(V^1A\right)_{i,j}(z)=B_{i,j}(z)\text{ if }i,j1,$$
and (6.13) in (iii) follows.
(iii) $``$ (iv): Using
(6.23)
$$T_i^{(A)}\xi (z)=\frac{1}{N}\underset{w^N=z}{}\overline{m_i^{(A)}(w)}\xi (w)\text{}\xi L^2(𝕋),$$
and introducing the matrix product from (iii), we get
(6.24)
$$\{\begin{array}{c}A_{i,0}(z)=V_{i,0},\hfill \\ A_{i,j}(z)=_{k=1}^{n1}V_{i,k}B_{i,k}(z)\text{}j>0.\hfill \end{array}$$
Using again (2.24), we thus get
(6.25)
$$T_i^{(A)}e_0=\overline{V_{i,0}}e_0,$$
where $`V_{i,j}`$ denotes the matrix entries of $`V`$ ($`:=A(1)`$) in $`\mathrm{U}(N)`$.
(iv) $``$ (v) $``$ (i): By (6.23),
(6.26)
$$T_i^{(A)}e_0=\overline{A_{i,0}(z)}𝒫(𝕋,)\text{;}$$
and so if (iv) is assumed, then
(6.27)
$$A_{i,0}^{(k)}=0\text{ for }k>0.$$
Since
(6.28)
$$\underset{i}{}\overline{A_{i,0}(z)}A_{i,0}(z)=1$$
in general, we get
(6.29)
$$\lambda _0(A)=\underset{i}{}\overline{A_{i,0}^{(0)}}A_{i,0}^{(0)}=1.$$
From (6.26), we get
(6.30)
$$T_i^{\left(A\right)}e_0=\underset{k=0}{\overset{g1}{}}\overline{A_{i,0}^{\left(k\right)}}e_k,$$
and by (5.47),
(6.31)
$$\sigma _𝒦^{\left(A\right)}\left(E_{0,0}\right)=\underset{k=0}{\overset{g1}{}}\underset{l=0}{\overset{g1}{}}\lambda _{0,0}^{\left(kl\right)}E_{k,l}.$$
Note the coefficient of $`E_{0,0}`$ in (6.31) is $`\lambda _0\left(A\right)=\lambda _{0,0}^{\left(00\right)}`$. For the other coefficients, we have
(6.32)
$$\left|\lambda _{0,0}^{\left(kl\right)}\right|^2\lambda _{0,0}^{\left(kk\right)}\lambda _{0,0}^{\left(ll\right)}$$
and
(6.33)
$$\lambda _{0,0}^{\left(kk\right)}=\underset{i=0}{\overset{N1}{}}\left|A_{i,0}^{\left(k\right)}\right|^2.$$
Hence, the vanishing of the non-$`(0,0)`$ coefficients in (6.31) is equivalent to condition (6.27), which we already showed is equivalent to both (iv) and (i). ∎
For $`\xi 𝒦`$, let $`_+^{\left(A\right)}\left[\xi \right]`$ denote the cyclic subspace in $`L^2\left(𝕋\right)`$ generated by $`\xi `$ and the operators $`T_i^{\left(A\right)}`$, and let $`E_+^{\left(A\right)}\left[\xi \right]`$ be the corresponding projection.
###### Corollary 6.4.
Suppose condition (i) of Theorem 6.2 holds. Then the (unique) operator $`W`$ in the commutant of $`T^{\left(A\right)}`$ which satisfies
(6.34)
$$P_𝒦WP_𝒦=E_{0,0}$$
is $`W=E_+^{\left(A\right)}\left[e_0\right]`$.
###### Proof.
Let $`v^N`$ be the vector with coordinates $`v_i:=A_{i,0}^{\left(0\right)}`$. Then (i) holds if and only if $`v=1`$. In that case, the restriction of $`T^{\left(A\right)}`$ to $`_+^{\left(A\right)}\left[e_0\right]`$ is the subrepresentation which is induced by the Cuntz state \[Cun77\] $`\omega _v`$. Specifically, if $`I=(i_1,\mathrm{},i_k)`$, $`J=(j_1,\mathrm{},j_l)`$ are multi-indices, then it follows from (iv) that
(6.35)
$$e_0\text{ }T_I^{\left(A\right)}T_J^{\left(A\right)}e_0=v_{i_1}\mathrm{}v_{i_k}\overline{v_{j_1}}\mathrm{}\overline{v_{j_l}}=\omega _v\left(s_Is_J^{}\right).$$
Hence $`_+^{\left(A\right)}\left[e_0\right]`$ is the closed subspace spanned by $`e_0`$ and $`\left\{T_I^{\left(A\right)}e_0\right\}`$. Since
(6.36)
$$T_i^{\left(A\right)}e_0\left(z\right)=\underset{j=0}{\overset{N1}{}}A_{i,j}\left(z^N\right)z^j𝒫(𝕋,),$$
$`_+^{\left(A\right)}\left[e_0\right]`$ is contained in the Hardy space $`_+:=\overline{span}\left\{e_n;n0\right\}`$. Since $`𝒦=span\{e_r,\mathrm{},e_1,e_0\}`$, it follows that $`E_+^{\left(A\right)}\left[e_0\right]`$ satisfies (6.34). Since $`\sigma ^{\left(A\right)}\left(E_{0,0}\right)=E_{0,0}`$, the operator $`W`$ in the commutant of $`T^{\left(A\right)}`$ satisfying (6.34) is unique by Theorem 5.1, and so $`W=E_+^{\left(A\right)}\left[e_0\right]`$. ∎
###### Remark 6.5.
The Hilbert space prior to Corollary 6.4 is of the form $`_+^{\left(A\right)}\left[\xi \right]`$ for some $`\xi 𝒦`$ where
(6.37)
$$T_i^{\left(A\right)}\xi =v_i\xi $$
for some $`v=\left(v_i\right)^N`$. Since $`\omega _v`$ is a Cuntz state, the corresponding representation is unique from $`v`$ up to unitary equivalence, by \[Cun77\]. It also follows from \[Jor99, Theorem 6.3\] that every $`\xi 𝒦`$ which satisfies (6.37) for some $`v^N`$ must be a monomial, i.e., $`\xi \left(z\right)=z^k`$ for some $`k\{0,1,\mathrm{},r\}`$. In our special setting, this can be proved as follows: If $`T_i^{}\xi =\overline{v}_i\xi `$ for a unit vector $`\xi 𝒦`$, where $`_i\left|v_i\right|^2=1`$, then $`\xi =_iT_iT_i^{}\xi =_i\overline{v}_iT_i\xi `$, or, spelled out, $`\xi \left(z\right)=_i\overline{v}_im_i\left(z\right)\xi \left(z^N\right)`$. Hence, putting $`m\left(z\right)=_i\overline{v}_im_i\left(z\right)`$, we have $`\xi \left(z\right)=m\left(z\right)\xi \left(z^N\right)`$. Now we may use the argument from \[Jor99, page 104\] or \[BrJo97b, Theorem 3.1\] to conclude that $`\left|m\left(z\right)\right|=1`$ and $`\left|\xi \left(z\right)\right|=1`$ for almost all $`z𝕋`$. But both $`m\left(z\right)`$ and $`\overline{\xi \left(z\right)}`$ are polynomials, and we deduce from Lemma 3.1 that they are monomials. If $`\xi \left(z\right)=z^k`$ it follows that $`m\left(z\right)=\xi \left(z\right)\overline{\xi \left(z^N\right)}=z^{\left(N1\right)k}`$. Thus a very restrictive necessary and sufficient condition for the Cuntz state $`\omega _v`$ to occur is that
(6.38)
$$\underset{i}{}\overline{v}_im_i\left(z\right)=z^{\left(N1\right)k}\text{ for some }k\{0,1,\mathrm{},r\}.$$
The other details of the decompositions may be spelled out as follows: Let $`=L^2\left(𝕋\right)`$. Note first that Corollary 6.4 holds whenever $`E_{0,0}`$ is replaced by any one-dimensional projection in the fixed-point set of $`\sigma `$ in the generality of Theorem 5.1, since $`E_{0,0}\left(𝒦\right)=E_{0,0}\left(\right)`$ is then a one-dimensional $`T^{\left(A\right)}`$-invariant space. If $`e`$ is a general (not necessarily one-dimensional) projection in the fixed-point set, it follows from Lemma 3.3 in \[BJKW00\] that $`e`$ commutes with all the operators $`PT_i^{\left(A\right)}P`$ and $`PT_i^{\left(A\right)}P`$ and hence the projection $`E`$ onto $`\left[𝒪_Ne𝒦\right]`$ in $``$ is in the commutant of the representation, and $`PEP=e`$. This justifies claims in Remark 5.11. In the case $`g=2`$, which is completely enumerated in \[BEJ00\] (for $`N=2`$) and in Table 2 (for $`N>2`$), the ranges of the projections in the fixed-point set are spanned by subsets of the orthonormal set $`\{e_0,e_1,\mathrm{},e_r\}`$, and hence each Hilbert space in the decomposition has the form $`E_+^{\left(A\right)}\left[\xi \right]`$ for a suitable $`\xi `$ in this set. The vector $`\xi `$ defines a Cuntz state if and only if the corresponding projection in $`\left(𝒦\right)^\sigma `$ is one-dimensional. Note also that the result \[Jor99, Theorem 6.3\] spelled out above implies that all one-dimensional projections in $`\left(𝒦\right)^\sigma `$ are diagonal in the standard basis. This gives a partial explanation of these results.
### Other examples of decompositions of $`L^2\left(𝕋\right)`$
Generically, the wavelet examples do not have minimal projections $`e`$ in $`\left(𝒦\right)^\sigma `$ of dimension one, and even if there is one such projection in $`\left(𝒦\right)^\sigma `$, there may be others that do not have dimension one. For illustration, let us repeat some of the examples in \[BEJ00\], rephrased in the setting of the present paper. Take, for example, $`N=2=g`$. Then, by (6.38), only the two cases $`\left(\begin{array}{cc}1& 0\\ 0& z\end{array}\right)`$ and $`\left(\begin{array}{cc}z& 0\\ 0& 1\end{array}\right)`$ have the Cuntz-state vectors. For $`A\left(z\right)=\left(\begin{array}{cc}1& 0\\ 0& z\end{array}\right)`$, even though $`\lambda _0\left(A\right)=1`$, this example has a minimal $`2`$-dimensional projection $`e`$ in $`\left(𝒦\right)^\sigma `$, namely $`e=E_{1,1}+E_{2,2}`$. It is minimal in the sense that the restriction of $`T^{\left(A\right)}`$ to $`E_e=\left[𝒪_2e𝒦\right]`$ is irreducible; see \[BEJ00, (4.52)\]. The matrix $`A\left(z\right)=\left(\begin{array}{cc}0& 1\\ z& 0\end{array}\right)`$ has $`\lambda _0\left(A\right)=0`$, and its representation $`T^{\left(A\right)}`$ decomposes into a sum of two irreducibles associated with respective $`2`$-dimensional projections $`e`$ and $`f`$ in $`\left(𝒦\right)^\sigma `$: $`e=E_{0,0}+E_{1,1}`$ and $`f=E_{2,2}+E_{3,3}`$; see \[BEJ00, (4.46)\]. Hence in this case, $`=L^2\left(𝕋\right)=\left[𝒪_2e𝒦\right]\left[𝒪_2f𝒦\right]`$ with $`T^{\left(A\right)}`$ restricting to an irreducible representation on each of the two subspaces. A direct application of (5.15)–(5.16) in Theorem 5.1 shows that these two subrepresentations are disjoint, i.e., inequivalent. While these examples have $`g=2`$, the cases $`g>2`$ entail a richer decomposition structure.
Note that the result, Theorem 5.1, which is used in analyzing the representations $`T^{\left(A\right)}`$, applies also when $`g>2`$. Moreover, the factorization result, Corollary 5.2, yields in principle a way of understanding the general case, i.e., arbitrary $`g`$ and $`N`$, but unpublished experimentation with examples for $`g=3`$ (i.e., two projections in $`^N`$) has not so far yielded decomposition structures more general than the above mentioned ones.
###### Remark 6.6.
The proof of (6.3) shows more generally that
(6.39)
$$\sigma ^{(B,A)}\left(E_{0,0}\right)=\left(\underset{i=0}{\overset{N1}{}}\overline{A_{i,0}^{\left(0\right)}}B_{i,0}^{\left(0\right)}\right)E_{0,0}.$$
Referring to Theorem 5.1, we note then that the necessary and sufficient condition for $`E_{0,0}`$ to induce an operator in $`L^2\left(𝕋\right)`$ which intertwines the two representations $`T^{\left(A\right)}`$ and $`T^{\left(B\right)}`$ is $`\left(A^{\left(0\right)}B^{\left(0\right)}\right)_{0,0}=1`$. Moreover, if this holds, then both $`T^{\left(A\right)}`$ and $`T^{\left(B\right)}`$ must satisfy the equivalent conditions in Theorem 6.2, and we must then further have $`A_{i,0}^{\left(0\right)}=B_{i,0}^{\left(0\right)}`$ for all $`i`$. Then the intertwining operator $`W`$ which is induced by $`E_{0,0}`$ via $`\sigma ^{(B,A)}\left(E_{0,0}\right)=E_{0,0}`$ is the one which results from the uniqueness of the $`\omega _v`$-representation where $`v_i=A_{i,0}^{\left(0\right)}=B_{i,0}^{\left(0\right)}`$.
###### Acknowledgements.
This research was done when one of us (P.E.T.J.) visited the University of Oslo with support from the university and NFR. He is grateful for kind hospitality. We wish to thank Brian Treadway for his skillful typesetting, and Rune Kleveland and Brian Treadway for suggestions and elimination of mistakes in earlier versions of this paper. We are also indebted to Rong-Qing Jia and the referee for bringing the references \[JiMi91\], \[JiSh94\], \[LLS96\], \[RiSh91\], \[RiSh92\] to our attention, and thus enabling us to complete the discussion between (2.30) and (2.31).
This is an expanded version of the invited lecture delivered by P.E.T.J. at the International Conference on Wavelet Analysis and Its Applications held at Zhongshan University, Guangzhou, China, in November of 1999 (the conference announcement may be seen in The Wavelet Digest 8 (1999), on the World Wide Web at http://www.wavelet.org/wavelet/digest\_08/digest\_08.04.html#13). |
warning/0001/cond-mat0001212.html | ar5iv | text | # High – pressure Raman study of L-alanine crystal
## Abstract
Pressure-dependent Raman scattering studies in the range 0.0 – 32 kbar were carried out in L-alanine in order to investigate its external mode phonon spectra in relation to the phase transitions in the crystal. A careful analysis of the spectra shows that the low-energy Raman modes exhibit variation both in frequency and in intensity and between 26 and 28 kbar it is observed a splitting of a external mode, indicating that the D<sub>2</sub> normal phase undergoes a transition. Pressure coefficients for external modes are also given.
I. INTRODUCTION
In recent times one has witnessed a growing interest in the study of amino acid crystals. This interest has been stimulated by the perspective of understanding a system where the hydrogen bond plays a fundamental role and, as a consequence of this understanding, a better knowledge of some important biological molecules, e.g. proteins, can be obtained. L-alanine, NH<sub>2</sub>CHCH<sub>3</sub>COOH, is one of 20 amino acid serving as building blocks for the proteins of living beings. In the crystalline form L-alanine crystal has been subjected to extensive investigations by many authors using various experimental techniques. X-ray diffraction studies revealed the crystal structure and the position of the heavy atoms of the amino acid. Neutron diffraction studies confirmed the results of Ref. 1 and determined the hydrogen atoms positions. A number of studies have been carried out on the vibrational spectra of L-alanine crystal focusing on the assignment of modes and the effects of various intermolecular potentials on the crystal vibration , the effect of temperature on the crystal structure , the NH<sub>3</sub>-torsion temperature dependence and the dynamic localization of vibrational energy . Some studies have been made on the vibrational spectra of L-alanine crystal. Thermal conductivity and phonon echo in L-alanine were also performed. NMR experiments and spin-lattice relaxation time T<sub>1</sub> have suggested that the crystal undergoes a phase transition at 178 K. However, Raman spectra do not confirm the existence of this transition. It was also suggested that L-alanine crystals should undergo a phase transition induced by the application of pressure, but no experimental investigation was made up to now.
It is well established that the use of hydrostatic pressure can diminish the interatomic and the intermolecular distances of ions and molecules in the crystal structure producing, eventually, a change of structure. Such effects were detected recently by Raman scattering experiments performed on monohydrated L-asparagine crystal, (and confirmed by X-ray diffraction synchrotron measurements.), where three different phase transitions were observed for pressures up to 2.0 GPa.
In this letter we investigate the influence of high hydrostatic pressure on a L-alanine crystal observing the low-frequency region of the Raman spectra. By exploring the fact that the modes of the low-frequency region gives important information about the structure’s stability we performed Raman measurements in the range of pressure 0.0 – 32 kbar. We also give pressure coefficients for the external modes, fundamental data to determine Grüneisen parameter.
II. EXPERIMENTAL
High-pressure Raman experiments at room temperature were performed on a small piece of L-alanine sample pressed in a standard Diamond Anvil Cell. As hydrostatic medium we have used a 4:1 methanol:ethanol mixture. The pressure was calibrated using the shifts of ruby lines. The pressure calibration is expected to be accurate by $`\pm `$ 0.3 kbar. The Raman spectra were excited with a 514.5 nm line of an Argon ion laser working at 30 mW in the backscattering geometry. The laser beam was focused on the sample surface using a lens with f = 20.5 mm. To ensure focusing of laser on the sample when Raman spectra were recorded or on the ruby chip when pressure was calibrated a image of the hole in the gasket, the pressure compartment, was recorded by a CCD camera. The backscattering light was analyzed in a triplemate Jobin Yvon spectrometer (T64000) equipped with a N<sub>2</sub> \- cooled CCD system. The frequency of the Raman bands are expected to be accurate by $`\pm `$ 2 cm<sup>-1</sup>.
III. RESULTS AND DISCUSSIONS
The L-alanine crystal has an orthorhombic symmetry with four molecules per unit cell with a space group P2<sub>1</sub>2<sub>1</sub>2<sub>1</sub> (D$`{}_{}{}^{4}{}_{2}{}^{}`$) and cell parameters a = 6.023, b = 12.343 and c = 5.784 $`\AA `$. The molecule is found in the zwitterion form and all three hydrogen atoms of the ammonium group form hydrogen bonds with three different neighboring molecules. In fact, the main link of molecules in the crystal structure occurs along the c-axis, where it is observed a chain of hydrogen bonds. Also, two other hydrogen bonds bind these chains together producing a three dimensional network. Such link mechanism is common to most amino acid crystals as L-threonine and monohydrated L-asparagine .
It is well established that the intramolecular modes are in the high frequency region of the Raman spectra of molecular crystals while the external phonon modes, representing essentially the translations and the rotations of rigid molecules, are in the low frequency region. The pressure effect on the crystal vibrations are mainly of two types: (a) slight modifications on the energy of modes in the high frequency region; (b) large shifts of the energy modes in the low frequency region. These modifications are due to the fact that the pressure should greatly decrease the intermolecular distances while the interatomic separations of each molecule are only slightly decreased. Eventually, if the structure is greatly modified by the external agent (pressure) via hard changes in the intermolecular distances, it will adapt itself to a new configuration and the crystal symmetry will be changed. Associated to these changes it will be observed drastic modifications in the external modes region of the Raman spectra as have been already observed in another amino acid crystal.
In Fig. 1 we show Raman spectra of L-alanine crystal for selected pressures. The spectra are unpolarized but the laser beam is reaching the sample surface along the c-axis. All pressure values are given in kbar units. The spectrum taken at 0.7 kbar has a profile similar to the spectrum recorded at atmospheric pressure, except for little differences in the relative intensities. In this spectrum we labeled the external mode peaks as A, B, C, D, E and F, in order to discuss the effect of pressure in a most clear way. The peak A has an initial frequency of 41 cm<sup>-1</sup>, while the peaks B, C, D, E and F have frequencies at P = 0.0 kbar equal to 48, 75, 105, 113 and 138 cm<sup>-1</sup>, respectively. The 41 and 48 cm<sup>-1</sup> modes have been assigned as $`w`$-axis libration, in a picture where $`u,v`$ and $`w`$ are perpendicular axes and, additionally, the $`v`$ is along the long molecular axis, $`w`$ is perpendicular to the plane of the molecule and $`u`$ is nearly parallel to the crystallographic c-axis. It was observed through temperature studies an instability in the 48 cm<sup>-1</sup> mode associated to localization of vibrational energy in the 41 cm<sup>-1</sup> mode. Although the localization of vibrational energy should not be related to the pressure effect in a direct way, it is worth to mention that as pressure is increased, the peak with frequency 41 cm<sup>-1</sup> gains intensity, similar to the temperature effect. In fact we observe that at 0.7 kbar the intensity of the bands of 41 and 48 cm<sup>-1</sup> modes are almost the same and as the pressure increases, the difference between the intensity of the two modes increases considerably. However the main effect of pressure on the Raman spectra of L-alanine crystal is observed in another spectral region as discussed in a forthcoming paragraph.
The dependence with pressure of all L-alanine low frequency bands, can be seen by analysis of the plots of frequency ($`\omega `$) versus pressure (P) given in Fig. 2. In the plots the circles represent data taken while compressing the L-alanine crystal up to 32.1 kbar and the solid lines are the least square fitting of the data to the function:
$$\omega =\omega _0+\alpha P$$
(1)
A linear behavior was observed for all the $`\omega `$ vs. P plots of low frequency modes. The values found for the intercept, $`\omega _0`$, the linear pressure coefficient, $`\alpha `$, and the experimental frequency of the mode at P = 0 kbar, $`\omega _{0(exp.)}`$, are listed in Table I. It is important to state that for peak A, the linear fitting is a good one for data from 0 to 18 kbar (the $`\omega _0`$ value for peak A in Table I corresponds to this first linear fitting). Also, it is observed that from 20 to 32 kbar a linear fitting adjust very well to experimental data of peak A, but with a different dependence of frequency with pressure. This would mean that a slow modification of the L-alanine structure is starting to occur in a pressure value of about 20 kbar.
This modification, in fact, will be completed at a pressure of about 28.4 kbar. In Fig. 1 the spectrum taken at 26.8 kbar shows that peak E is just a single band and when pressure reaches 28.4 kbar a splitting of the peak is verified. The spectrum of 28.4 kbar in Fig. 1 shows this new peak arising from the splitting of peak E, and it is marked as G. In Fig. 2 the plots of $`\omega `$ versus P also illustrates well the splitting of peak E above 28 kbar. Such a splitting is accounted for as a result of change of the lattice symmetry caused by a modification of the elementary cell. In other words, the number of external modes increases when the L-alanine crystal is compressed above 28 kbar and a phase transition to a new structure takes place.
It may be worthwhile to mention and to compare the high pressure L-alanine crystal study with that did in another amino acid crystal, monohydrated L-asparagine, NH<sub>2</sub>CO(CH<sub>2</sub>)CH(NH<sub>2</sub>)COOH $``$ H<sub>2</sub>O . Monohydrated L-asparagine crystal undergoes phase transitions at 1.0, 6.0 and 13 kbar, pressure values below the pressure value of the phase transition in L-alanine. The phase transition observed on monohydrated L-asparagine was ascribed to be a consequence of shortening of the length of the hydrogen (H) bonds in the crystal structure. Indeed, there are seven hydrogen bonds involving all the hydrogen atoms attached to the two nitrogen atoms and to oxygen of water molecule. These H bonds are fundamental to link together adjacent chains of molecules, including the bonds originated from water molecule. On the other hand, the zwitterionic L-alanine does not involve any water molecule. It is likely that the greater stability under high pressure of L-alanine compared with the other amino acid crystal should be originated from the fact that the little glue group of monohydrated L-asparagine, the water molecule, is absent in the material we discussed in this letter. Central to our understanding of the difference of the pressure effect on the two amino acid crystals is the possibility of localization of the atoms participating of the hydrogen bonds. Of course, a more precise description of the shortening of the molecular links and the right mechanism of the phase transition in L-alanine crystal will be achieved by the localization of the heavy atoms of the molecule by an x-ray diffraction study. In summary, comparing the monohydrated L-asparagine crystal with the L-alanine crystal a more stable structure under pressure is achieved for the last material.
IV. CONCLUSIONS
The results of our pressure Raman light scattering studies performed for the L-alanine crystal have led to the following conclusion: for pressures up to 32 kbar the crystal undergoes a phase transition between 26.8 and 28.4 kbar. The transition observed in our work was identified by a drastic change in the Raman spectra, i.e., an increasing in the number of observed Raman peaks in the external mode region. The pressure coefficients for external mode was also given.
It is clear that deductions concerning structural changes based on the spectra reported here, although very clear, is not complete. X-ray diffraction measurements over this pressure range, both on the room pressure phase and on the new phase would be of great value for determining the new space group and understanding the mechanisms of the changes revealed in this study.
###### Acknowledgements.
We acknowledge the financial support from the FUNCAP, CAPES and FINEP, Brazilian agencies. We are indebted to Prof. O. Pilla for discussions and to Prof. M.P. Almeida for a reading of the manuscript. |
warning/0001/astro-ph0001001.html | ar5iv | text | # Statistics of Turbulence from Spectral-Line Data Cubes
## 1 Introduction
The interstellar medium is turbulent and the turbulence is crucial for understanding of various interstellar processes. Interstellar turbulence occurs in magnetized fluid and magnetic field establishes a connection between ISM phases (McKee & Ostriker 1977) thus making the turbulent cascade much more complex and coupling together cosmic rays and gas. Theoretical understanding of such a multiphase media with the injection of energy at different scales (Scalo 1987) is extremely challenging.
In terms of the topic of the present meeting, turbulence is important both for accelerating cosmic rays and for their diffusion. Indeed, whatever mechanism of cosmic ray acceleration we consider, its understanding requires proper accounting for scattering of cosmic rays by turbulent magnetic field. The same is true for the propagation of cosmic rays. For instance, if it were not for magnetic field lines wandering, the diffusion of cosmic rays perpendicular to the magnetic field direction would be suppressed (see Jokipii 1999). Moreover, it is becoming clear that particle streaming along magnetic field lines is also substantially influenced by magnetic turbulence.
In view of a broader picture, turbulence is widely believed to be an important element of molecular cloud dynamics and star formation process, although various authors disagree on the degree of its importance (see discussion in Vazques-Semadeni & Passot 1999). Undoubtedly turbulence is essential for heat transfer in the interstellar medium. It has been recently suggested that turbulence is also a key element to understanding various chemical reactions (Gredel 1999) and of the fundamental problem of MHD, namely, to the problem of fast magnetic reconnection (Lazarian & Vishniac 1999). This very limited and incomplete list of processes for which turbulence is essential explains the motivation behind the attempts to study interstellar turbulence.
Unfortunately interstellar turbulence remains a mystery in spite of all the attempts to study it. Substantial progress in numerical research (see Ostriker 1999, Vazquez-Semadeni & Passot 1999) is not adequate to reproduce the flows comparable in complexity and in Reynolds numbers, and the situation will not change in any foreseeable future. Thus only direct observational studies of interstellar turbulence may provide us with the crucial information on this phenomenon. Approaching the problem one would like to know at least the statistics of density, velocity and magnetic field. In this review I briefly discuss what information emission lines can supply us with. I would like to quote Alyssa Goodman, who believes that present day technology made spectral-line mapping of large portions of interstellar media “a booming cottage industry”. Attempts to use this wealth of observational data via visual inspection become fruitless and this calls for the introduction of more sophisticated techniques.
Statistical description is a nearly indispensable strategy when dealing with turbulence and a big advantage of statistical techniques is that they extract underlying regularities of the flow and reject incidental details. Attempts to study interstellar turbulence with statistical tools date as far back as the 1950s (see Horner 1951, Kampe de Feriet 1955, Munch 1958, Wilson et al. 1959) and various directions of research achieved various degree of success (see reviews by Kaplan & Pickelner 1970, Dickman 1985, Lazarian 1992, Armstrong, Rickett & Spangler 1995). Studies of turbulence statistics of ionized media were successful (see Spangler & Gwinn 1990, Narayan 1992) and provided the information of the statistics of plasma density<sup>1</sup><sup>1</sup>1Incidentally the found spectrum was close to a Kolmogorov one. at scales $`10^8`$-$`10^{15}`$ cm. This research profited a lot from clear understanding of processes of scintillations and scattering achieved by theorists (see Goodman & Narayan 1985, Narayan & Goodman 1989). At the same time the intrinsic limitations of the scincillations technique are due to the limited number of sampling directions and difficulties of getting velocity information.
Deficiencies in the theoretical description have been, to our mind, the major impediments to studies of turbulence using emission lines. For instance, important statistical studies of molecular clouds (Dickman 1985, Dickman & Kleiner 1985, Miesch & Bally 1994) have not achieved the success parallel to that in scintillation studies.
Potentially, studies of interstellar turbulence via emission lines can provide statistics of turbulence in various interstellar phases, including neutral gas. More importantly, velocity information allows one to distinguish between static structures and dynamical turbulence.
The difficulty of studying Doppler broadened lines stems from the fact that one has to account for both velocity and density fluctuations. Indeed, at any given velocity the fluctuation of emissivity may arise both from the actual blobs of gas moving at this velocity and from parcels of gas with different spatial positions but accidentally having the same component of velocity along the line of sight. Therefore fluctuations of emissivity at a given velocity would be expected even if the media were completely incompressible.
There exist various ways of dealing with position-position-velocity (henceforth PPV) data cubes. One of them is to identify clumps and to describe their statistics (see Stutzki & Gusten 1990, Williams, de Geus & Blitz 1994). Another is use a more traditional set of hydrodynamic tools like power spectra, structure functions etc. The two statistics are interrelated (see Stutzki et al. 1998), but in general the relation between various tools is non-trivial. It seems that for answering various questions different statistical tools are more suitable. Therefore it is very encouraging that a number of techniques, including Principal Component Analysis (see Heyer & Schloerb 1997) and Spectral Correlation Functions (Goodman 1999, Rosolowsky et al. 1999) have been recently introduced to the field.
In what follows I depart from a traditional statistical hydrodynamics and describe how the 3D velocity and density power spectra can be extracted from position-position-velocity (PPV) data cubes. This choice reflects my subjective preferences and partially motivated by the fact that this approach relates the long-studied 3D density and velocity statistics (e.g. power spectra) with the observational data. Even with this limitation the scope the subject is too broad and I shall mostly talk about atomic hydrogen (HI) studies, that can be viewed as a test case for the technique. I discuss advantages of using HI as a test case in section 2, the problem of space-velocity mapping in section 3 and spectra in velocity slices in section 4. The interpretation of 21 cm Galactic and SMC data is given in section 5. Possible anisotropies of statistics stemming from magnetic field are dealt with in section 6, where a new technique for statistical studies of magnetic field is suggested. I consider formation of emissivity enhancements that can be identified as filaments and clouds in section 7 and discuss the generalization of the technique in section 8. Being aware of the limitations of the traditional hydrodynamic description of turbulence, we describe alternative approaches, i.e. 2D Genus statistics, Spectral Correlation Functions and Bispectrum in section 9. A short discussion of the results is given in section 10.
## 2 HI as a Test Case
Atomic hydrogen is an important component of the interstellar media and many efforts have been devoted to its studies (see Burton 1992). In terms of turbulence studies it has a number of advantages. For one thing, when dealing with HI one may in most cases disregard self-absorption. There are two major reasons for that: self-absorption is small (Braun 1997, Higgs 1999) and as shown in Lazarian (1995, henceforth L95), small localized absorption features typical to HI only marginally influence the statistics on the scales larger than their size. For another thing, the pervasive distribution of neutral hydrogen presents a sharp contrast to the localized distribution of molecular species, and this alleviates problems related to averaging. Moreover, atomic hydrogen emissivity is proportional to the first power of atomic density and this simplifies the analysis.
HI has a substantial filling factor ($`20\%`$ or larger) in the Galactic disc and therefore its motions should reflect large scale galactic supersonic turbulence. At the same time, its statistics may have connection with the statistics of molecular clouds. An additional advantage of HI is that it can be studied not only within our Galaxy but for the nearby galaxies as well.
Another motivation for studies of HI statistics stems from the recent attempt to describe the structures in the Galactic hydrogen in order to estimate the fluctuations of microwave polarization arising from interstellar dust. This contribution is extremely important in view of present-day efforts in the Cosmic Microwave Background (CMB) research (see Prunet & Lazarian 1999, Draine & Lazarian 1999). Some of the studies, for instance one by Sethi, Prunet and Bouchet (1998), attempts to relate the statistics of density observed in the velocity space and the statistics of polarization fluctuations. If such a relation were possible, it would greatly alleviate the efforts to study polarization of cosmological origin. As an intermediate step in this work, however, one should relate the statistics in emissivity in the PPV space and density of HI in real space.
The timing for developing statistical tools for HI studies is also influenced by the fact that new large data cubes, e.g. the Canadian Galactic Survey data (see Higgs 1999), should become available soon.
## 3 Basic equations
### 3.1 Space-Velocity Mapping
Problem
The notion that the velocity fluctuation can influence emissivity within PPV data cubes is not new. Since the early-seventies Butler Burton on numerous occasions claimed the importance of velocity fluctuations for the interpretation of 21 cm data (Burton 1970, 1971). The ambiguities of inferring cloud properties from CO emission lines were discussed by Adler & Roberts (1992). Using N-cloud simulations of spiral disks they showed that many spurious effects appear because of velocity blending along line of sight. Recently a number of researchers doing numerics (Pichardo et al. 1999, Vazques-Semadeni 1999) pointed out that pixel-to-pixel correlation between the channel maps and the velocity slices of PPV data cubes tends to be larger with the velocity rather than the density field.
To describe power spectra of velocity and density fields, i.e. to express the interstellar statistics using the language that was so successful in hygrodynamics (Monin & Yaglom 1972), one needs to disentangle velocity and denisty contributions to 21 cm emissivity fluctuations.
Approach
A quantitative treatment of the effects of space-velocity mapping is given in Lazarian & Pogosyan (1999, henthforth LP99). There it is assumed that the velocity of a gas element can be presented as a sum of the regular part $`𝐯^{reg}`$ which can arise for instance from Galactic rotation, and a random, turbulent, part $`𝐮`$, so that $`𝐯^{obs}=𝐯^{reg}+𝐮`$. The mapping from real space to PPV coordinates corresponds to a transformation
$`𝐗_𝐬`$ $`=`$ $`𝐗`$
$`z_s`$ $`=`$ $`A\left[f^1z𝐮(𝐱)\widehat{𝐳}\right],`$ (1)
where henceforth we use large letters to denote vectors in the Position-Position plane (i.e. $`xy`$-plane) and use $`z_s`$ for the velocity coordinate. The parameter $`A`$ is just a conversion factor which specifies the units of $`z_s`$ coordinate and it is intuitively clear that this factor should not enter any final expressions for turbulence statistics. On the contrary, the shear parameter $`f=\left(\delta v_z^{reg}/\delta z\right)^1`$ is an important characteristic of the mapping and one expects it to influence our final results. For Galactic disc mapping it is convenient to choose $`A=f`$, while studies of isolated clouds correspond to a zero shear, i.e. $`f^10`$. As most work on HI has been done so far on Galactic disc HI, to simplify our presentation we use the former choice. With this definition of space-velocity mapping LP99 obtain the power spectrum $`P_s`$ in the PPV space:
$`\rho _s(𝐤)\rho _s^{}(𝐤^{})`$ $`=`$ $`P_s(𝐤)\delta (𝐤𝐤^{})`$
$`P_s(𝐤)`$ $`=`$ $`e^{f^2k_z^2v_T^2}{\displaystyle }d^3𝐫e^{i𝐤𝐫}\mathrm{\Xi }(𝐤,𝐫),𝐫=𝐱𝐱,`$ (2)
where the kernel is
$$\mathrm{\Xi }(𝐤,𝐫)=e^{ifk_z(u_z(𝐱)u_z(𝐱))}\rho (𝐱)\rho (𝐱).$$
(3)
In derivation of this expression it is explicitly assumed that the turbulence is statistically homogeneous in the real space coordinates and the average denoted by angular brackets $`\mathrm{}.`$ depends only on the vector separation between points. The density Fourier modes in PPV space $`\rho _s(𝐤)`$ are uncorrelated<sup>2</sup><sup>2</sup>2A treatment of turbulence within individual clouds is slightly different (LP99). which is reflected in $`\delta `$ function presence in the right-hand side of the first equation in (2). The factor $`e^{f^2k_z^2v_T^2}`$, where $`v_T`$ is a thermal velocity of atoms originates from averaging over thermal distribution and it shows that only supersonic motions are readily available for statistical studies. Note, that expressions similar to (2) and (3) were earlier obtained by Scoccimarro et al. (1999) in the framework of studies of Large Scale Structure of the Universe and this confirms the similarity of the problems studied in the two fields. However, the problem of “redshift-space” corrections to the statistics of galaxy distribution (Kaiser 1987) has been addressed either in the linear regime when perturbations are small or when the velocity contribution to the Fourier spectrum can be factorized by a Maxwellian factor (see Hamilton 1998). The problem that is dealt in turbulence case is much richer as one has to deal with non-linear density fields transformed by coherent velocities.
Note that velocity and density enter eq. (3) in a different way: velocity is in the exponent and density enters as the product $`\rho (𝐱)\rho (𝐱)`$. This provides an opportunity to disentangle the two contributions.
### 3.2 Spectrum in PPV Space
LP99 proves that in terms of final results for Lognormal distribution of density fluctuations and Gaussian distribution of velocity fluctuations it is safe to separate velocity and density in the following way
$$e^{if\mathrm{}}\rho (𝐱)\rho (𝐱)=e^{if\mathrm{}}\rho (𝐱)\rho (𝐱+𝐫),$$
(4)
even if density and velocity are strongly correlated. It is interesting to check the degree of uncertainty that the assumption (4) entails using numerically generated density and velocity fields.
For the sake of simplicity the density correlation function and velocity correlation tensor are considered to be isotropic in Galactic coordinates ($`xyz`$ space), i.e.
$$\rho (𝐱)\rho (𝐱+𝐫)=\xi (r)=\xi (𝐫).$$
(5)
$$\mathrm{\Delta }u_i\mathrm{\Delta }u_j=\left(D_{LL}(r)D_{NN}(r)\right)\frac{r_ir_j}{r^2}+D_{NN}(r)\delta _{ij},$$
(6)
where $`D_{LL}`$, $`D_{NN}`$ are longitudinal and transverse correlation functions respectively (Monin & Yaglom 1972), and $`\delta _{ik}`$ equals 1 for $`i=k`$ and zero otherwise. These assumptions are not necessary, as the treatment can be provided for instance for axisymmetric turbulent motions (see Oughton, Radler & Matthaeus 1997) as it is discussed in (L95).
The general expression for the 3D spectrum in PPV space is
$$P_s(𝐤)=e^{f^2k_z^2v_T^2/2}d^3𝐫e^{i𝐤𝐫}\xi (r)\mathrm{exp}\left[\frac{1}{2}f^2k_z^2D_z(𝐫)\right],$$
(7)
where
$$D_z(𝐫)\mathrm{\Delta }u_i\mathrm{\Delta }u_j\widehat{z}_i\widehat{z}_j=D_{NN}(r)+[D_{LL}(r)D_{NN}(r)]\mathrm{cos}^2\theta ,\mathrm{cos}\theta \widehat{𝐫}\widehat{𝐳}$$
(8)
is the projection of structure tensor to the $`z`$-axis. Expression (7) is quite general and can be used to relate arbitrary velocity and density statistics in galactic coordinates with the HI emissivity in the PPV space.
## 4 Spectra in Velocity Slices
Observations of Galactic HI (Green 1993) revealed two dimensional spectrum of intensity fluctuations (see L95) and this spectrum shows power-law behaivior. Similar power laws for Galactic data were obtained by Crovisier & Dickey (1983), Kalberla & Mebold (1983), Kalberla & Stenholm (1983) and for Small Magellanic Clouds (SMC) by Stanimirovic et al. (1999). Thus LP99 considered power law statistics, namely, of velocity $`𝒫_{3v}k^\nu `$ and density $`𝒫_{3\rho }k^n`$, where $`𝒫`$ is used to denote spectra in Galactic coordinates. Note, that $`n<0`$ and $`\nu <0`$ and $`D_zCr^m`$, where $`m=\nu 3`$. Power law spectra were also reported for molecular <sup>12</sup>CO (data from Heithausen & Thaddeus 1990 and Falgarone et al. 1998) and <sup>13</sup>CO (data from Heyer & Schloerb 1997) lines and it looks that power law spectra are quite generic for interstellar turbulence (Armstrong et al. 1997). Thus the assumption of a power law statistics does not tangibly constrain the range of applicability of the developed theory<sup>3</sup><sup>3</sup>3It is rather unnatural to expect that velocity and density spectra not being power laws conspire to produce power law emissivity..
For power-law spectra of density with $`n>3`$ the correlation functions are also power-law:
$$\xi (r)=\rho ^2\left(1+\left(\frac{r_0}{r}\right)^\gamma \right),\gamma =n+3>0.$$
(9)
Substituting Eq. (9) in (7) one can see that
$$P_s(|𝐊|,k_z)=\rho ^2\left[P_v(|𝐊|,k_z)+P_\rho (|𝐊|,k_z)\right],$$
(10)
where the part $`P_v`$ comes from integrating unity in Eq. (9) and the part $`P_\rho `$ comes from integration the $`\left(\frac{r_0}{r}\right)^\gamma `$ part. As we may see, the $`P_\rho `$ part is influenced by both velocity and density fluctuations, while $`P_v`$ part arises only from density fluctuations. LP99 show that an expression is similar to (10) valid for $`n<3`$.
The relation between 2D spectrum in velocity slices and the underlying 3D emissivity spectrum in the PPV space is given by
$$P_2(𝐊)|_{}\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}𝑑k_zP_s(𝐊,k_z)\mathrm{\hspace{0.17em}2}\left[(1\mathrm{cos}(k_z))/(k_z)^2\right],$$
(11)
where $`𝐊`$ denotes a 2D wavevector defined in the planes perpendicular to the line-of-sight and $`|_{}`$ reflects the dependence on the slice thickness. Equation (11) represents the one dimensional integral of the three dimensional spectrum with the window function given by the expression in square brackets. It is easy to see that the thinner is the velocity slice $``$, the larger the $`k_z`$ range for which the window function is close to unity and therefore more 3D modes contribute to 2D spectrum.
Substituting (10) into (11) one can see that the two dimensional spectrum can be presented as a sum
$$P_2(|𝐊|)=\rho ^2\left[P_{2v}(|𝐊|)+P_{2\rho }(|𝐊|)\right],$$
(12)
where the expressions for $`P_{2v}`$ and $`P_{2\rho }`$ are self-evident. To avoid possible misunderstanding I would like to stress that $`P_{2v}`$ and $`P_{2\rho }`$ are not 2D velocity and density spectra and, for instance, $`P_{2\rho }`$ depends both on velocity and density statistics.
Velocity fluctuations are most important for supersonic turbulence which is the case for cold HI. In this situation the following power-law asymptotics can be obtained (see Table 1):
$$\mathrm{thin}\mathrm{slice}:C|𝐊|^m\delta V^2$$
(13)
$$\mathrm{thick}\mathrm{slice}:C|𝐊|^m\delta V^2$$
(14)
In other words, if the velocity dispersion $`Cr^m`$ on the scale $`|𝐊|^1`$ is larger than the squared width of the channel (in velocity units) the slice is termed thin. If the opposite is true the slice is termed thick.
Thin Slices
It is easy to see from Table 1 that in thin slices the velocity mapping makes the spectra more shallow (as $`m>0`$). This means that velocity creates a lot of small scale structures. It is also evident that if the $`n<3`$ the $`P_{2v}`$ contribution dominates. In the opposite regime $`P_{2\rho }`$ contribution is important.
Thick Slices
The relative importance of $`P_{2v}`$ and $`P_{2\rho }`$ depends on whether $`n>3m/2`$ or $`n<3m/2`$. In the former case $`P_{2\rho }`$ dominates, while for the latter $`P_{2\rho }`$ becomes dominant only when the slice is very thick, i.e. a substantial portion of the line is integrated over. Indeed, it is easy to see that integration over the line-width leaves only density information. In the intermediate situation if the density spectrum is steep, i.e. $`n<3m/2`$ $`P_{2v}`$ provides most of the contribution to $`P_2`$.
For warm HI the thermal velocity dispersion is comparable with the turbulent one. Thus fluctuations of intensity arise mostly from density inhomogeneities and the analysis in L95 is applicable. The amplidude of fluctuations arising from the warm phase of HI is suppressed due thermal velocity smearing effects. Therefore in the mixture of the warm and cold phase the contribution of the cold phase to the measured spectrum is likely to dominate (LP99).
## 5 Statistics of Diffuse HI
### 5.1 Analysis of data
One of the most thorough jobs of obtaining 21 cm statistics was done by Green (1993). His observations of the HI emission were accomplished with the Synthesis Telescope of the Dominion Radio Astrophysical Observatory (DRAO) towards $`l=140^{},b=0^{}`$ ($`03^h03^m23^s,+58^{}06^{}20^{}`$, epoch 1950.0) and they revealed a power law spectrum of 2D intensity. This spectrum is proportional to $`P_2(𝐊)|_{}`$ and its interpretation depends on whether the slicing is thick or thin. To answer this question one has to estimate the dispersion of turbulent velocity at the scales under study and compare it to the velocity thickness of the slice (see eq.(13) and (14). Assuming that velocity variations at the scale $`30`$ pc amount to $`10`$ km/s and arise from the Kolmogorov turbulence, the structure functions of velocity are
$$D_{LL}(r)100\left(\frac{r}{30\mathrm{p}\mathrm{c}}\right)^{2/3}\mathrm{km}^2\mathrm{s}^2,$$
(15)
The width of the interferometer channels combined to give a single data point in Green’s dataset is $`\delta V=5.94\mathrm{km}/\mathrm{s}`$. The slice thickness in parsec is $`\delta Vf`$ pc, and varies from $`600`$ pc for the closest slices to $`2200`$ pc for the distant ones<sup>4</sup><sup>4</sup>4Note, that the cut-off due to thermal velocity (see section 2.1) in Warm Neutral Medium (see a table of idealized phases in Draine & Lazarian (1999) ) is $`6`$ km/s. If the WNM constitutes the dominant fraction of the neutral phase (Dickey 1995) then the velocity resolution above is optimal and no further decrease in $`\delta V`$ will result in getting new information. However, if close to Galactic plane Cold Neutral Media constitutes a substantial portion of mass, the increase of velocity resolution up to 1 km/s is desirable.. The wavenumber of transition from thin to think slice given by eq. (14) is equal $`0.16pc^1`$.
In figure 1 the turbulence scales covered by Green’s study are shown. The smallest $`|𝐊|`$ span from $`1/3`$ pc<sup>-1</sup> for the closest slices to $`1/200`$ pc<sup>-1</sup> for the distant ones. It is obvious that most of the measurements correspond to the thin slice regime.
As we mentioned earlier, whether $`P_{2v}`$ or $`P_{2\rho }`$ dominates the observed emissivity spectra in the thin slice regime depends on whether $`n`$ is larger or smaller than $`3`$. If $`n>3`$, $`P_{2\rho }`$ is dominant and the observations by Green (1993) reveal the spectrum with index $`n+m/2`$. For $`m=2/3`$ the spectrum of emissivity obtained by Green (1993), namely the emissivity with the index $`2.7`$, corresponds to $`n3`$. If, however, the density spectrum is steep (i.e. $`n<3`$), the fluctuations of 21 cm intensity observed by Green (1993) can arise from velocity fluctuations. In this case the spectral index is $`3+m/2`$. For $`m=2/3`$ one gets the slope $`8/32.7`$ which is exactly what is observed. The question now is whether nature conspires to create the density spectrum with $`n3`$ and thus make the slopes of $`P_{2v}`$ or $`P_{2\rho }`$ identical in thin slices or we observe $`P_{2\rho }`$, while $`n<3`$.
To answer this question one has to consider thick slices of data. Unfortunately, for thick slices of data one has to account for lines of sight being not parallel. This problem was studied by Lazarian (1994) for the case of density statistics, but the study has not been generalized so far for the case when both density and velocity contribute to emissivity. Thus external galaxies provide a better case for study thick slices. Data for Small Magellanic Cloud (SMC) in Stanimirovic et al. (1999) shows the steepening of the slope from $`3`$ for the slices obtained with the maximal velocity resolution to $`n3.5`$ for data integrated over the entire emission line (Stanimirovic, private communication) which corresponds to our theoretical predictions for long-wave dominated density spectrum with the index $`n3.5`$. Thus the set of Green’s and Stanimirovic’s data is consistent with an interpretation that both velocity and density exhibit spectra close to Kolmogorov.
A potential difficulty of “Kolmogorov” cascade interpretation is that the SMC data show power-law slope up to 4 kpc scale. To explain Doppler broaderning of molecular lines one has to accept that the energy is being injected at large scales. However scales of several kpc look excessive. Injection of energy at such large scales is possible in the form of superbubles, but the details and the very possibility of the cascade in these circumstances is unclear.
### 5.2 Further tests
The “Kolmogorov” interpretation above apparently needs further testings. There are various pieces of evidence that could be interpreted in favor of the spectrum of density being shallow, i.e. $`3`$. However, our analysis shows that this interpretation is not substantiated.
For instance, Braun (1999) reports a power law index of $`3`$ for the spectrum of 21 cm emission from structures near the North-East major axis of M31 galaxy. However, he uses not the whole spectral lines, as we do in our treatment but the maximal values of velocity only (Braun 1999, private communication). The interpretation of this result in terms of the power spectra as discussed in L95 and LP99 is impossible, as the treatment of data is very different.
Shallow spectrum of Far Infrared emission from dust (Wall & Waller 1998, Waller et al 1998) does not support the the shallow HI density spectrum either. According to Stanimirovic (1999, private communication) the shallow spectrum of Far Infrared emission when converted into dust column density provides a steep “Kolmogorov”-type spectrum.
One can argue that a possible hint in favor of HI density being short wave dominated comes from molecular data discussed in Stutzki et al. (1998). There for both <sup>12</sup>CO (data from Heithausen & Thaddeus (1990) and Falgarone et al. (1998)) and <sup>13</sup>CO (data from Heyer & Schloerb (1997)) transitions the spectrum of intensity was observed to have a power law index $`2.8`$. As the data is averaged over velocity, the fluctuations of intensity are due to density fluctuations and the spectrum of density should have the same slope as the spectrum of emissivity (L95), provided that the transitions are optically thin. The problem is that they are not thin and therefore the analysis above is not applicable.
Attempting to establish the actual underlying spectrum one may use one dimensional spectrum introduced in LP99
$$P_1(k_z)=𝑑𝐊P_s(𝐊,k_z)$$
(16)
Similar to two dimensional spectrum $`P_1`$ can be presented as a sum of $`P_{1\rho }`$, which scales as $`k_z^{2(n+2)/m}`$ and $`P_{1v}`$, which scales as $`k_z^{2/m}`$. Naturally, if $`n<3`$ $`P_{1v}`$ dominates, while $`P_{1\rho }`$ dominates for $`n>3`$. The analysis of data using $`P_1(k_z)`$ has not been done so far.
## 6 Anisotropies and Magnetic field
### 6.1 Goldreich-Shridhar Turbulence
It is natural to expect that dynamically important magnetic field makes interstellar turbulence anisotropic (Montgomery 1982, Higton 1984). Indeed, it gets difficult for hydrodynamic motions to bend magnetic fields at small scales if the energy density of the magnetic field and hydrodynamic motions are comparable at large scales. The turbulence in ionized gas has been found to be anisotropic and its Kolmogorov-type spectrum of plasma density fluctuations observed via radio scintillations and scattering (see Armstrong et al. 1995 and references therein) has been interpreted recently as a consequence of a new type of MHD cascade by Goldreich & Sridhar (1995). The Goldreich-Shridhar model of turbulence<sup>5</sup><sup>5</sup>5A qualitative discussion of the model and the role of reconnection for the cascade can be found in Lazarian & Vishniac (1999). differs considerably from the Kraichnan one (Iroshnikov 1963, Kraichnan 1964). It accounts for the fact that hydrodynamic motions can easily mix up magnetic field lines in the plane perpendicular to the direction of the mean field (see discussion in Lazarian & Vishniac 1999). Such motions provide eddies elongated in the field direction and the velocity spectrum close to the Kolmogorov one.
The Goldreich-Shridhar turbulence is anisotropic with eddies stretched along magnetic field. The wavevector component parallel to magnetic field $`k_{}`$ scales as $`k_{}^{2/3}`$, where $`k_{}`$ is a wavevector component perpendicular to the field. Thus the degree of anisotropy increases with the decrease of scale.
### 6.2 Anisotropies and magnetic field direction
It is both challenging and important to determine the degree of anisotropy for the HI statistics for various parts of the Galaxy. This information can provide an insight to the nature of HI turbulence and may be used as a diagnostic for the interstellar magnetic field. For instance, measurements of the structure functions of HI intensity
$$S(\theta ,\varphi )=(I(𝐞_1)I(𝐞_2))^2$$
(17)
as a function of a positional angle $`\varphi `$ for individual subsets of data should reveal magnetic field direction in various portions of the sky, if the turbulence is anisotropic as we expect it to be. This technique is somewhat analogous to a technique of finding magnetic field direction using the fluctuations of synchrotron radiation (see Lazarian 1992) but its applicability may be much wider.
So far, the attempts to measure anisotropy in HI are limited to the Green (1994) study, where no anisotropy was detected. Apparently a better analysis is needed. For the slices with high degree of anisotropy the statistical technique can be improved as suggested in L95.
## 7 Structures in PPV space
PPV data cubes, e.g. HI data cubes, exhibit a lot of small scale emissivity structure<sup>6</sup><sup>6</sup>6It was noticed by Langer, Wilson & Anderson (1993) that more structure is seen in PPV space than in the integrated intensity maps.. The question is what part of them is real, i.e. is associated with density enhancements in galactic coordinates and what part of them is produced by velocity fluctuations. A related question is whether the structures we see are produced dynamically, through forces, e.g. self-gravity, acting on the media or they may be produced statistically exhibiting the properties of random field. The second question was partially answered in Lazarian & Pogosyan (1997), where it was shown that density fluctuations with Gaussian distribution and power spectra result in filamentary structures. The structures become anisotropic and directed towards the observer when the velocity effects are accounted for<sup>7</sup><sup>7</sup>7There is a distant analogy between this effect and the “fingers of God” effect (see Peebles 1971) in the studies of large scale structure of the Universe. (see Lazarian 1999, fig. 2).
The issue of density enhancements produced by velocity fluctuations is closely related to the statistics of “clouds” observed in PPV space. The results on velocity mapping that we discussed earlier suggest that spectra of fluctuations observed in PPV velocity slices are more shallow than the underlying spectra. This means more power on small scales or, in other words, more small scale structure (“clouds”) appears in the PPV slices due to velocity fluctuations.
It has been believed for decades that emission cloud surveys (see Casoli et al. 1984, Sanders et al. 1985, Brand & Wouterloot 1995) provide a better handle on the actual spectrum of cloud mass and sizes than the extinction surveys (see Scalo 1985) because the velocity resolution is available. These two sorts of survey present different slopes for clump sizes and the difference cannot be accounted through occlusion of small clouds in the extinction surveys by larger ones (Scalo & Lazarian 1996). In view of the discussion above it looks that extinction survey may be closer to the truth, while a lot of structure detected via analyzing PPV cubes is due to velocity caustics. Paradoxically enough, emission data integrated over the spectral lines may provide a better handle on the distribution of cloud sizes compared to high resolution spectral-line data cubes. Averaging over velocity results in the distortions of the cloud size spectra due to occlusion effects, but these effects can be accounted for using the formulae from Scalo & Lazarian (1996).
## 8 Generalization of the technique
The formalism was described above in terms of HI power law statistics. It is obvious that it can be modified to deal with arbitrary statistics and with a variety of emission transitions. Here we briefly discuss complications which a generalization of the technique in order to be applicable to molecular clouds and ionized media may encounter.
### 8.1 Forward and inverse problems
A considerable number of researchers believes that self-similar behavior reflected in power law statistics is a characteristic feature of the interstellar turbulence including the molecular cloud turbulence (see Elmegreen & Falgarone 1996). However, some researchers (e.g. Williams 1999) see departures from a power law, e.g. signatures of a characteristic scale. In those cases, one still can find the underlying density and velocity statistics solving forward and inverse problems (see Lazarian 1999).
To solve the forward problem one needs to use expressions for the observable statistics, e.g. expressions for 2D and 1D spectra in PPV space (see eqs (7) and (16)) and fit the observable statistics varying the input of the velocity and density statistics. Naturally, the question of uniqueness for such solutions emerge, but with a reasonable choice of input parameters one may hope to avoid ambiguities.
A different approach involves the inversion of input data. Inversion also requires a model, but for the case of turbulence studies the model can be quite general and usually includes some symmetry assumptions, like the assumption of isotropy or axial symmetry of turbulence statistics (Lazarian 1994a, L95). For the case of turbulence the inversion has been developed for statistics of density (L95, Lazarian 1993) and magnetic field (Lazarian 1992, Lazarian 1994a). A remarkable property of the inversion for turbulence statistics is that it allows analytical solutions, which shows that the inversion is a well posed procedure in the mathematical sense. The procedure for inverting velocity data should be analogous to inverting density & magnetic field statistics, but has not been developed so far. We expect wide application of forward and especially inverse problem technique when the deviations from self-similar behavior become apparent in the data.
### 8.2 Various Transitions
As we discussed earlier, one of the advantages of using HI as a test case is that the emissivity is proportional to the column density. This is true for some optically thin transitions in molecular clouds, but fails when absorption is important. My analysis showed that the absorption is relatively easy to account for if it arises from dust, but much harder to deal with if it is self-adsorption. In the former case most of the analysis above is valid provided that the turbulence scale under study is much smaller than the extinction length.
Homogeneity of sample is another major concern for studies of turbulence in molecular clouds (see Houlahan & Scalo 1990). Filtering the data (see Miesch & Bally 1994), application of wavelets (see Stutzki et al. 1998) or both are required. However, it seems that as the resolution of data improves the effects of cloud edges get less important and easier to take care of.
Some emissivities, e.g. those of H$`\alpha `$ lines are proportional to the squared density of species. However, it is possible to generalize the technique above for those transitions and provide a quantitative treatment of turbulence in ionized emitting media, e.g. of HII regions (O’Dell 1986, O’Dell & Castaneda 1987).
## 9 Beyond Power Spectra
The approach that we discussed so far can be characterized as an interpretation of the emissivity spectra<sup>8</sup><sup>8</sup>8 In fact we do not distinguish between spectra and correlation functions. The two statistics are related via Fourier transform and provide an equivalent description. In some particular circumstances one or the other may be more convenient, however. in terms of the underlying statistics of velocity and density.
The advantage of this approach is that no numerical inversion (see Lazarian 1999) is performed and thus one should not worry about increase of the data noise. The power spectra are widely used in hydrodynamics and therefore there is hope to relate the statistics of interstellar turbulence with something simple and more familiar like Kolmogorov-type cascade.
In spite of all these advantages, the information that this approach can supply us with is limited. Indeed, media clustered by self-gravity and more diffuse media may have the same index of power spectrum, while being very different. In general, statistical measures borrowed from hydrodynamics may not be adequate while dealing with interstellar media. Indeed, we have to address particular questions, e.g. the question the identification of star-forming regions, which are beyond the standard hydrodynamic description. Therefore attempts to introduce new descriptors are worthy of high praise. It may happen that in answering specific questions one has to use particular descriptors.
### 9.1 Spectral Correlation Function
A new tool termed “spectral correlation function” or SCF has been recently introduced to the field (Goodman 1999, Rosolowsky et al. 1999). It compares neighboring spectra with each other. For this purpose the following measure is proposed:
$$S(T_1,T_0)_{s,l}1\left(\frac{D(T_1,T_2)}{s^2T_1^2(v)𝑑v+T_0^2(v)𝑑v}\right),$$
(18)
where the function
$$D(T_1,T_2)_{s,l}\left\{[sT_1(v+l)T_0(v)]^2𝑑v\right\}$$
(19)
and the parameters $`s`$ and $`l`$ can be adjusted. One way to choose them is to minimize $`D`$ function. In this case $`S`$ function gets sensitive to similarities in the shape of two profiles. Fixing $`l`$, $`s`$ or both parameters one can get another 3 function that are also sensitive to similarities in amplitude, velocity offset or to both parameters.
The purpose of those functions is to distinguish regions with various star forming activity and to compare numerical models with observations. To do this histograms of SCF are compared with histograms of SCF obtained for the randomized spatial positions. This allows to models to be distinguished on the basis of their clustering properties. First results reported by Rosolowsky et al. (1999) are very encouraging. It was possible to find differences for simulations corresponding to magnetized and unmagnetized media and for those data sets for which an earlier analysis by Falgarone et al. (1994) could not find the difference. The mathematical development of this new tool is under way (Padoan et al. 2000) and we expect new exciting results to be obtained in the near future.
A few comments about spectral correlation functions may be relevant. First of all, by its definition it is a very flexible tool. In the analysis of Rosolowsky et al. (1999) the SCF were calculated for the subcubes over which the original data was divided. In this way SCF preserves the spatial information and in some sense is similar to cloud-finding algorithms (see Stutzki & Gusten 1990, Williams, de Geus & Blitz 1994). However, one may fix the angular separation between the studied spectra and then the technique will be more similar to the traditional correlation function analysis that is sensitive to turbulence scale rather than to positional information. I believe that this avenue should be explored in future along with other more sophisticated techniques that can be applied to SCF. At first glance, it looks counterproductive to get a whole lot of various correlations using SCF as the input data. However, we must find a way of distinguishing regions with various physical properties and we are still in search for the best descriptors.
At the moment the distinction between various interstellar regions and the sets of simulated data is made by eye examining the histograms of SCF. With more information available it seems feasible to use wavelets that will emphasize some characteristics of the histograms in order to make the distinction quantitative. Construction these wavelets will be the way of “teaching” SCF to extract features that distinguish various sets of data.
### 9.2 Genus Statistics
The topology of ISM is an essential characteristic of the media. Genus analysis has been proved to be a useful tool for characterizing topology of the Universe (see Gott et al. 1989) and therefore it is tempting to apply it to the ISM studies.
The two dimensional genus analysis can be directly related to the media topology. By 2D maps we mean here maps integrated over the emission line, i.e. total intensity maps.
A two-dimensional genus is (Melott et al. 1989)
$$G_2(\nu _t)=(\mathrm{number}\mathrm{of}\mathrm{isolated}\mathrm{high}\mathrm{density}\mathrm{regions})(\mathrm{number}\mathrm{of}\mathrm{isolated}\mathrm{low}\mathrm{density}\mathrm{regions})$$
where $`\nu _t`$ denotes the dependence of genus on the threshold density in units of standard deviations from the mean. It is obvious that if one raises the density threshold from the mean, the low density regions coalesce and the genus becomes more positive. The opposite is true if $`\nu _t`$ decreases. Thus for Gaussian fluctuations one expects genus to be antisymmetric about zero, but the actual distributions should be able to reveal “bubble” or “meatball” topology of various parts of the ISM. Algorithms exist for calculating genus for 2D maps, e.g. microwave background maps (Colley et al. 1996) and therefore the application of genus statistics to interstellar maps is straightforward (and long overdue).
The 3D genus statistics (see Gott et al. 1989) in PPV space is less easy to interpret. As we discussed earlier, a lot of structures there are due to velocity caustics and the relation of the structures in galactic coordinates and PPV space is not obvious. However, it seems interesting to apply genus analysis to the PPV space in search for another statistical tool to distinguish various interstellar regions. After all, SCF introduced by Alyssa Goodman do not have a straightforward relation to the known parameters, but are very useful.
### 9.3 Bispectrum
Attempts to use multipoint statistics are a more traditional way to remove the constraints that the use of two point statistics, e.g. power spectra entails. Unfortunately, very high quality data is needed to obtain the multipoint statistics.
Among multipoint statistics, bispectrum (see Scoccimarro 1997) seems the most promising. This is partially because it has been successfully used in the studies of the Large Scale Structure of the Universe.
Bispectrum is a Fourier transform of the three point correlation function and if the power spectrum $`P(𝐤)`$ is defined as
$$\delta \rho (𝐤_1)\delta \rho (𝐤_2)=P(𝐤)\delta _D(𝐤_1+𝐤_2)$$
(20)
where $`\delta _D`$ is the Dirac delta function that is zero apart from the case when $`𝐤_1+𝐤_2=0`$, the bispectrum $`B_{123}`$ is
$$\delta \rho (𝐤_1)\delta \rho (𝐤_2)\delta \rho (𝐤_3)=B_{123}\delta _D(𝐤_1+𝐤_2+𝐤_3)$$
(21)
It is advantageous to use “hierarchical amplitude” (Fry & Seldner 1982) statistics
$$\mathrm{\Phi }_{123}\frac{B_{123}}{P(𝐤_1)P(𝐤_2)+P(𝐤_2)P(𝐤_3)+P(𝐤_3)P(𝐤_1)}$$
(22)
which is for power law spectra is a scale independent quantity.
In the studies of Large Scale Structure the hierarchical amplitudes were calculated for various initial conditions to compare with observations. In interstellar media it is advisable to compare various regions of sky using the tool. Impediments for the use of the technique stem from the increase of noise with the use of multipoint statistics and the problems of averaging along the line of sight. The problems should be addressed in the future.
### 9.4 Other techniques
A wavelet technique (see Gill & Henriksen 1990, Langer, Wilson & Anderson 1993, Rauzy, Lachieze-Rey & Henkiksen 1993) is discussed in this volume by Stutzki who proposed a so called $`\mathrm{\Delta }`$\- variance technique (Stutzki et al. 1998) which is is related to wavelet transforms (Zielinsky & Stutzki 1999). Wavelets potentially are a versatile tool that can filter out the large scale inhomogeneities of the data and concentrate the analysis on the scales of interest (see Stutzki, this volume).
Another useful statistical tool is the Principal Component Analysis (PCA). This tool was employed to spectral line imaging studies of the interstellar medium by Heyer & Schloerb (1997). The goal of the PCA is to determine a set of orthogonal “axes” $`u_{kl}`$ for which the the variance of the data is maximized. In the case of the data in $`n`$ points with $`p`$ velocity (spectrometer) channels for each point the data can be presented as $`\delta T_{ij}=T_{ij}T_{ij}_n`$, where $`T_{ij}`$ is the temperature at the channel $`j`$ at a position $`i`$ and $`\mathrm{}_n`$ denote averaging over positions. Maximazing the variance means maximizing the expression $`y_{ij}y_{ij}=u_{ik}S_{jk}u_{ij}`$, where summation over the repeating indexes is implied and $`S_{ik}=\delta T_{ij}\delta T_{jk}_n`$. In practice finding of $`u_{ij}`$ amounts to solving a set of eigenvalue equations $`S_{ik}u_{kj}=\lambda u_{ij}`$. To visualize the variance related to $`l`$-th principal component eigenimages are constructed from the projections of $`T_{ij}`$ onto the eigenvector, i.e. $`l`$th eigenimage at point $`(r_i)`$ is $`\delta T_{ij}u_{l,j}`$. Heyer & Schloerb (1997) showed that using PCA technique it is possible to decompose large-scale spectroscopic images of molecular clouds. Their analysis enabled them to calculate the velocity - scale relations for a number of cloud complexes. In terms of the statistical analysis presented above, PCA provides a means of filtering out large scale features responsible for the largest contribution to the global variance. This makes the sample more homogeneous and suitable for describing using correlation functions and power spectra. The potential of this statistical tool is to be further explored. It is likely that combining the various set of data (for instance, HI and CO) more interesting correlations can be obtained via PCA technique.
## 10 Discussion
It is not possible to cover all the various interesting approaches that have been tried in order to study interstellar turbulence via emission lines. For instance, we omitted a discussion of 3D correlation functions in PPV space introduced in Perault et al. (1986). We did not cover studies of turbulence using centroids of velocity (see Dickman 1985) either. One reason for this is that I believe that the statistics of velocity centroids have to be described in terms of underlying velocity and density.
A search for tools to deal with the interstellar turbulence has been intensified recently and this shows that there is deep understanding in the community that the wealth of observational data must be explored and it is essential to compare observations and numerical simulations. I personally believe that the development of theoretical approaches to dealing with data has become at this point not less important than obtaining the data.
Most of the present review I devoted to dealing with power spectra which reflects my personal preferences. Although far from being unambiguous, the power spectra were most intensively studied in hydrodynamics and the MHD theory and therefore they provide a bridge between idealizations that we partially understand and terra incognita of interstellar turbulence. Whether this approach is useful for a particular phase of the interstellar media is not clear a priori. We may or may not have any self-similarity indicating a turbulent cascade. However, at least for HI it seems that the approach is promising. Indeed, we managed to relate, although tentatively, the statistics of 21 cm emission with the statistics of Goldreich-Shridhar cascades. The next class of objects to study using the technique should be molecular clouds.
Although studies of molecular clouds are expected to face more problems, some of them mentioned earlier on, it is likely that the underlying 3D statistics will be soon obtained for the optically thin molecular lines. Comparison of this statistics with that in diffuse media should provide an insight to the nature of interstellar turbulent cascade and turbulent support of molecular clouds.
However, the limitations of the power-spectrum approach make it necessary to use alternative tools such as wavelets, genus statistics, principal component analysis and develop new ones such as spectral correlation function even though the relation between their output and the familiar notions from hydrodynamics is not always clear. In studies of interstellar medium one has to address particular questions, e.g. the question of star formation and therefore appearance of specialized tools is only natural.
## Summary
1. Velocity and density power spectra can be obtained from observed emissivities. Velocity fluctuations make emissivity spectra in velocity slices shallower. This results in much of small scale structure in PPV space that can be erroneously interpreted as interstellar clouds or clumps.
2. Turbulence is likely to be anisotropic with magnetic field defining the anisotropy direction. This should allow a new way of studying magnetic field.
3. The available wealth of observational data motivates the development of new tools for data handling.
## Acknowledgments
This review is partially based on the results that I obtained together with Dmitry Pogosyan. Discussions with Chris McKee, John Scalo, Steve Shore and Enrique Vazquez-Semadeni are acknowledged. I am grateful to Robert Braun and Snezana Stanimirovic for helpful input on data reduction procedures.
## References |
warning/0001/cond-mat0001179.html | ar5iv | text | # Quantum Dynamics in Nanoscale Magnets in Dissipative Environments
## 1 Introduction
Hysteresis phenomena of ferromagnets have been one of the most interesting problems in the magnetism and statistical physics. Mechanism of the coercive force has been investigated by studying the processes that lead to the critical nucleation and motion of the domain wall.$`^{\text{?}\text{}\text{?}\text{)}}`$ From the point of view of free energy of the system, the hysteresis phenomena have been discussed in terms of the relaxation process of the metastable state to the true equilibrium state in the picture shown in Fig. 1 (a). Usually the end point of hysteresis is related to the so-called spinodal point where the metastability disappears (see Fig. 1 (b)).
However there is some probability of relaxation from the metastable state B to the equilibrium state A. For example at finite temperatures the probability of thermal excitation to the top of barrier C is proportional to $`\mathrm{exp}(\mathrm{\Delta }E_\mathrm{B}/k_\mathrm{B}T)`$ and thus the relaxation rate $`p_{\mathrm{th}}`$ of the metastable state through this activation process is given by
$$p_{\mathrm{th}}=\frac{1}{\tau _0}e^{\frac{\mathrm{\Delta }E_\mathrm{B}}{k_\mathrm{B}T}},$$
(1.1)
i.e., the Arrhenius law. At low temperatures the relaxation time $`\tau =1/p_{\mathrm{th}}`$ of this process diverges exponentially.
The shapes of the free energy in Fig. 1 are given by the mean field theory, which gives a good intuitive picture of the metastability. Here it should be noted that $`\mathrm{\Delta }E_\mathrm{B}`$ should be a microscopic quantity. Because in a bulk system $`\mathrm{\Delta }E_\mathrm{B}`$ is of order the system size, the activation rate vanishes, i.e., $`p_{\mathrm{th}}=0`$. For a more quantitative understanding, we must look at the system microscopically. There the system is not uniform and we have to consider a microscopic break through of the metastable state. Such break through occurs as a process of creating the critical nucleus.$`^{\text{?)}}`$ For this microscopic process, Figs. 1(a)-(b) represent effective potentials of the size of the nucleus. The relaxation time of metastable states has been classified according to the size of system and the generating rate of the critical nuclei. There are two regions, i.e., a single nucleation region (stochastic region) and a multi nucleation region (Avrami region).$`^{\text{?}\text{)}}`$
When the size of the magnets becomes smaller than the width of the domain wall, the nucleus can not be defined. In such cases, the magnetization of the system changes uniformly and this process of breakdown of the metastable state is called ”coherent process”. Relaxations in this situation also have been studied extensively.$`^{\text{?}\text{)}}`$
It has been pointed out that quantum fluctuation may play an important role in such small systems. In order to detect such quantum processes, several experiments have been proposed.$`^{\text{?}\text{)}}`$ However distribution of particle sizes prevents to analyze their processes in simple ways.$`^{\text{?}\text{)}}`$ Studies on single magnetic particles have been also performed but clear evidence of quantum processes has not yet been found.$`^{\text{?}\text{)}}`$
In this respect, nanoscale molecular magnets such as Mn<sub>12</sub>,$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ Fe<sub>8</sub>,$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ and V<sub>15</sub>$`^{\text{?}\text{)}}`$ etc. are more promising. These molecules Mn<sub>12</sub> and Fe<sub>8</sub> consist of small number of atoms. The low energy state of the system is represented by an effective $`S=10`$ spin. Because interactions among molecules are very small, each atom can be regarded as a $`S=10`$ single spin. The Hamiltonian of the spin is generally given by
$$=DS_z^2HS_z+Q,$$
(1.2)
where $`S_z=10,9,\mathrm{}10`$ and $`Q`$ denotes a term which causes the quantum fluctuation, such as $`S_x,S_x^2S_y^2`$, or $`(S^+)^4+(S^{})^4`$. In these systems the energy levels as a function of the field have a discrete structure (Fig.2(a)). There we expect an explicit quantum mechanical dynamics.
Due to the term $`Q`$, a small energy gap is formed at each crossing point as shown in Fig.2. This structure is called avoided level crossing. When the field is swept through an avoided level crossing point, so called nonadiabatic transition occurs. Nonadiabatic transition plays important roles in microscopic quantum dynamics such as level dynamics of semiconductor, chemical reaction and optics. Nonadiabatic transitions in various cases have been reviewed by Nakamura.$`^{\text{?}\text{)}}`$
In the uniaxial magnets, the nonadiabatic transition as shown in Fig.2(b) occurs. Here the population coming in from the channel A is scattered to the channels B and C with probabilities $`p`$ and $`1p`$, respectively. Here the channel B is the ground state. Thus the scattering to B is an adiabatic change and corresponds to the tunneling. On the other hand, $`1p`$ is a probability to jump up to the channel C. This process corresponds to the case to stay in the metastable state (un-tunneling).
The states of the channels A and C are the same state when there is no quantum fluctuation i.e., $`Q=0`$. This unperturbed state is called the diabatic state. These channels A and C have similar states even in the presence of $`Q`$. The probability $`p`$ was studied by Landau,$`^{\text{?}\text{)}}`$ Zener$`^{\text{?}\text{)}}`$ and Stückelberg$`^{\text{?}\text{)}}`$ (LZS) and is given by
$$p=1\mathrm{exp}\left(\frac{\pi (\mathrm{\Delta }E)^2}{2c\mathrm{}g\mu _\mathrm{B}\mathrm{\Delta }m}\right),$$
(1.3)
where $`\mathrm{\Delta }E`$ is a gap at the avoided level crossing and $`\mathrm{\Delta }m`$ is the difference of magnetization of the levels, $`c`$ the speed of the sweeping field $`c=dH/dt`$. Thus the term $`g\mu _\mathrm{B}\mathrm{\Delta }mc`$ is the changing rate of the Zeeman energy.$`^{\text{?}\text{}\text{?}\text{}\text{?}\text{)}}`$ The probability $`p`$ plays an important role in quantum mechanical relaxation of the present system.
In this LZS type nonadiabatic transition the transition occurs only in the vicinity of the crossing points. The first property is an essential ingredient for the relaxation at discrete points of the magnetic field in the molecular magnets.$`^{\text{?)}}`$ Making use of this dependence we can estimate $`\mathrm{\Delta }E`$ from the change of magnetization $`\mathrm{\Delta }M`$ . Let the crossing levels have magnetizations of the corresponding adiabatic states $`m`$ and $`m^{}`$. The change of magnetization is given by
$$\mathrm{\Delta }M=pm^{}+(1p)mm=p(m^{}m).$$
(1.4)
In the cases where $`\mathrm{\Delta }E`$ is observed by other methods such as AC-susceptibility, this sweeping dependence of $`\mathrm{\Delta }M`$ would give a method to confirm that $`\mathrm{\Delta }E`$ really comes from the tunneling gap.$`^{\text{?)}}`$ However such confirmation has not yet been observed in small magnetic particles such as magnetic dots and ferritin. On the other hand the two characteristics have been observed in molecular magnets at least qualitatively. As a more peculiar property of the nonadiabatic transition, it has also been pointed out that the application of an alternate field at the resonant points will cause a nontrivial oscillation of magnetization due to phase interference.$`^{\text{?}\text{)}}`$
In order to study experimental data quantitatively, we need to incorporate the effects of environments. Effects of noise on the LZS transition have been studied by Kayanuma.$`^{\text{?}\text{)}}`$ The dynamics of molecular magnets in dissipative environments has also been investigated quite extensively.$`^{\text{?}\text{)}}`$
A thermal bath causes enhancement of the relaxation, e.g. the thermally assisted resonant tunneling, where resonant tunnelings of excited states play an important role.
In this paper we study the characteristics of the resonant tunneling affected by thermal disturbance. In particular, we study the effect of environment at very low temperatures such that relaxation process does not depend on the temperature. Even at such low temperatures, contact with the bath causes relaxation between levels which have magnetizations of the same sign. These states belong to the same valley in the potential picture of Fig. 1(a). The energy barrier does not exist between them. It is found that the relaxation between them easily occurs even with very weak disturbance. In pure quantum mechanical motion, transitions between these levels are almost prohibited except near the avoided level crossing points. Thus the magnetization curve with the dissipative effect is different from the one of a pure quantum case, although it does not depend on the temperature. We call such a process ”deceptive nonadiabatic transition”.$`^{\text{?}\text{)}}`$
At higher temperatures, excitation levels begin to contribute to the relaxation phenomena. At higher temperatures, alternate enhancements of relaxation at resonant points are observed,$`^{\text{?)}}`$ which is called ’the parity effect’. We consider the mechanism of such an alternation in a view point of nonadiabatic transition of excited states and find it as a general property of resonant tunneling of excited states reflecting the structure of energy levels.$`^{\text{?}\text{)}}`$ We also discuss the $`\sqrt{t}`$-dependence of initial decay at resonant points. Furthermore we study various cases of the LZS process in fluctuating random environments.
## 2 Numerical Method
The most standard method to study quantum dynamics in dissipative environments is the quantum master equation (QME) which describes the equation of motion of the reduced density matrix of the system $`\rho (t)`$ which is derived by tracing out the degrees of the freedom of the environment from the density matrix of the total system. The total system consists of a system $`_\mathrm{S}`$, a thermal bath $`_\mathrm{B}`$ and an interaction between them $`_\mathrm{I}`$:$`^{\text{?}\text{)}}`$
$$=_\mathrm{S}+_\mathrm{I}+_\mathrm{B}.$$
(2.1)
The reduced matrix is given by
$$\rho (t)=\mathrm{Tr}_\mathrm{B}e^\beta .$$
(2.2)
We have the following equation of motion for $`\rho (t)`$ in the limit of weak coupling, assuming that the correlation time of the bath variable is very short (Markovian approximation)
$$\frac{d}{dt}\rho (t)=\frac{1}{i\mathrm{}}[,\rho (t)]+\mathrm{\Gamma }\rho (t),$$
(2.3)
where $`\mathrm{\Gamma }`$ is a linear operator acting on $`\rho (t)`$. This equation has been used to study quantum dynamics of optical process, etc. In most cases $`\mathrm{\Gamma }\rho `$ has the so-called Lindblad form$`^{\text{?}\text{)}}`$
$$\mathrm{\Gamma }\rho =A^{}A\rho +\rho A^{}A+A^{}\rho A+A\rho A^{},$$
(2.4)
where $`A`$ is an operator of the system. However in multileveled phenomena $`\mathrm{\Gamma }\rho `$ has a more general form.
In the cases where the bath consists of an infinite number of bosons, a general expression can be derived.$`^{\text{?}\text{)}}`$
$$\frac{\rho (t)}{t}=\mathrm{i}[,\rho (t)]\lambda \left([X,R\rho (t)]+[X,R\rho (t)]^{}\right),$$
(2.5)
where
$`\overline{k}|R|\overline{m}`$ $`=`$ $`\zeta ({\displaystyle \frac{E_{\overline{k}}E_{\overline{m}}}{\mathrm{}}})n_\beta (E_{\overline{k}}E_{\overline{m}})\overline{k}|X|\overline{m},`$
$`\zeta (\omega )`$ $`=`$ $`I(\omega )I(\omega ),\mathrm{and}n_\beta (\omega )={\displaystyle \frac{1}{e^{\beta \omega }1}}.`$
Here $`\beta `$ is the inverse temperature of the reservoir $`1/T`$, and we set $`\mathrm{}`$ to be unity. $`|\overline{k}`$ and $`|\overline{m}`$ are the eigenstates of $``$ with the eigenenergies $`E_{\overline{k}}`$ and $`E_{\overline{m}}`$, respectively. $`I(\omega )`$ is the spectral density of the boson bath. Here we adopt the form $`I(\omega )=I_0\omega ^\alpha `$. When $`\alpha =1`$, it corresponds to the so called Ohmic bath and when $`\alpha =2`$, it corresponds to the phonon bath (super-Ohmic). As a more realistic bath for the experimental situation at very low temperature, we may take the dipole-field from other molecules and interactions with the nuclear spins.$`^{\text{?}\text{)}}`$ $`X`$ is an operator of system which is attached to bosons of the reservoir linearly, representing the interaction between the system and the thermal bath. In the present study we take $`X=\frac{1}{2}\left(S_x+S_z\right)`$. The relaxation process depends on the form of $`X`$. Generally the coupling with the transverse component $`X=S_x`$ is more efficient than that of the longitudinal one $`X=S_z`$ for the relaxation. Detailed comparisons among choices of the form of the coupling will be presented elsewhere.
For strong noise caused by fluctuating forces we can simulate quantum dynamics by solving the Schrödinger equation in random fields.$`^{\text{?}\text{)}}`$
## 3 Quantum Dynamics in Dissipative Environment
### 3.1 Deceptive nonadiabatic transition
In the lowest avoided level crossing point $`(S,S)`$ the change of magnetization $`\mathrm{\Delta }M`$ is given by (1.4) . However at higher crossing points $`(m,m^{})`$ with $`m^{}<S`$, the population scattered from $`m`$ to $`m^{}`$ is found to decay easily to the ground state, i.e., $`m^{}S`$, even when the dissipative effect is so small that the population at the metastable level of $`m`$ hardly decays. This difference can be easily understood from the intuitive picture of Fig. 1(a). That is, the relaxation in the same valley, i.e., $`m^{}S`$, is easy while the relaxation over the barrier $`mS`$ is hard. In this situation, we can not apply the relaxation (1.4) directly to estimate the LZS probability $`p`$. However we can still estimate $`p`$ using $`\mathrm{\Delta }M`$ because the relaxation from the level of $`m`$ occurs with the LZS probability and the relaxation to the ground state occurs in a rather short time. Taking these points into account, we modify the relation (1.4) by replacing the final magnetization $`m^{}`$ by $`S`$:
$$\mathrm{\Delta }M=pS(1p)mm=p(Sm).$$
(3.1)
In order to confirm these processes we performed simulations using the QME. First we confirmed that relaxation from the metastable point is unlikely to occur when the coupling between the bath and system is weak and the temperature is low. On the other hand, a fast relaxation is observed between levels with magnetizations of the same sign, which are in the same valley in Fig. 1(a). Furthermore when we sweep the field we find a step-wise magnetization curve whose step heights do not depend on the temperature but are definitely different from the pure quantum case. In Fig. 3, we show an example of magnetization process for $`T=0.1,\mathrm{\Gamma }=0.5`$ with very small effects of environment ($`\lambda =0.00001`$: a solid line) and that of pure quantum system ($`\lambda =0`$: dashed line). Both of them are not temperature dependent within this temperature range. We call this stepwise structure in dissipative environments ‘the deceptive nonadiabatic transition’. We find that we can correctly estimate the pure quantum transition probabilities using the relation (3.1). Thus even at very low temperatures the effect of the environment can not be excluded, but quantum mechanical processes and dissipative effects due to environments can be disentangled, and the information on the LZS probabilities can be extracted.
For the phenomena described above, the existence of the environment is important but the detailed nature is not important as long as it leads to fast relaxation to the ground state. If the environment causes a change of LZS probability, which would be possible when the sweeping rate is very slow, further consideration is necessary.$`^{\text{?)}}`$
### 3.2 Parity effect
At higher temperatures, excitation levels begin to contribute to the relaxation leading to temperature dependent phenomena. These processes would depend on the detailed characteristics of the bath and the ways of the coupling between the system and the bath. Therefore general description is difficult. However here we point out a general property of relaxation under these conditions.
As a characteristic of resonant tunneling at rather high temperatures, it has been observed that the amount of relaxation at the resonant points changes alternatively.$`^{\text{?)}}`$ Along a diabatic line, the energy gap increases monotonically as the difference of magnetizations $`|mm^{}|`$ of levels decreases. Thus the transition probabilities at the resonant points increase monotonically. In a perturbational treatment the energy gap depends on the difference as$`^{\text{?}\text{)}}`$
$$\mathrm{\Delta }E\left(\frac{\mathrm{\Gamma }}{D}\right)^{|mm^{}|}.$$
(3.2)
Thus we have to consider a mechanism of the alternate enhancements. Here we interpret it from the view point of resonant tunneling of excited state. The transition probabilities at resonant points with the same value of $`|mm^{}|`$ are nearly the same. Those points are located at the same horizontal level in Fig. 2. For example the values of $`p`$ given by (1.3) for the case of $`\mathrm{\Gamma }=0.45`$ with the sweeping speed $`c=0.0001`$ at the points, $`(8,5),(9,4)`$, and $`(10,3)`$ are 0.91, 0.64 and 0.99, respectively. On the other hand, those at $`(8,6),(9,5)`$ and $`(10,4)`$ are 0.72, 0.037 and 0.01, which are very small. Thus most of the population at the levels $`m=8,9`$ and $`10`$ decays at the former points. These decays cause enhancements of relaxation at $`H=0.3,0.5`$ and 0.7, which gives the parity effect. In Fig. 4, we show the magnetization of this case with its time derivative.
Because the energy structure shown in Fig. 2 is general for uniaxial magnets, we expect that the alternate enhancement of relaxation, i.e., the parity effect, is a general property of resonant tunneling in the thermal environment. We have also pointed out that if we change the sweeping rate the enhanced sequence is shifted. For example if we sweep much slower, the probabilities at $`(8,6),(9,5)`$ and $`(10,4)`$ become large and populations on the lines decay there, which causes the shift of the enhanced sequence at $`H=0.2,0.4`$ and 0.6.
### 3.3 Non-exponential decay at the resonant point
The magnetization which is initially polarized upward decays rather fast at a resonant point. Here the field is set at this point and is not swept. There are several paths for the magnetization to relax at this point. First let us consider relaxation by the nonadiabatic transitions at the lowest resonant point. Because at the resonant point the energy gap is very narrow, the field fluctuates around the point as shown in Fig. 5(a). If we regard the motion of the field as a Brownian motion, it is known that the number of the times the field crosses the resonant point is proportional to $`\sqrt{t}`$, i.e., the recurrence time of one-dimensional Brownian motion..$`^{\text{?}\text{)}}`$
At each crossing, the population moves to the other branch by the LZS transition probability
$$p_i=1\mathrm{exp}(\frac{\pi (\mathrm{\Delta }E)^2}{4c_i})\frac{\pi (\mathrm{\Delta }E)^2}{4c_i}.$$
(3.3)
In a pure quantum mechanical process, quantum mechanical interference occurs among the transitions.$`^{\text{?)}}`$ But assuming a fast decoherence, the total transition probability is expected to be given by
$$p_{\mathrm{total}}(t)\underset{i}{}p_i\sqrt{t}\frac{1}{c_i}\frac{\pi (\mathrm{\Delta }E)^2}{4}\alpha \sqrt{t}.$$
(3.4)
Thus we naturally expect that the magnetization decays as
$$\mathrm{\Delta }M=M_0(12\alpha \sqrt{t})$$
(3.5)
at the initial stage.
For a longer time scale, the field does not fluctuate freely but is confined near the resonant point. Thus $`p_{\mathrm{total}}`$ for long time is proportional to $`t`$. Therefore the magnetization decays in exponentially
$$M(t)e^{t/\tau },\tau =\frac{p_{\mathrm{total}}(t)}{t}.$$
(3.6)
This mechanism described above may give the simplest explanation of the $`\sqrt{t}`$ behavior. A more detailed analysis has been given in the reference ,$`^{\text{?)}}`$ taking into account the explicit nature of the fluctuation of the external field.
### 3.4 Modification of the transition rate
When the amplitude of the external disturbance is strong, we have to use another estimation of the transition probability. Kayamura and Nakayama have investigated LZS transition in fluctuating field and obtained expression of the transition probability. $`^{\text{?)}}`$ In the case where the sweeping rate is slow and the transient time through the resonant point is much longer than the phase coherence time, then they obtained the transition probability as
$$p=p_{\mathrm{SD}}=\frac{1}{2}(1\mathrm{exp}(\pi (\mathrm{\Delta }E)^2/2c\mathrm{}g\mu _\mathrm{B}\mathrm{\Delta }m)).$$
(3.7)
On the other hand, when the transient time is very short, transition probability does not change from that of the pure LZS transition
$$p=p_{\mathrm{LZS}}=1\mathrm{exp}(\pi (\mathrm{\Delta }E)^2/4c\mathrm{}g\mu _\mathrm{B}\mathrm{\Delta }m).$$
(3.8)
They confirmed such dependences by numerical simulation.
Next, let us consider a case where excited levels contribute to the relaxation. If the frequency of contacts between the system and the bath is high, a tunneling through the excited state would enhance the relaxation rate even the population at the excited state is very little. $`^{\text{?}\text{)}}`$
Let us consider the case where the LZS transition probability at the lowest level $`p_0`$ is very small and the one at the excited level $`p_1`$ is of the order 1. Let us consider a case where a state is excited to the exited level. In the off-resonant region this excited state decays to the original state very rapidly. The population at the excited state $`n_\mathrm{E}`$ is determined by the balance equation
$$n_\mathrm{E}R_{\mathrm{E}\mathrm{G}}=n_GR_{\mathrm{G}\mathrm{E}},$$
(3.9)
where $`n_\mathrm{G}`$ is the population in the ground state and $`R_{\mathrm{E}\mathrm{G}}`$ and $`R_{\mathrm{G}\mathrm{E}}`$ are transition rate from the excited state to the ground state and vise versa. Although $`n_\mathrm{E}=R_{\mathrm{G}\mathrm{E}}/R_{\mathrm{E}\mathrm{G}}n_G`$ is very small, $`R_{\mathrm{E}\mathrm{G}}`$ and $`R_{\mathrm{G}\mathrm{E}}`$ themselves can be very large.
At the resonant point, small fluctuations of the field would cause the crossing point as we see in Fig. 5(a). Thus the population pumped to the excited state A’ can be transferred to the state B’ and then it decays to B instead of A. (Fig.5(b)) This path (A$``$ A’$``$ B’$``$ B) becomes dominant when $`p_1p_0`$ and fluctuation of the field is rather fast. The opposite path (B$``$ B’$``$ A’$``$ A) is also larger. In the initial stage most of the population is at A and therefore the population moves along the former path. Thus the effective transition rate A $``$ B is enhanced very much.
$$p_{\mathrm{eff}}R_{\mathrm{G}\mathrm{E}}\nu p_1,$$
(3.10)
where $`\nu `$ is the frequency of the crossing. For a short time $`\nu `$ is proportional to $`\sqrt{t}`$ and for a long time it is proportional to $`t`$ as we saw above.
To study the transition probability in this case, the correct information for $`R_{\mathrm{G}\mathrm{E}}`$ and $`\nu `$ is necessary, but it is generally difficult. However, if we could estimate these quantities from the enhancement, it would yield detailed knowledge of the bath. In nanoscale molecular magnets, it would be possible to study such a detailed property, which is a very interesting research area in the future.
## 4 Summary and Discussion
Nanoscale molecular magnets display several phenomena which originate from explicit quantum mechanical transitions between discrete levels. In this paper we studied effects of dissipative environment which smears out the pure quantum processes. So far the relaxation processes have been a kind of black-box and have been treated only phenomenologically. But it would be possible to begin to study explicit processes of relation in nanoscale magnets because of their simple form.
So far we studied the $`S=10`$ spin representing the low energy structure of magnetic energy levels. Let us consider the structure of the full energy level. The molecular magnets have complicated structures. For example Mn<sub>12</sub> includes 12 Mn molecules with many other atoms which have nuclear spins. Thus the dimension of total Hamiltonian is $`5^84^4\times I`$, where $`I`$ comes from the degree of freedom of nuclear spins. This degree of nuclear spins causes random effects on each Mn atom. It would be an interesting problem to study how this random field on individual atoms causes changes of the energy levels at low temperatures.
Even without the effects of nuclear spins, there are dipole-dipole couplings among the molecules which cause random noise on the whole molecule. Thus it should be taken into account even in the view point of $`S=10`$ spin. Effects of this field are studied as ‘feedback-effect’ on the LZS process of magnetization.$`^{\text{?}\text{)}}`$ It would be interesting to study natures of noises explicitly in nanoscale molecular magnets.
## Acknowledgements
We would like to thank Professors Bernard Barbara, W. Wernsdorfer, Y. Kayanuma, H. Nakamura and H. Shibata for encouraging discussions. We also thank the Grant-in-Aid from the Ministry of Education, Science and Culture for the international collaboration program, which has helped the work discussed in the present article. The simulations have been made using the computational facility of the Super Computer Center of Institute for Solid State Physics, University of Tokyo, which is also appreciated. |
warning/0001/math0001033.html | ar5iv | text | # Askey-Wilson polynomials: an affine Hecke algebraic approach
## 1. Introduction
### 1.1.
Due to work of Cherednik , Macdonald , Noumi and Sahi , one can associate to every irreducible affine root system certain families of orthogonal polynomials (all closely related to the Macdonald polynomials), and prove their basic properties using a fundamental representation of the affine Hecke algebra in terms of difference-reflection operators. In this paper, we consider a rank one example of this theory in detail. The example is connected with a rank one non-reduced irreducible affine root system which has four orbits under the action of the associated affine Weyl group. The family of symmetric orthogonal polynomials associated to this particular affine root system is the celebrated four parameter family of Askey-Wilson polynomials, see .
### 1.2.
The four parameter family of Askey-Wilson polynomials has played an important and central role in the theory of basic hypergeometric orthogonal polynomials. In fact, up to date they seem to be the most general family of basic hypergeometric orthogonal polynomials which satisfy the additional requirement that they are joint eigenfunctions of a second-order $`q`$-difference operator. We use the link between the Askey-Wilson polynomials and the most general non-reduced affine root system of rank one (see 1.1) to derive in this paper the basic properties of the Askey-Wilson polynomials (and more!) from the algebraic structure of the associated (double) affine Hecke algebra.
### 1.3.
We introduce a Cherednik-Dunkl type difference-reflection operator $`Y`$ using the fundamental representation of the affine Hecke algebra of type $`\stackrel{~}{A}_1`$. The fundamental representation was defined by Noumi in the higher rank case, see also Sahi . Sahi’s non-symmetric Askey-Wilson polynomials are defined as the eigenfunctions of the Cherednik-Dunkl operator $`Y`$. They form a linear basis of the Laurent polynomials in one variable. We explicitly indicate their connection with the symmetric Askey-Wilson polynomials as originally defined by Askey and Wilson in . In particular, we give explicit expressions for the non-symmetric Askey-Wilson polynomials as a sum of two terminating balanced $`{}_{4}{}^{}\varphi _{3}^{}`$’s (here $`{}_{r}{}^{}\varphi _{s}^{}`$ is the basic hypergeometric series, see Gasper and Rahman for the definition). All the other results in this paper are derived without using the explicit series expansions of the (non-)symmetric Askey-Wilson polynomials in terms of basic hypergeometric series.
### 1.4.
We derive bi-orthogonality relations for the non-symmetric Askey-Wilson polynomials by explicitly computing the adjoint of the Cherednik-Dunkl operator $`Y`$ with respect to an explicit, complex measure. By a kind of symmetrization procedure, the bi-orthogonality relations imply Askey and Wilson’s orthogonality relations for the symmetric Askey-Wilson polynomials.
### 1.5.
We shortly describe Sahi’s results how an anti-isomorphism of the double affine Hecke algebra gives rise to a duality between the spectral and the geometric parameter of the (non-)symmetric Askey-Wilson polynomial. We use the duality to determine explicit intertwining properties of the action of the double affine Hecke algebra under the non-symmetric Askey-Wilson transform. Here the non-symmetric Askey-Wilson transform is a generalized Fourier transform which is defined with respect to the bi-orthogonality measure of the non-symmetric Askey-Wilson polynomials. This leads to the evaluation of the diagonal terms of the (bi-)orthogonality relations in terms of certain residues of the complex weight function. This technique is motivated by Cherednik’s approach to the study of the diagonal terms for non-symmetric Macdonald polynomials, in which he rewrites part of the action of the double affine Hecke algebra on non-symmetric Macdonald polynomials in terms of explicit operators acting on the spectral parameter.
### 1.6.
To complete the explicit computation of the diagonal terms, we need the evaluation of the constant term and the evaluation of the non-symmetric Askey-Wilson polynomial in a specific point (the latter playing a fundamental role in the duality arguments). We evaluate the constant term, which is essentially the well-known Askey-Wilson integral (see ), using shift operators. For the evaluation of the non-symmetric Askey-Wilson polynomial in a specific point we use a Rodrigues type formula for the non-symmetric Askey-Wilson polynomial in terms of Sahi’s intertwiners.
### 1.7.
The purpose of this paper is two-fold. First of all, we would like to show the power of the Cherednik-Macdonald theory in the study of basic hypergeometric orthogonal polynomials. It not only shows that all the basic properties of the Askey-Wilson polynomials can be obtained by natural algebraic manipulations, but it also reveals new and important insights in the structure of the Askey-Wilson polynomials.
Secondly, the affine Hecke algebraic approach also works for multivariable Askey-Wilson polynomials, the so called Koornwinder polynomials (which are associated with a higher rank non-reduced affine root system). The structure of the proofs in the higher rank setting are essentially the same, although the technicalities are more involved. Only part of the Cherednik-Macdonald theory associated with non-reduced affine root systems has been written down explicitly at this moment, see Noumi and Sahi . In our opinion, this paper can serve as one of the building blocks for obtaining a full understanding of the Cherednik-Macdonald theory in the case of non-reduced affine root systems. The higher rank case will be treated in an upcoming paper of the second author.
### 1.8.
In view of our aims described in 1.7, we have chosen to make this paper fairly self-contained. In particular, we have included some of the proofs of Noumi and Sahi , restricted to our present rank one setting. Furthermore, we have included several proofs which are fairly straightforward modifications from Cherednik’s and Macdonald’s work in case of Macdonald polynomials, see also for instance Opdam’s lecture notes for the classical $`q=1`$ setting.
### 1.9.
Finally we would like to point out the close connection with the paper of Kalnins and Miller. In the Askey-Wilson second order $`q`$-difference operator is written as the composition of a first order $`q`$-difference operator with its adjoint. This decomposition leads naturally to proofs of the orthogonality relations and of the quadratic norm evaluations for the symmetric Askey-Wilson polynomials (using shift principles). In our paper we use similar techniques, but we decompose the Askey-Wilson second order $`q`$-difference operator now as a sum of a difference-reflection operator (the Cherednik-Dunkl operator $`Y`$) and its inverse. This decomposition has the advantage that the Cherednik-Dunkl operator $`Y`$ itself satisfies a self-adjointness property with respect to a suitable extension of the Askey-Wilson orthogonality measure to a complex measure space for non-symmetric functions. This extra property of $`Y`$ naturally leads to the introduction of non-symmetric and anti-symmetric Askey-Wilson polynomials and to their bi-orthogonality relations.
### 1.10.
Notations: We use Gasper and Rahman’s notations for basic hypergeometric series and $`q`$-shifted factorials. We write $`_+=\{0,1,2,\mathrm{}\}`$ for the positive integers and $`=\{1,2,\mathrm{}\}`$ for the strictly positive integers.
### 1.11.
Acknowledgments: The second author is supported by a NWO-TALENT stipendium of the Netherlands Organization for Scientific Research (NWO). Part of the research was done while the second author was supported by the EC TMR network “Algebraic Lie Representations”, grant no. ERB FMRX-CT97-0100.
## 2. The Dunkl-Cherednik difference-reflection operators
### 2.1.
Let $`\widehat{}`$ be the vector space consisting of affine, linear transformation from $``$ to $``$. We identify $`\widehat{}`$ with $`\delta `$, where $`\delta `$ is the function identically one, and where $``$ acts by multiplication on itself. We introduce and study in this section a particular example of a rank one affine root system $`S\widehat{}`$. See Macdonald for the general discussion of affine root systems.
### 2.2.
Let $`.,.`$ be the positive semi-definite form on $`\widehat{}`$ defined by
$$\lambda +\mu \delta ,\lambda ^{}+\mu ^{}\delta =\lambda \lambda ^{},\lambda ,\lambda ^{},\mu ,\mu ^{}.$$
Then we can associate to every $`0f=\lambda +\mu \delta \widehat{}\delta `$ the involution $`s_f:\widehat{}\widehat{}`$ defined by
$$s_f(g)=gg,f^{}f,g\widehat{},$$
where $`f^{}:=2f/f,f`$. Observe that $`s_f`$ is an isometry with respect to $`.,.`$. In fact, we have $`s_f(g)=g\stackrel{~}{s}_f`$, where $`\stackrel{~}{s}_f\widehat{}`$ is the reflection in $`f^1(0)`$.
### 2.3.
We define now a subset $`S\widehat{}`$ by
$$S=\{\pm 1+\frac{m}{2}\delta ,\pm 2+m\delta |m\},$$
and we write $`𝒲=𝒲(S)`$ for the subgroup of invertible linear transformations of $`\widehat{}`$ generated by $`s_f`$, $`fS`$. Then $`S\widehat{}`$ is an affine root system. In particular, $`f,g^{}`$ for all $`f,gS`$, and $`S`$ is stable under the action of the affine Weyl group $`𝒲(S)`$.
### 2.4.
The gradient root system $`\mathrm{\Sigma }`$ of $`S`$ is the projection of $`S`$ on $``$ along the direct sum decomposition $`\widehat{}=\delta `$. Here $`\mathrm{\Sigma }=\{\pm 1,\pm 2\}`$, which is a non-reduced root system of type $`BC_1`$, with associated Weyl group $`W=\{1,s_1\}=\{\pm 1\}`$. Observe that $`W𝒲`$ acts on $`\widehat{}`$ by $`(\pm 1)(\lambda +\mu \delta )=\pm \lambda +\mu \delta `$.
### 2.5.
Let $`a_0=\delta 2S`$ and $`a_1=2S`$. Observe that $`a_0^{}=a_0/2=\frac{1}{2}\delta 1S`$, and $`a_1^{}=a_1/2=1S`$. Then $`\{a_0^{},a_1^{}\}`$ forms a basis of the affine root system $`S`$. We write $`S^+`$ for the positive roots in $`S`$ with respect to $`\{a_0^{},a_1^{}\}`$, so that $`S=S^+(S^+)`$ disjoint union. We furthermore set $`\mathrm{\Sigma }^+=\{a_1^{},a_1\}`$, which are the positive roots of $`\mathrm{\Sigma }`$ with respect to the basis $`\{a_1^{}\}`$. Observe that $`S^+=\mathrm{\Sigma }^+\{fS|f(0)>0\}`$.
### 2.6.
The affine Weyl group $`𝒲`$ is generated by the simple reflections $`s_0:=s_{a_0}=s_{a_0^{}}`$ and $`s_1=s_{a_1}=s_{a_1^{}}`$, while $`W`$ is generated by $`s_1`$. Observe that $`s_1s_0=\tau (1)`$, where $`\tau (\mu )`$ is the translation operator $`\tau (\mu )f=f+\mu ,f\delta `$ for $`f\widehat{}`$ and $`\mu `$. In particular, $`s_1s_0`$ has infinite order in $`𝒲`$ and
$$𝒲=W\tau ().$$
Furthermore, $`𝒲`$ is isomorphic to the Coxeter group with two generators $`s_0`$, $`s_1`$ and relations $`s_0^2=1`$, $`s_1^2=1`$.
### 2.7.
The affine root system $`S`$ has four $`𝒲`$-orbits, namely
$$\begin{array}{cc}\hfill S_s^1& =𝒲a_0^{}=\left(\frac{1}{2}+\right)\delta \pm 1,S_s^2=𝒲a_1^{}=\delta \pm 1,\hfill \\ \hfill S_l^1& =𝒲a_0=\left(1+2\right)\delta \pm 2,S_l^2=𝒲a_1=2\delta \pm 2.\hfill \end{array}$$
### 2.8.
We have the disjoint union $`S=R^{}R`$, with $`R=S_l^1S_l^2`$ a reduced, irreducible affine root system with basis $`\{a_0,a_1\}`$, affine Weyl group $`𝒲`$, and gradient root system $`\{\pm a_1\}`$, and with $`R^{}=S_s^1S_s^2`$ the corresponding affine co-root system. The co-root system $`R^{}`$ is a reduced, affine root system with basis $`\{a_0^{},a_1^{}\}`$, affine Weyl group $`𝒲`$, and gradient root system $`\{\pm a_1^{}\}`$. Similarly as for $`S`$, see 2.5, the fixed choice of basis give rise to a decomposition of $`R`$ and $`R^{}`$ in positive and negative roots (the positive roots are denoted by $`R^+`$ and $`R^{,+}`$, respectively). For example, we have $`R^+=\{a_1\}\{\pm a_1+\delta \}`$.
### 2.9.
Let $`\omega \widehat{}`$ be the involution $`\omega (x):=\frac{1}{2}x`$ ($`x`$), and consider $`\omega `$ as an involution of $`\widehat{}`$ by $`\omega :ff\omega `$ for $`f\widehat{}`$. Then $`\omega `$ preserves $`S`$. Furthermore, $`\omega =s_1\tau (\frac{1}{2})`$, $`\omega (a_i)=a_{1i}`$ and $`\omega s_i\omega =s_{1i}`$ for $`i=0,1`$. The subgroup $`𝒲^e`$ of the invertible linear transformations of $`\widehat{}`$ generated by $`𝒲`$ and $`\omega `$ is called the extended affine Weyl group. It is isomorphic to $`𝒲\mathrm{\Omega }`$, where $`\mathrm{\Omega }`$ is the subgroup of order two generated by $`\omega `$.
### 2.10.
Set $`𝒜:=[x^{\pm 1}]`$ for the algebra of Laurent polynomials in one indeterminate $`x`$. We set $`x^f:=q^\lambda x^\mu 𝒜`$ for $`f=\mu +\lambda \delta +\delta `$, where $`q`$ is a fixed non-zero complex number. Observe that $`x^a𝒜`$ is well defined for $`aS`$ and that $`𝒲^e`$ preserves $`+\delta `$. Furthermore, $`w(x^\mu ):=x^{w(\mu )}`$ for $`\mu `$ and $`w𝒲^e`$ extends to an action of $`𝒲^e`$ on $`𝒜`$ by linearity. In particular,
$$s_0(x^m)=q^mx^m,s_1(x^m)=x^m,m.$$
Observe that $`\tau (\mu )`$ ($`\mu `$) acts as a $`q`$-difference operator: $`\tau (\mu )(x^m)=q^{\mu m}x^m`$ for all $`m`$.
### 2.11.
A multiplicity function $`\underset{¯}{t}=\{t_a\}_{aS}`$ of $`S`$ is a choice of non-zero complex numbers $`t_a`$ ($`aS`$) such that $`t_{w(a)}=t_a`$ for all $`aS`$ and all $`wW`$. We use the convention that $`t_f=1`$ for all $`f\widehat{}S`$. A multiplicity function $`\underset{¯}{t}`$ of $`S`$ is determined by the values $`k_0:=t_{a_0}`$, $`u_0:=t_{a_0^{}}`$, $`k_1:=t_{a_1}`$ and $`u_1:=t_{a_1^{}}`$, see 2.7. Later on it will be necessary to impose (generic) conditions on the parameters $`t_f`$ ($`fS`$) and on the deformation parameter $`q`$. Until section 6 it suffices to assume that $`|q|1`$ and that $`k_0^2,k_1^2,u_1^2\pm q^{}`$. These conditions are assumed to hold throughout the remainder of the paper, unless specified explicitly otherwise.
### 2.12.
The Hecke algebra $`H_0=H_0(k_1)`$ of type $`A_1`$ is the unital, associative $``$-algebra with generator $`T_1`$ and relation $`(T_1k_1)(T_1+k_1^1)=0`$. Observe that $`\{1,T_1\}`$ is a linear basis of $`H_0`$ and that $`T_1`$ is invertible in $`H_0`$ with inverse $`T_1^1=T_1+k_1^1k_1`$.
### 2.13.
The affine Hecke algebra $`H=H(R;k_0,k_1)`$ of type $`\stackrel{~}{A}_1`$ is the unital $``$-algebra with generators $`T_0`$ and $`T_1`$ and relations
$$(T_ik_i)(T_i+k_i^1)=0,i=0,1.$$
Similarly as for $`H_0`$ we have that $`T_i`$ is invertible in $`H`$ with inverse $`T_i^1=T_i+k_i^1k_i`$.
### 2.14.
For $`w𝒲`$, let $`w=s_{i_1}s_{i_2}\mathrm{}s_{i_r}`$ be a reduced expression, i.e. a minimal expression of $`w`$ as product of the simple reflections $`s_0`$ and $`s_1`$. Then $`T_w:=T_{i_1}T_{i_2}\mathrm{}T_{i_r}`$ is well-defined and $`\{T_w\}_{w𝒲}`$ is a linear basis of $`H(R;k_0,k_1)`$, see . In particular, we may regard $`H_0`$ as a subalgebra of $`H`$.
### 2.15.
We set $`Y:=T_{\tau (1)}=T_1T_0H`$. By , we known that $`Y`$ is algebraically independent in $`H`$. Let $`[Y^{\pm 1}]H`$ be the commutative subalgebra generated by $`Y^{\pm 1}`$. Then
$$H_0(k_1)[Y^{\pm 1}]H(R;k_0,k_1)[Y^{\pm 1}]H_0(k_1)$$
as linear spaces, where the isomorphisms are given by multiplication. In particular, $`\{Y^m,Y^nT_1\}_{m,n}`$ and $`\{Y^m,T_1Y^n\}_{m,n}`$ are linear bases of $`H(R;k_0,k_1)`$, see \[13, proposition 3.7\].
In the remainder of the paper we identify $`𝒜`$ with $`[Y^{\pm 1}]`$ as algebra by identifying the indeterminate $`x`$ of $`𝒜`$ with $`Y`$. In particular, we write $`f(Y)=_kc_kY^k[Y^{\pm 1}]`$ for $`f(x)=_kc_kx^k𝒜`$.
### 2.16.
By Lusztig \[13, proposition 3.6\], we have the fundamental commutation relations
$$T_1f(Y)f(Y^1)T_1=\left((k_1k_1^1)Y^2+(k_0k_0^1)Y\right)\left(\frac{f(Y^1)f(Y)}{1Y^2}\right)$$
in $`H(R;k_0,k_1)`$ for all $`f(Y)[Y^{\pm 1}]`$. Indeed, observe that if the formula holds for $`f(Y)`$ and $`g(Y)`$, then it also holds for $`f(Y)g(Y)`$. It thus suffices to prove it for $`f(Y)=Y^{\pm 1}`$, in which case it follows from an elementary computation using the definition of $`Y`$ and the quadratic relations for the $`T_i`$, see 2.15 and 2.13, respectively.
### 2.17.
The following result was observed by Sahi \[19, theorem 5.1\] in the higher rank setting.
###### Corollary (The non-affine intertwiner).
Set $`S_1=[T_1,Y]=T_1YYT_1H`$. Then $`f(Y)S_1=S_1f(Y^1)`$ in $`H`$ for all $`f(Y)[Y^{\pm 1}]`$.
###### Proof.
This follows immediately from the definition of $`S_1`$ and from Lusztig’s commutation relation 2.16. ∎
### 2.18.
Another important consequence of Lusztig’s commutation relation 2.16 is the following result.
###### Corollary.
The affine Hecke algebra $`H=H(R;k_0,k_1)`$ acts on $`[Y^{\pm 1}]`$ by
$$\begin{array}{cc}\hfill T_1.g(Y)& =k_1g(Y^1)+\left((k_1k_1^1)Y^2+(k_0k_0^1)Y\right)\left(\frac{g(Y^1)g(Y)}{1Y^2}\right)\hfill \\ & =k_1g(Y)+k_1^1\frac{(1k_0k_1Y^1)(1+k_0^1k_1Y^1)}{(1Y^2)}(g(Y^1)g(Y)),\hfill \\ \hfill f(Y).g(Y)& =f(Y)g(Y)\hfill \end{array}$$
for all $`f(Y),g(Y)[Y^{\pm 1}]`$.
###### Proof.
Let $`\chi `$ be the character of $`H_0(k_1)`$ which maps $`T_1`$ to $`k_1`$. By 2.15 we may identify the representation space of the induced representation $`\text{Ind}_{H_0}^H(\chi )=H_\chi `$ with $`[Y^{\pm 1}]`$. By 2.16 the corresponding induced action of $`H`$ on $`[Y^{\pm 1}]`$ is as indicated in the statement of the corollary. ∎
### 2.19.
We define linear operators $`\widehat{T}_i\text{End}_{}(𝒜)`$ by
$$\begin{array}{cc}\hfill \widehat{T}_i:=& k_i+k_i^1\frac{\left(1k_iu_ix^{a_i^{}}\right)\left(1+k_iu_i^1x^{a_i^{}}\right)}{1x^{a_i}}(s_i\text{id})\hfill \\ \hfill =& k_is_i+\frac{(k_ik_i^1)+(u_iu_i^1)x^{a_i^{}}}{(1x^{a_i})}(\text{id}s_i),i=0,1.\hfill \end{array}$$
The following theorem was proved by Noumi \[16, section 3\] in the higher rank setting, see also \[19, section 2.3\].
###### Theorem.
The application $`T_i\widehat{T}_i`$ ($`i=0,1`$) extends uniquely to an algebra homomorphism $`\pi _{\underset{¯}{t},q}:H(R;k_0,k_1)\text{End}_{}(𝒜)`$.
###### Proof.
We identify $`[Y^{\pm 1}]H(R;u_1,k_1)`$ with $`𝒜`$ as algebra by identifying $`Y`$ with $`x^1`$. Then it follows from corollary 2.18 (applied to $`H(R;u_1,k_1)`$) that $`(\widehat{T}_1k_1)(\widehat{T}_1+k_1^1)=0`$ in $`\text{End}_{}(𝒜)`$. Conjugating $`\widehat{T}_1`$ with the involution $`\omega `$, see 2.9, and replacing $`k_1`$ and $`u_1`$ by $`k_0`$ and $`u_0`$ respectively, we see that $`(\widehat{T}_0k_0)(\widehat{T}_0+k_0^1)=0`$ in $`\text{End}(𝒜)`$. The theorem follows, since we have shown that all the defining relations 2.13 of $`H(R;k_0,k_1)`$ are satisfied by the linear operators $`\widehat{T}_i\text{End}_{}(𝒜)`$ ($`i=0,1`$). ∎
### 2.20.
Observe that the linear operator $`\widehat{T}_0`$ on $`𝒜`$ has a reflection and a $`q`$-difference part, while $`\widehat{T}_1`$ has only a reflection part, see 2.10. The operators $`\widehat{T}_i\text{End}_{}(𝒜)`$ ($`i=0,1`$) are called the difference-reflection operators associated with $`S`$. In the remainder of the paper, we simply write $`T_i`$ for the difference-reflection operators $`\widehat{T}_i`$ ($`i=0,1`$) if no confusion is possible. The operator $`Y=T_1T_0\text{End}_{}(𝒜)`$ is called the Cherednik-Dunkl operator associated with $`S`$.
### 2.21.
The representation $`\pi _{\underset{¯}{t},q}`$ has two extra degrees of freedom $`u_0`$ and $`u_1`$ besides the deformation parameter $`q`$ (which already lives on the affine Weyl group level, see 2.10). The motivation to label these two degrees of freedom in this particular way comes from the theory of double affine Hecke algebras. The double affine Hecke algebra $`(S;\underset{¯}{t};q)`$ associated with the affine root system $`S`$ (see ) is the subalgebra of $`\text{End}_{}(𝒜)`$ generated by $`\pi _{\underset{¯}{t},q}(H(R;k_0,k_1))`$ and $`𝒜`$, where we consider $`𝒜`$ as a subalgebra of $`\text{End}_{}(𝒜)`$ via its regular representation. We write $`f(z)\text{End}\left(𝒜\right)`$ for the Laurent polynomial $`f(x)𝒜`$ regarded as a linear endomorphism of $`𝒜`$. By the second formula for the difference-reflection operators $`T_i`$ in 2.19 we have that
$$f(z)T_iT_i(s_if)(z)=\frac{(k_ik_i^1)+(u_iu_i^1)z^{a_i^{}}}{(1z^{a_i})}\left(f(z)(s_if)(z)\right),f𝒜$$
in $`(S;\underset{¯}{t},q)`$ for $`i=0,1`$.
### 2.22.
The labeling of the extra degrees of freedom in the representation $`\pi _{\underset{¯}{t},q}`$ is now justified by the following theorem, together with 2.11.
###### Theorem.
$`(S;\underset{¯}{t};q)`$ is isomorphic as algebra to the unital, associative $``$-algebra $`(\underset{¯}{t};q)`$ with generators $`V_0^{},V_0,V_1,V_1^{}`$ and relations:
1. The application $`T_iV_i`$ for $`i=0,1`$ extends to an algebra homomorphism $`H(R;k_0,k_1)(\underset{¯}{t};q)`$.
2. The application $`T_iV_i^{}`$ for $`i=0,1`$ extends to an algebra homomorphism $`H(R;u_0,u_1)(\underset{¯}{t};q)`$.
3. (Compatibility). $`V_1^{}V_1V_0V_0^{}=q^{1/2}`$.
The isomorphism $`\varphi :(\underset{¯}{t};q)(S;\underset{¯}{t};q)`$ is explicitly given by $`\varphi (V_i)=T_i`$ ($`i=0,1`$), $`\varphi (V_0^{})=T_0^1z^{a_0^{}}=q^{1/2}T_0^1z`$ and $`\varphi (V_1^{})=z^{a_1^{}}T_1^1=z^1T_1^1`$.
The existence of the algebra homomorphism $`\varphi `$ follows by direct computations using 2.21. It is immediate that $`\varphi `$ is surjective. The injectivity of $`\varphi `$ requires a detailed study of the difference-reflection operators $`T_i`$ associated with $`S`$. We give the proof in 8.3.
## 3. Non-symmetric Askey-Wilson polynomials
### 3.1.
For $`aR`$, let $`(a)\text{End}(𝒜)`$ be the difference-reflection operator defined by
$$(a):=t_as_a+t_a^1\frac{(1t_at_{a/2}x^{a/2})(1+t_at_{a/2}^1x^{a/2})}{(1x^a)}\left(1s_a\right).$$
Then it is immediate that $`(a_i)=T_is_i`$ for $`i=0,1`$, where $`T_i`$ is the difference-reflection operator associated with $`S`$ (see 2.19). Furthermore, we have $`w(a)w^1=(w(a))`$ for all $`w𝒲`$ and all $`aR`$. Since any $`aR`$ is conjugate to $`a_0`$ or $`a_1`$ under the action of $`𝒲`$, we obtain from the quadratic relations for the difference-reflection operators $`T_i`$ (see 2.13 and 2.19) that
$$(a)^1=(a)+(t_at_a^1)s_a,aR.$$
### 3.2.
We define a total order $``$ on the basis of monomials $`\{x^m\}_m`$ of $`𝒜`$ by
$$1x^1xx^2x^2\mathrm{}.$$
Observe that this order is not well behaved under multiplication of the monomials: if $`x^{m_i}x^{n_i}`$ ($`i=1,2`$), then not necessarily $`x^{m_1+m_2}x^{n_1+n_2}`$.
### 3.3.
Let $`ϵ:\{\pm 1\}`$ be the function which sends a positive integer to $`1`$ and a strictly negative integer to $`1`$.
###### Lemma.
Let $`aR`$ be of the form $`a=2+k\delta `$ with $`k`$ (see 2.8) and let $`m`$. Then
$$(a)\left(x^m\right)=t_a^{ϵ(m)}x^m+\text{lower order terms w.r.t.}.$$
###### Proof.
For $`aR`$, let $`D_a\text{End}(𝒜)`$ be the divided difference operator defined by
$$D_af:=\frac{fs_af}{1x^a},f𝒜.$$
Then $`D_a(1)=0`$ and
$$\begin{array}{c}\hfill D_a\left(x^m\right)=\{\begin{array}{cc}x^{ma}x^{m2a}\mathrm{}x^{mm,a^{}a}\hfill & \text{ if }m,a^{},\hfill \\ x^m+x^{m+a}+\mathrm{}+x^{m(1+m,a^{})a}\hfill & \text{ if }m,a^{}.\hfill \end{array}\end{array}$$
Observe that $`m,a^{}=m`$ when $`a=2+k\delta `$ for some $`k`$. The lemma is now immediate when $`m_+`$. For $`m`$, we first observe that the coefficient of $`x^m`$ in the expansion of $`(a)(x^m)`$ in terms of monomials is zero. Indeed, the coefficient of $`x^m`$ in
$$t_a^1(1t_at_{a/2}x^{a/2})(1+t_at_{a/2}^1x^{a/2})D_a(x^m)$$
is $`t_aq^{mk}`$, which cancels with the coefficient of $`t_as_a(x^m)=t_aq^{mk}x^m`$. Hence the highest order term of $`(a)(x^m)`$ is $`t_a^1x^m`$ when $`m`$. This completes the proof of the lemma. ∎
### 3.4.
Lemma 3.3 implies the following triangularity property of the Cherednik-Dunkl operator $`Y`$. Set $`\gamma _m:=k_0^{ϵ(m)}k_1^{ϵ(m)}q^m`$ for $`m`$.
###### Proposition.
For all $`m`$, we have
$$Y(x^m)=\gamma _mx^m+\text{lower order terms w.r.t.}.$$
###### Proof.
Observe that $`Y=T_1T_0=(a_1)s_1(a_0)s_0=(a_1)(s_1(a_0))\tau (1)`$. Now $`s_1(a_0)=2+\delta `$ and $`\tau (1)(x^m)=q^mx^m`$ (see 2.10), so the proposition follows from lemma 3.3. ∎
### 3.5.
The diagonal terms $`\gamma _m`$ ($`m`$) of the triangular operator $`Y`$ are pair-wise different by the generic conditions 2.11 on $`q`$ and on the multiplicity function $`\underset{¯}{t}`$. Hence proposition 3.4 leads immediately to the following proposition (compare with Sahi \[19, section 6\] for the higher rank setting).
###### Proposition.
There exists a unique basis $`\{P_m()=P_m(;\underset{¯}{t};q)|m\}`$ of $`𝒜`$ such that
1. $`P_m(x)=x^m+`$ lower order terms with respect to $``$,
2. $`Y(P_m)=\gamma _mP_m`$
for all $`m`$.
###### Definition.
The Laurent polynomial $`P_m=P_m(;\underset{¯}{t};q)`$ ($`m`$) is called the monic, non-symmetric Askey-Wilson polynomial of degree $`m`$.
We will justify this terminology in section 5, where we relate the non-symmetric Askey-Wilson polynomials with the well-known symmetric Askey-Wilson polynomials by a kind of symmetrization procedure.
## 4. The fundamental representation
### 4.1.
In the previous section we have diagonalized the action of the “translation part” $`[Y^{\pm 1}]`$ of the affine Hecke algebra $`H=H(R;k_0,k_1)`$ under the fundamental representation $`\pi _{\underset{¯}{t},q}`$ (see 2.19). The corresponding eigenfunctions are exactly the non-symmetric Askey-Wilson polynomials. Since $`H`$ is generated as algebra by $`Y`$ and the difference-reflection operator $`T_1`$, see 2.15, it suffices to understand the action of $`T_1`$ on the non-symmetric Askey-Wilson polynomials in order to completely decompose $`𝒜`$ as an $`H`$-module. Recall the notation $`\gamma _m=k_0^{ϵ(m)}k_1^{ϵ(m)}q^m`$ ($`m`$) for the eigenvalues of $`Y`$.
###### Proposition.
For $`m`$ we have
$$T_1P_m=\alpha _mP_m+\beta _mP_m,$$
with
$$\alpha _m=\frac{(k_1^1k_1)\gamma _m^2+(k_0^1k_0)\gamma _m}{1\gamma _m^2}.$$
If $`m`$ then $`\beta _m=k_1`$, and
$$\beta _m=k_1\underset{\xi =\pm 1}{}\frac{(1+k_0k_1^1\gamma _m^\xi )(1k_0^1k_1^1\gamma _m^\xi )}{(1\gamma _m^{2\xi })}$$
if $`m_+`$.
###### Proof.
The formula for $`m=0`$ reduces to $`T_1(P_0)=k_1P_0`$, which is clear. For $`0m`$, we derive from Lusztig’s formula 2.16 and from the definition 3.5 of the non-symmetric Askey-Wilson polynomials that for all $`f(Y)[Y^{\pm 1}]`$,
$$\left(f(Y)f(\gamma _m)\right)T_1P_m=\alpha _m\left(f\left(\gamma _m\right)f\left(\gamma _m\right)\right)P_m$$
with $`\alpha _m`$ as given in the statement of the proposition. Since $`\gamma _m`$ ($`m`$) are mutually different by the conditions 2.11 on the parameters, we derive from proposition 3.5 that $`T_1P_m=\alpha _mP_m+\beta _mP_m`$ for some $`\beta _m`$.
If $`m`$, then we have $`x^ms_1(x^m)=x^m`$. Combined with the formula $`T_1=s_1(a_1)^1+k_1k_1^1`$ (see 3.1) and with the triangularity of $`(a_1)`$ (see lemma 3.3), we obtain that the coefficient of $`x^m`$ in the expansion of $`T_1(x^m)`$ with respect to the basis of monomials is equal to $`k_1`$. By the definition 3.5 of the non-symmetric Askey-Wilson polynomials, we conclude that $`\beta _m=k_1`$ for $`m`$.
Now act by $`T_1`$ on both sides of the formula $`T_1P_m=\alpha _mP_m+\beta _mP_m`$ and use the quadratic relation for $`T_1`$, see 2.12. It follows that the $`\alpha _m`$ and the $`\beta _m`$ satisfy the relation
$$\beta _m\beta _m=(k_1\alpha _m)(k_1^1+\alpha _m),0m.$$
This allows us to compute $`\beta _m`$ with $`m`$ from the known expressions for $`\alpha _m`$ and $`\beta _m`$, which yields the desired result. ∎
### 4.2.
A uniform formula for the action of $`T_1`$ on the non-symmetric Askey-Wilson polynomials can be obtained by renormalizing the non-symmetric Askey-Wilson polynomials in a suitable way. A natural renormalization, together with a new proof of proposition 4.1, is given in section 10.
### 4.3.
As a consequence of proposition 4.1, we can compute the action of the non-affine intertwiner $`S_1`$ (see 2.17) on the non-symmetric Askey-Wilson polynomials explicitly.
###### Corollary.
We have $`S_1(P_m)=\left(\gamma _m\gamma _m\right)\beta _mP_m`$ for $`m`$, where $`\beta _m`$ is as in proposition 4.1.
###### Proof.
By proposition 4.1 and by the definition 3.5 of the non-symmetric Askey-Wilson polynomial, we have
$$\begin{array}{cc}\hfill S_1P_m=& (T_1YYT_1)P_m=\left(\gamma _mY\right)T_1P_m\hfill \\ \hfill =& \left(\gamma _mY\right)\left(\alpha _mP_m+\beta _mP_m\right)=\left(\gamma _m\gamma _m\right)\beta _mP_m.\hfill \end{array}$$
### 4.4.
We set $`𝒜(0)=\text{span}\{P_0\}`$ and $`𝒜(m)=\text{span}\{P_m,P_m\}`$ for $`m`$.
###### Theorem.
(i) The representation $`(\pi _{\underset{¯}{t},q},H(R;k_0,k_1))`$ is faithful.
(ii) The center $`𝒵(H)`$ of $`H=H(R;k_0,k_1)`$ is equal to $`[Y^{\pm 1}]^W=[Y+Y^1]`$.
(iii) The decomposition $`𝒜=_{m_+}𝒜(m)`$ is the multiplicity-free, irreducible decomposition of $`𝒜`$ as $`(\pi _{\underset{¯}{t},q},H)`$-module. It is also the decomposition of $`𝒜`$ in isotypical components for the action of the center, where the central character of $`𝒜(m)`$ is given by $`\chi _m(f(Y))=f\left(\gamma _m\right)`$ for $`f(Y)𝒵(H)=[Y^{\pm 1}]^W`$.
###### Proof.
(i) Suppose that $`h=f(Y)+T_1g(Y)`$ acts as zero on $`𝒜`$ under the representation $`\pi _{\underset{¯}{t},q}`$, where $`f,g𝒜`$. Let $`h`$ act on the non-symmetric Askey-Wilson polynomials $`P_m`$ ($`m<0`$) and use proposition 4.1 together with the fact that the coefficients $`\beta _m`$ ($`m<0`$) in 4.1 are non-zero. Then we conclude that $`g(\gamma _m)=0`$ for all $`m`$. By the conditions 2.11 on the parameters, this implies $`g=0`$ in $`𝒜`$. But $`h=f(Y)`$ acting on $`P_m`$ shows that $`f\left(\gamma _m\right)=0`$ for all $`m`$, hence $`f=0`$ in $`𝒜`$. Combined with 2.15, this shows that $`\pi _{\underset{¯}{t},q}`$ is faithful.
(ii) Clearly any element from $`[Y+Y^1]`$ commutes with $`[Y^{\pm 1}]`$, but also with $`T_1`$ by Lusztig’s formula 2.16. Hence 2.15 gives $`[Y+Y^1]𝒵(H)`$. Suppose $`0h=f(Y)+T_1g(Y)𝒵(H)`$, where $`f,g𝒜`$. Then $`h`$ acts as a constant on each of the $`P_m`$ ($`m`$). In view of proposition 4.1, this implies that $`g(\gamma _m)=0`$ for all $`m0`$, hence $`g=0`$ in $`𝒜`$. By corollary 4.3 we then have for $`m0`$,
$$f\left(\gamma _m\right)S_1P_m=S_1(hP_m)=h(S_1P_m)=f\left(\gamma _m^1\right)S_1P_m.$$
Furthermore, $`S_1P_m0`$ by the conditions 2.11 on the parameters. Hence $`f(\gamma _m)=f(\gamma _m^1)`$ for $`0m`$, i.e. $`h=f(Y)[Y+Y^1]`$.
(iii) The second statement follows directly from proposition 3.5 and from the fact that the central character values $`\chi _m(Y+Y^1)=\gamma _m+\gamma _m^1`$ ($`m_+`$) are pair-wise different by the conditions 2.11 on the parameters. For the first statement, it then suffices to show that $`𝒜(m)`$ ($`m_+`$) are irreducible $`H`$-modules. This follows without difficulty from 2.15, proposition 4.1, corollary 4.3 and the fact that $`\beta _m0`$ for all $`0m`$ by the conditions 2.11 on the parameters. ∎
## 5. The (anti-)symmetric Askey-Wilson polynomials
### 5.1.
In the present rank one setting, the representation theory of the underlying two-dimensional Hecke algebra $`H_0=H_0(k_1)`$ is extremely simple: the trivial representation $`\chi _+`$ and the alternating representation $`\chi _{}`$ exhaust its irreducible representations, where $`\chi _\pm `$ are uniquely determined by $`\chi _\pm (T_1)=\pm k_1^{\pm 1}`$. The corresponding mutually orthogonal, primitive idempotents are given by
$$C_+=\frac{1}{1+k_1^2}\left(1+k_1T_1\right),C_{}=\frac{1}{1+k_1^2}\left(1k_1^1T_1\right).$$
So $`\{C_{},C_+\}`$ is a partition of the unity for $`H_0`$. In particular, we have $`C_{}+C_+=1`$.
### 5.2.
The partition of the unity of $`H_0`$ introduced in 5.1 gives the decomposition $`𝒜=𝒜_{}𝒜_+`$ of $`𝒜`$ in isotypical components for the action of $`(\pi _{\underset{¯}{t},q}|_{H_0},H_0)`$, where $`𝒜_\pm =C_\pm 𝒜`$. Observe that $`𝒜_\pm `$ consists of the Laurent polynomials $`f𝒜`$ which satisfy $`(T_1k_1^{\pm 1})f=0`$. By the explicit expression 2.19 for the difference-reflection operator $`T_1`$ we have $`T_1k_1=\varphi _1(x)(s_1\text{id})`$ for some non-zero rational function $`\varphi _1(x)`$, so that $`𝒜_+`$ coincides with the algebra $`𝒜^W=[x+x^1]`$ consisting of $`W`$-invariant Laurent polynomials. We call $`𝒜_{}`$ the subspace of anti-symmetric Laurent polynomials. The decomposition $`f=C_{}f+C_+f`$ for $`f𝒜`$ is its unique decomposition as a sum of an anti-symmetric and a symmetric Laurent polynomial.
### 5.3.
The irreducible $`H`$-module $`𝒜(m)𝒜`$ decomposes under the action of $`H_0`$ by $`𝒜(m)=𝒜_{}(m)𝒜_+(m)`$, where $`𝒜_\pm (m)=C_\pm 𝒜(m)`$.
###### Proposition.
(i) Let $`m_+`$, then $`\text{dim}(𝒜_+(m))=1`$. More precisely, there exists a unique $`P_m^+𝒜_+(m)`$ of the form $`P_m^+(x)=x^m+`$ lower order terms with respect to $``$. In terms of non-symmetric Askey-Wilson polynomials, we have
$$P_m^+=P_m+\frac{(1+k_0k_1^1\gamma _m)(1k_0^1k_1^1\gamma _m)}{(1\gamma _m^2)}P_m,m_+.$$
(ii) We have $`𝒜_{}(0)=\{0\}`$ and $`\text{dim}(𝒜_{}(m))=1`$ for $`m`$. More precisely, there exists for all $`m`$ a unique $`P_m^{}𝒜_{}(m)`$ of the form $`P_m^{}(x)=x^m+`$ lower order terms with respect to $``$. In terms of non-symmetric Askey-Wilson polynomials, we have
$$P_m^{}=P_m\frac{(1+k_0k_1^1\gamma _m^1)(1k_0^1k_1^1\gamma _m^1)}{(1\gamma _m^2)}P_m,m.$$
###### Proof.
The statements for $`m=0`$ are immediate since $`T_1(1)=k_11`$, where $`1𝒜`$ is the Laurent polynomial identically equal to one. For $`m`$, we can write $`C_\pm P_m𝒜_\pm (m)`$ explicitly as
$$\begin{array}{cc}\hfill C_\pm P_m& =\frac{1\pm k_1^{\pm 1}\alpha _m}{1+k_1^{\pm 2}}P_m\pm \frac{k_1^{\pm 1}\beta _m}{1+k_1^{\pm 2}}P_m\hfill \\ & =k_1^{\pm 1}\frac{(k_1^1\pm \alpha _m)}{(1+k_1^{\pm 2})}\left(P_m\pm k_1^1(k_1^{\pm 1}\alpha _m)P_m\right)\hfill \end{array}$$
in view of (the proof of) proposition 4.1. Observe that the coefficient of $`P_m`$ is non-zero by the conditions 2.11 on the parameters. In particular, $`𝒜_\pm (m)=\text{span}\{C_\pm P_m\}`$ are one-dimensional subspaces for all $`m`$. Dividing out the non-zero coefficient of $`P_m`$ in the expansion of $`C_\pm P_m`$ and using proposition 3.5, we conclude that there exist unique elements $`P_m^\pm 𝒜_\pm (m)`$ ($`m`$) satisfying $`P_m^\pm (x)=x^m+`$ lower order terms w.r.t. $``$. The explicit formulas for $`P_m^\pm `$ in terms of non-symmetric Askey-Wilson polynomials follow now by substituting the explicit expression for $`\alpha _m`$ in the above expansion of $`C_\pm P_m`$, see proposition 4.1. ∎
###### Definition.
(i) The polynomial $`P_m^+=P_m^+(;\underset{¯}{t};q)𝒜_+`$ ($`m_+`$) is called the monic, symmetric Askey-Wilson polynomial of degree $`m`$.
(ii) The polynomial $`P_m^{}=P_m^{}(;\underset{¯}{t};q)𝒜_{}`$ ($`m`$) is called the monic, anti-symmetric Askey-Wilson polynomial of degree $`m`$.
Askey and Wilson defined a very general family of basic hypergeometric orthogonal polynomials which are nowadays known as the Askey-Wilson polynomials. In 5.9 we justify our terminology for the Laurent polynomials $`P_m^\pm `$ by showing that the $`P_m^+`$ ($`m`$) coincide with the Askey-Wilson polynomials as defined in .
### 5.4.
We can also express the non-symmetric Askey-Wilson polynomial in terms of the symmetric Askey-Wilson polynomial in the following way.
###### Lemma.
We have $`P_0=P_0^+`$, and for $`m`$,
$$\begin{array}{cc}\hfill P_m& =\frac{1}{\gamma _m\gamma _m}\left(Y\gamma _m\right)P_m^+,\hfill \\ \hfill P_m& =\frac{\gamma _m}{(1+k_0k_1^1\gamma _m)(1k_0^1k_1^1\gamma _m)}\left(Y\gamma _m\right)P_m^+.\hfill \end{array}$$
###### Proof.
The statement for $`m=0`$ is trivial. For $`m`$ the formulas follow directly from proposition 3.5 and from the expansion of $`P_m^+`$ as linear combination of non-symmetric Askey-Wilson polynomials, see proposition 5.3. ∎
### 5.5.
Observe that the affine Weyl group $`𝒲`$ acts on $`𝒜`$ by algebra automorphisms (see 2.10). This action can be uniquely extended to an action (by automorphisms) of $`𝒲`$ on the rational functions $`(x)`$ in the indeterminate $`x`$. Since $`|q|1`$ (see 2.11), we have
$$\underset{w𝒲}{}(x)w=\underset{m,\sigma W}{}(x)\tau (m)\sigma $$
as a subalgebra of $`\text{End}_{}\left((x)\right)`$.
### 5.6.
Any $`X\text{Im}(\pi _{\underset{¯}{t},q})\text{End}_{}(𝒜)`$ can be uniquely written as a finite $`(x)`$-linear combination of the automorphisms $`w𝒲`$ of $`𝒜`$. By 5.5, we may regard $`X`$ as a linear endomorphism of $`(x)`$. We thus have a unique decomposition $`X=X_{}s_1+X_+`$ where $`X_\pm _m(x)\tau (m)`$ are $`q`$-difference operators with rational coefficients. We write $`X_{sym}:=X_{}+X_+`$, so that $`Xf=X_{sym}f`$ for all $`f𝒜_+=𝒜^W`$.
### 5.7.
In order to make the connection between the symmetric Askey-Wilson polynomials $`P_m^+`$ ($`m`$) and the four parameter family of Askey-Wilson polynomials as originally defined in , it is convenient to reparametrize the multiplicity function $`\underset{¯}{t}(u_0,u_1,k_0,k_1)`$ (see 2.11) by
$$a=k_1u_1,b=k_1u_1^1,c=q^{\frac{1}{2}}k_0u_0,d=q^{\frac{1}{2}}k_0u_0^1.$$
### 5.8.
Using the parameters $`a,b,c,d`$ (see 5.7), we can give the following explicit expression for the $`q`$-difference operator $`(Y+Y^1)_{sym}`$, see for the higher rank result.
###### Proposition.
We have
$$\left(Y+Y^1\right)_{sym}=A(x)\left(\tau (1)1\right)+A(x^1)\left(\tau (1)1\right)+k_0k_1+k_0^1k_1^1,$$
with
$$A(x)=k_0^1k_1^1\frac{(1ax)(1bx)(1cx)(1dx)}{(1x^2)(1qx^2)}.$$
###### Proof.
We write $`T_i=k_i+\varphi _i(x)(s_i1)`$ with
$$\varphi _0(x)=k_0^1\frac{(1cx^1)(1dx^1)}{(1qx^2)},\varphi _1(x)=k_1^1\frac{(1ax)(1bx)}{(1x^2)}$$
for the difference-reflection operators associated with $`S`$, see 2.19. Since $`Y=T_1T_0`$, $`s_0=t(1)s_1`$ and $`T_i^1=T_i+k_i^1k_i`$ for $`i=0,1`$, we have
$$(Y+Y^1)_{sym}=B(x)(\tau (1)1)+C(x)(\tau (1)1)+D(x)$$
for unique coefficients $`B,C,D(x)`$. Observe that $`D(x)=(Y+Y^1)(1)=k_0k_1+k_0^1k_1^1`$ where $`P_0=1A_+`$ is the Laurent polynomial identically equal to one, since $`T_i(1)=k_i1`$ for $`i=0,1`$. To compute the coefficient $`B(x)`$ (respectively $`C(x)`$), we need to compute the coefficient of $`\tau (1)`$ (respectively $`\tau (1)`$) in $`(Y+Y^1)_{sym}`$. Recall that $`s_0=s_1\tau (1)=\tau (1)s_1`$, so that the $`\tau (1)`$-term (respectively $`\tau (1)`$-term) of $`Y_{sym}`$ has coefficient $`\varphi _1(x)\varphi _0(x^1)=A(x)`$ (respectively has coefficient $`(k_1\varphi _1(x))\varphi _0(x)`$). The $`\tau (1)`$-term (respectively $`\tau (1)`$-term) of $`(Y^1)_{sym}`$ is zero (respectively has coefficient $`\varphi _0(x)\varphi _1(qx^1)+\varphi _0(x)(k_1^1\varphi _1(qx^1))=k_1^1\varphi _0(x)`$). Adding the contributions, we see that $`B(x)=A(x)`$ and that $`C(x)=\varphi _0(x)\left(k_1+k_1^1\varphi _1(x)\right)=A(x^1)`$, which completes the proof of the proposition. ∎
The second order $`q`$-difference operator $`L=(Y+Y^1)_{sym}`$ is called the Askey-Wilson second-order $`q`$-difference operator, cf. \[1, (5.7)\].
### 5.9.
The symmetric Askey-Wilson polynomial $`P_m^+`$ ($`m_+`$) lies in the irreducible $`H(R;k_0,k_1)`$-module $`𝒜(m)`$, hence the central element $`Y+Y^1𝒵(H)`$ acts on $`P_m^+`$ as the scalar $`\gamma _m+\gamma _m^1`$, see theorem 4.4. Combined with proposition 5.8, we conclude that $`P_m^+`$ is an eigenfunction of the Askey-Wilson second-order $`q`$-difference operator $`L`$ with eigenvalue $`\gamma _m+\gamma _m^1`$. The eigenvalues $`\gamma _m+\gamma _m^1`$ ($`m_+`$) are mutually different by the conditions 2.11 on the parameters and $`\{P_m^+\}_{m_+}`$ is a linear basis of $`𝒜_+=𝒜^W`$, so that $`P_m^+`$ is the unique $`W`$-invariant Laurent polynomial satisfying $`LP_m^+=(\gamma _m+\gamma _m^1)P_m^+`$. A comparison with \[1, (5.7)\] yields the following result.
###### Theorem.
The $`W`$-invariant Laurent polynomial $`P_m^+`$ ($`m_+`$) coincides with the monic Askey-Wilson polynomial of degree $`m`$ as described in . In particular, we have in terms of basic hypergeometric series,
$$P_m^+(x)=\frac{(ab,ac,ad;q)_m}{a^m(abcdq^{m1};q)_m}{}_{4}{}^{}\varphi _{3}^{}(\begin{array}{cc}& q^m,q^{m1}abcd,ax,ax^1\\ & ab,ac,ad\end{array};q,q).$$
### 5.10.
Theorem 5.9 and lemma 5.4 can be used to write the non-symmetric and the anti-symmetric Askey-Wilson polynomials as a sum of two terminating balanced $`{}_{4}{}^{}\varphi _{3}^{}`$’s. It is convenient to write $`P_m^+(x)=P_m^+(x;a,b,c,d)`$ for the symmetric Askey-Wilson polynomial $`P_m^+(x)=P_m^+(x;\underset{¯}{t};q)`$ when we want to emphasize the dependence of $`P_m^+`$ on the (reparametrized) multiplicity function $`(a,b,c,d)`$, see 5.7.
###### Proposition.
(i) For $`m_+`$ we have
$$\begin{array}{cc}& P_m(x)=q^m\frac{(1abcdq^{m1})}{(1abcdq^{2m1})}P_m^+(x;a,b,c,d)\hfill \\ & +q^{(m1)/2}\frac{(1cx^1)(1dx^1)x(1q^m)}{(1abcdq^{2m1})}P_{m1}^+(q^{1/2}x;q^{1/2}a,q^{1/2}b,q^{1/2}c,q^{1/2}d),\hfill \end{array}$$
where the second term should be read as zero when $`m=0`$.
(ii) For $`m`$ we have
$$\begin{array}{cc}\hfill P_m(x)& =\frac{1}{(1cdq^{m1})}P_m^+(x;a,b,c,d)\hfill \\ & q^{(m1)/2}\frac{(1cx^1)(1dx^1)x}{(1cdq^{m1})}P_{m1}^+(q^{1/2}x;q^{1/2}a,q^{1/2}b,q^{1/2}c,q^{1/2}d).\hfill \end{array}$$
(iii) For $`m`$ we have
$$\begin{array}{cc}& P_m^{}(x)=\frac{(1abcdq^{m1})}{ab(1cdq^{m1})}P_m^+(x;a,b,c,d)\hfill \\ & +q^{(m1)/2}\frac{(1cx^1)(1dx^1)x(ab1)}{ab(1cdq^{m1})}P_{m1}^+(q^{1/2}x;q^{1/2}a,q^{1/2}b,q^{1/2}c,q^{1/2}d).\hfill \end{array}$$
###### Proof.
Recall the rational function $`\varphi _0(x)`$ defined in the proof of proposition 5.8. The proof of proposition 5.8 shows that $`(Y^1)_{sym}=k_1^1\varphi _0(x)\left(\tau (1)1\right)+k_0^1k_1^1`$.
(i) The formula for $`m=0`$ is trivial. Let $`m`$. By lemma 5.4, we have $`P_m=(\gamma _m\gamma _m)^1(\gamma _m(Y^1)_{sym})P_m^+`$. In view of the explicit formula for the $`q`$-difference operator $`(Y^1)_{sym}`$, we need to write $`(t(1)1)P_m^+`$ as a terminating balanced $`{}_{4}{}^{}\varphi _{3}^{}`$. By a direct computation using the explicit expression of $`P_m^+`$ in terms of a terminating balanced $`{}_{4}{}^{}\varphi _{3}^{}`$, see theorem 5.9, we have
$$\begin{array}{cc}\hfill \left((\tau (1)1)P_m^+(.;a,b,c,d)\right)(x)=& (q^{1/2}x^1q^{1/2}x)(q^{m/2}q^{m/2})\hfill \\ & .P_{m1}^+(q^{1/2}x;q^{1/2}a,q^{1/2}b,q^{1/2}c,q^{1/2}d),\hfill \end{array}$$
cf. \[1, (5.6)\] or \[8, (7.7.7)\]. This leads to the desired result.
(ii) We have $`P_m(x)=(1+k_0k_1^1\gamma _m)^1(1k_0^1k_1^1\gamma _m)^1(1\gamma _m(Y^1)_{sym})P_m^+`$ by lemma 5.4. The proof is now similar to the proof of (i).
(iii) This follows from (i) and (ii), together with proposition 5.3. ∎
## 6. (Bi-)orthogonality relations
### 6.1.
We assume in this section that the multiplicity function $`\underset{¯}{t}`$ and the deformation parameter $`q`$ satisfy the additional conditions that $`0<|q|<1`$ and that $`efq^{}`$ for all $`e,f\{a,b,c,d\}`$.
### 6.2.
Let $`C`$ be a continuous rectifiable Jordan curve such that $`aq^k,bq^k,cq^k,dq^k`$ ($`k_+`$) are in the interior of $`C`$ and such that $`a^1q^k,b^1q^k,c^1q^k,d^1q^k`$ ($`k_+`$) are in the exterior of $`C`$. By the conditions 6.1 on the parameters, such a contour exists. We give $`C`$ the counterclockwise orientation. Let $`.,.=.,._{\underset{¯}{t},q}`$ and $`(.,.)=(.,.)_{\underset{¯}{t},q}`$ be the bilinear forms on $`𝒜`$ defined by
$$f,g=\frac{1}{2\pi i}_{xC}f(x)g(x^1)\mathrm{\Delta }(x)\frac{dx}{x},(f,g)=\frac{1}{2\pi i}_{xC}f(x)g(x^1)\mathrm{\Delta }_+(x)\frac{dx}{x},$$
where the weight functions $`\mathrm{\Delta }(x)=\mathrm{\Delta }(x;\underset{¯}{t};q)`$ and $`\mathrm{\Delta }_+(x)=\mathrm{\Delta }_+(x;\underset{¯}{t};q)`$ are given by the infinite products
$$\begin{array}{cc}\hfill \mathrm{\Delta }(x)& =\underset{aR^+}{}\frac{\left(1x^a\right)}{\left(1t_at_{a/2}x^{a/2}\right)\left(1+t_at_{a/2}^1x^{a/2}\right)},\hfill \\ \hfill \mathrm{\Delta }_+(x)& =\underset{aR:a(0)0}{}\frac{\left(1x^a\right)}{\left(1t_at_{a/2}x^{a/2}\right)\left(1+t_at_{a/2}^1x^{a/2}\right)}.\hfill \end{array}$$
The conditions 6.1 on the parameters ensure that the weight functions are well-defined. In terms of $`q`$-shifted factorials, we can rewrite the weight function $`\mathrm{\Delta }_+(x)`$ as
$$\mathrm{\Delta }_+(x)=\frac{(x^2,x^2;q)_{\mathrm{}}}{(ax,ax^1,bx,bx^1,cx,cx^1,dx,dx^1;q)_{\mathrm{}}}$$
using 5.7, 2.7 and 2.10. Hence $`\mathrm{\Delta }_+()`$ coincides with the weight function of the orthogonality measure of the symmetric Askey-Wilson polynomials as defined in (see also proposition 6.9). Observe that $`\mathrm{\Delta }(x)=\alpha (x)\mathrm{\Delta }_+(x)`$ with $`\alpha (x)=\alpha (x;k_1,u_1)`$ given by
$$\alpha (x)=\frac{(1k_1u_1x^1)(1+k_1u_1^1x^1)}{(1x^2)}=\frac{(1ax^1)(1bx^1)}{(1x^2)}.$$
### 6.3.
Using Cauchy’s theorem we can rewrite $`.,.`$ and $`(.,.)`$ as an integral over the unit circle $`T`$ in the complex plane plus a finite sum of residues of the integrand. The residues of the weight functions $`\mathrm{\Delta }()`$ and $`\mathrm{\Delta }_+()`$ can be computed explicitly, see \[1, section 2\] or \[8, section 7.5\] for more details.
### 6.4.
Observe that the factor $`\alpha (x)`$ in the weight function $`\mathrm{\Delta }(x)`$ satisfies the identity $`\alpha (x)+\alpha (x^1)=1ab`$. Since $`\mathrm{\Delta }_+(x)`$ is furthermore invariant under $`xx^1`$, we see that the restrictions of the bilinear forms $`.,.`$ and $`(.,.)`$ to $`𝒜^W`$ coincide up to a constant:
###### Lemma.
For $`f,g𝒜^W`$, we have $`f,g=\frac{1}{2}(1ab)(f,g)=\frac{1}{2}(1+k_1^2)(f,g)`$.
### 6.5.
Let $`T`$ be a linear endomorphism of $`𝒜`$. Then there exists at most one linear endomorphism $`T^{}`$ of $`𝒜`$ such that $`Tf,g=f,T^{}g`$ for all $`f,g𝒜`$, since the bilinear form $`.,.`$ is non-degenerate. If $`T^{}`$ exists, then we call $`T^{}`$ the adjoint of $`T`$ with respect to $`.,.`$.
### 6.6.
We write $`T_0^{},T_1^{}`$ for the difference-reflection operators associated to $`S`$ with respect to inverse parameters $`(\underset{¯}{t}^1,q^1)`$, where $`\underset{¯}{t}^1`$ is the multiplicity function $`(t_a^1)_{aS}`$. More precisely, $`T_i^{}`$ is the image of the fundamental generator $`T_iH(R;k_0^1,k_1^1)`$ under the (faithful) representation $`\pi _{\underset{¯}{t}^1,q^1}`$, see 2.19. Furthermore, we set $`Y^{}=T_1^{}T_0^{}`$ for the associated Cherednik-Dunkl operator.
###### Proposition.
The adjoint of the difference-reflection operator $`T_i`$ ($`i=0,1`$) and of the Dunkl operator $`Y`$ exists. More precisely, we have $`T_i^{}=(T_i^{})^1`$ ($`i=0,1`$) and $`Y^{}=(Y^{})^1`$.
###### Proof.
We use the notation $`T_i=k_i+\varphi _i(x)(s_i1)`$ and $`T_i^{}=k_i^1+\varphi _i^{}(x)(s_i^{}1)`$ ($`i=0,1`$) for the difference reflection operator with respect to the parameters $`(\underset{¯}{t},q)`$ and $`(\underset{¯}{t}^1,q^1)`$ respectively. Here $`s_1^{}=s_1`$, $`(s_0^{}f)(x)=f(q^1x^1)`$ and $`\varphi _i(x)`$ is as in the proof of proposition 5.8, while $`\varphi _i^{}(x)`$ is $`\varphi _i(x)`$ with the parameters $`(\underset{¯}{t},q)`$ replaced by $`(\underset{¯}{t}^1,q^1)`$. In view of the analytic dependence on the parameters $`\underset{¯}{t}`$ and $`q`$, we may assume without loss of generality that $`0<q<1`$ and that the Jordan curve $`C`$ in the definition 6.2 of $`.,.`$ satisfies the following additional properties: $`C`$ has a parametrization of the form $`r_C(x)e^{2\pi ix}`$ with $`r_C:[0,1](0,\mathrm{})`$, and $`C`$ is $`W`$-invariant: $`C^1:=\{z^1|zC\}=C`$.
For $`i=1`$ it follows now from the $`W`$-invariance of $`\mathrm{\Delta }_+(x)`$ that $`T_1f,g=f,T_1^{}g`$ for all $`f,g𝒜`$ with
$$T_1^{}=k_1\varphi _1(x^1)+\frac{\alpha (x)}{\alpha (x^1)}\varphi _1(x)s_1.$$
Now $`\alpha (x)\varphi _1(x)=\alpha (x^1)\varphi _1^{}(x)`$ and $`\varphi _1(x^1)=\varphi _1^{}(x)`$, so that
$$T_1^{}=k_1+\varphi _1^{}(x)(s_1^{}1)=(T_1^{})^1.$$
For $`i=0`$, let $`f,g𝒜`$ and set $`h(x)=f(qx^1)g(x^1)`$. Observe that
$$(T_0f)(x)g(x^1)f(x)\left((T_0^{})^1g\right)(x^1)=\varphi _0(x)\left(h(x)(s_0h)(x)\right)$$
and that
$$\varphi _0(x)\mathrm{\Delta }(x)=k_0^1\frac{(x^2,q^2x^2;q)_{\mathrm{}}}{(ax,bx,cx,dx,qax^1,qbx^1,qcx^1,qdx^1;q)_{\mathrm{}}}$$
is invariant under $`xqx^1`$. By the specific properties of $`C`$, we obtain
$$\begin{array}{cc}\hfill T_0f,gf,(T_0^{})^1g& =\frac{1}{2\pi i}_{xC}\left(h(x)(s_0h)(x)\right)\varphi _0(x)\mathrm{\Delta }(x)\frac{dx}{x}\hfill \\ & =\frac{1}{2\pi i}_{xCqC}h(x)\varphi _0(x)\mathrm{\Delta }(x)\frac{dx}{x}=0,\hfill \end{array}$$
where the last equality follows from Cauchy’s theorem since $`\varphi _0(x)\mathrm{\Delta }(x)`$ is analytic on and within $`CqC`$.
The statement for the Dunkl operator $`Y`$ is immediate since $`Y=T_1T_0`$. ∎
### 6.7.
We write $`P_m^{}`$ ($`m`$) for the non-symmetric Askey-Wilson polynomials with respect to the inverse parameters $`(\underset{¯}{t}^1,q^1)`$.
###### Proposition.
The two bases $`\{P_m\}_m`$ and $`\{P_n^{}\}_n`$ of $`𝒜`$ form a bi-orthogonal system with respect to the non-degenerate bilinear form $`.,.`$, i.e. $`P_m,P_n^{}=0`$ for $`m,n`$ if $`mn`$.
###### Proof.
By proposition 6.6 and proposition 3.5 we have
$$\gamma _mP_m,P_n^{}=YP_m,P_n^{}=P_m,(Y^{})^1P_n^{}=\gamma _nP_m,P_n^{}.$$
It follows that $`P_m,P_n^{}=0`$ if $`mn`$ since the eigenvalues $`\gamma _m`$ ($`m`$) of $`Y`$ are pair-wise different by the conditions 2.11 on the parameters. ∎
### 6.8.
We write $`P_m^+`$ and $`P_m^{}`$ for the symmetric and anti-symmetric Askey-Wilson polynomial with respect to inverse parameters $`(\underset{¯}{t}^1,q^1)`$.
Let $``$ be the basis of $`𝒜`$ consisting of $`P_m^+`$ ($`m_+`$) and $`P_n^{}`$ ($`n`$), and let $`^{}`$ be the basis of $`𝒜`$ consisting of $`P_m^+`$ ($`m_+`$) and $`P_m^{}`$ ($`m`$).
###### Proposition.
The pair $`(,^{})`$ forms a bi-orthogonal system of $`𝒜`$ with respect to the non-degenerate bilinear form $`.,.`$.
###### Proof.
It follows from proposition 6.7 and from the fact that $`P_m^\pm 𝒜(m)=\text{span}\{P_m,P_m\}`$ for $`m_+`$ (with the convention $`P_0^{}0`$), that $`P_m^\pm ,P_n^\pm =0`$ if $`mn`$.
By proposition 6.6 we have $`C_\pm f,g=f,C_\pm ^{}g`$ for all $`f,g𝒜`$, where $`C_\pm ^{}=(1+k_1^2)^1(1\pm k_1^1T_1^{})`$ are the mutually orthogonal, primitive idempotents of $`H_0(k_1^1)`$ (see 5.1) which act on $`𝒜`$ via $`\pi _{\underset{¯}{t}^1,q^1}`$. Hence $`P_m^\pm ,P_n^{}=0`$ for all $`m,n_+`$. ∎
### 6.9.
The bi-orthogonality relations of proposition 6.8 restricted to $`𝒜_+=𝒜^W`$ reduce to the well-known orthogonality relations \[1, theorem 2.3\] of the symmetric Askey-Wilson polynomials:
###### Proposition.
For all $`m_+`$, we have $`P_m^+=P_m^+`$. In particular, $`P_m^+,P_n^+=(P_m^+,P_n^+)=0`$ if $`mn`$.
###### Proof.
Recall that $`P_m^+`$ is the unique $`W`$-invariant Laurent polynomial of the form $`x^m+`$ lower order terms with respect to $``$ which is an eigenfunction of the Askey-Wilson $`q`$-difference operator $`L=(Y+Y^1)_{sym}`$ with eigenvalue $`\gamma _m+\gamma _m^1`$. Then $`P_m^+=P_m^+`$ follows from the fact that $`L`$ and the eigenvalue $`\gamma _m+\gamma _m^1`$ are invariant under $`(\underset{¯}{t},q)(\underset{¯}{t}^1,q^1)`$. The second statement follows now from proposition 6.8 and lemma 6.4. ∎
## 7. The generalized Weyl character formula
### 7.1.
The generalized Weyl character formula relates the anti-symmetric Askey-Wilson polynomial with the symmetric Askey-Wilson polynomial via the generalized Weyl denominator. The generalized Weyl denominator $`\delta ()`$, which we define in the following lemma, is an explicit anti-symmetric Laurent polynomial of minimal degree with respect to the total order $``$ on the monomials $`\{x^m\}_m`$.
###### Lemma.
We have $`𝒜_{}=\delta (z)\left(𝒜_+\right)`$, where $`\delta =\delta (;k_0,k_1)𝒜_{}`$ is given by
$$\delta (x)=x^1(xk_0^1k_1^1)(x+k_0k_1^1)=x^1(xa^1)(xb^1).$$
###### Proof.
By 2.21 we have $`(T_1+k_1^1)\delta (z)=\delta (z^1)(T_1k_1)`$. Combined with 5.2 this implies $`\delta (z)\left(𝒜_+\right)𝒜_{}`$. Let now $`f𝒜_{}`$, and set $`g=\delta ^1f(x)`$. Using the extended action of $`T_1`$ on $`(x)`$, see 5.6, we derive that
$$\delta (x^1)\left((T_1k_1)g\right)(x)=\left((T_1+k_1^1)f\right)(x)=0,$$
so that $`\left((T_1k_1)g\right)(x)=0`$. Since $`T_1k_1=\varphi _1(x)(s_11)`$ with $`0\varphi _1(x)(x)`$, we conclude that $`g`$ is $`W`$-invariant in $`(x)`$. In particular, $`\delta (x^1)f(x)=\delta (x)f(x^1)`$ in $`𝒜`$. Since $`\delta (x^1)=a^1b^1x^1(xa)(xb)`$ and $`\delta (x)`$ are relative coprime in the unique factorisation domain $`𝒜`$ by the conditions 2.11 on the parameters, we conclude that $`f`$ is divisible by $`\delta `$ in $`𝒜`$. Hence, $`g=\delta ^1f𝒜`$. Since $`g`$ is furthermore $`W`$-invariant, we conclude that $`g𝒜_+`$, hence $`𝒜_{}\delta (z)\left(𝒜_+\right)`$. ∎
### 7.2.
The bilinear form $`.,.`$ restricted to $`𝒜_{}`$ can now be related to the bilinear form $`(.,.)`$ on $`𝒜_+`$ using lemma 7.1. We identify $`\underset{¯}{t}`$ with $`(k_0,k_1,u_0,u_1)`$ in accordance with 2.11.
###### Lemma.
Assume that the parameters satisfy the additional conditions 6.1. Let $`\delta ^{}(x)=\delta (x;k_0^1,k_1^1)`$. Then for all $`f,g𝒜^W`$,
$$\delta (z)f,\delta ^{}(z)g_{\underset{¯}{t},q}=\frac{1}{2}(1+k_1^2)(f,g)_{k_0,qk_1,u_0,u_1,q}.$$
###### Proof.
Set $`\alpha ^{}(x)=\alpha (x;k_1^1,u_1^1)`$, see 6.2, then
$$\delta (x)\delta ^{}(x^1)\alpha (x)=(1ax)(1bx)(1ax^1)(1bx^1)\alpha ^{}(x).$$
By the explicit expression for the $`W`$-invariant weight function $`\mathrm{\Delta }_+(x;\underset{¯}{t};q)`$, see 6.2, we obtain
$$\delta (z)f,\delta ^{}(z)g_{\underset{¯}{t},q}=\frac{1}{2\pi i}_{xC}f(x)g(x^1)\alpha ^{}(x)\mathrm{\Delta }_+(x;k_0,qk_1,u_0,u_1;q)\frac{dx}{x}$$
for $`f,g𝒜^W`$. The result follows now by symmetrizing the integrand, cf. 6.4. ∎
### 7.3.
We are now in a position to relate the anti-symmetric Askey-Wilson polynomial $`P_m^{}`$ ($`m`$) with the symmetric Askey-Wilson polynomial $`P_{m1}^+`$ via the generalized Weyl denominator $`\delta `$. The result is as follows.
###### Proposition (Generalized Weyl character formula).
For $`m`$ we have
$$P_m^{}(x;\underset{¯}{t};q)=\delta (x;k_0,k_1)P_{m1}^+(x;k_0,qk_1,u_0,u_1;q).$$
###### Proof.
We first prove the proposition when $`|q|<1`$.
We assume for the moment that the multiplicity function $`\underset{¯}{t}`$ satisfies the additional conditions 6.1 and that $`(.,.)_{k_0,qk_1,u_0,u_1,q}`$ restricts to a non-degenerate bilinear form on $`𝒜_m^W:=\text{span}\{m_n|n=0,\mathrm{},m1\}`$, where $`m_0(x)=1`$ and $`m_n(x)=x^n+x^n`$ for $`n`$. These are generic conditions on the parameters, which can be removed by continuity at the end of the proof. Indeed, observe that the restriction of $`(.,.)_{\underset{¯}{t},q}`$ to $`𝒜_m^W`$ is non-degenerate when $`0<a,b,c,d,q<1`$, since then the bilinear form can be given as integration over the unit circle with respect to the positive weight function $`\mathrm{\Delta }_+(x)`$. By analytic continuation, it follows that the restriction of $`(.,.)_{\underset{¯}{t},q}`$ to $`𝒜_m^W`$ is non-degenerate for generic parameter values satisfying the conditions 6.1.
By proposition 6.9 we conclude that $`P_{m1}^+(x;k_0,qk_1,u_0,u_1;q)`$ is the unique $`W`$-invariant Laurent polynomial of the form $`x^{m1}+`$ lower order terms w.r.t. $``$ which is orthogonal to $`m_n`$ for $`n=0,\mathrm{},m2`$ with respect to the bilinear form $`(.,.)_{k_0,qk_1,u_0,u_1,q}`$. We show that $`p(x)=\delta (x;k_0,k_1)^1P_m^{}(x;\underset{¯}{t};q)`$ satisfies the same characterizing conditions.
By lemma 7.1 we have $`p𝒜^W`$. Since $`P_m^{}(x)=x^m+`$ lower order terms w.r.t. $``$ and $`\delta (x)=x+`$ lower order terms w.r.t $``$, we have $`p(x)=x^{m1}+`$ lower order terms w.r.t. $``$. By the triangularity properties of the anti-symmetric Askey-Wilson polynomials (see proposition 5.3(ii)) and by lemma 7.1, we see that $`\delta ^{}(z)m_n\text{span}\{P_k^{}|k=1,\mathrm{},m1\}`$ for $`n=0,\mathrm{},m2`$. By lemma 7.2 and proposition 6.8 we conclude that
$$\frac{1}{2}(1+k_1^2)(p,m_n)_{k_0,qk_1,u_0,u_1,q}=P_m^{},\delta ^{}(z)m_n_{\underset{¯}{t},q}=0,(n=0,\mathrm{},m2).$$
Hence $`p(x)=P_{m1}^+(x;k_0,qk_1,u_0,u_1;q)`$, as desired.
The proof for $`|q|>1`$ is a slight modification of the arguments for $`|q|<1`$. One uses now that $`P_m^+=P_m^+`$ for $`m_+`$ (see proposition 6.9) and a characterization of $`P_m^+(x;\underset{¯}{t};q)`$ in terms of the bilinear form $`(.,.)_{\underset{¯}{t}^1,q^1}`$. ∎
### 7.4.
We have now all the ingredients to express the diagonal terms $`P_m,P_m^{}`$ ($`m`$) and $`P_m^{},P_m^{}`$ ($`m`$) of the bi-orthogonality relations in proposition 6.7 and proposition 6.8 in terms of the “quadratic norms” $`P_m^+,P_m^+=\frac{1}{2}(1+k_1^2)(P_m^+,P_m^+)`$ ($`m_+`$).
Indeed, by the generalized Weyl character formula, lemma 7.2 and proposition 6.9 we have for $`m`$ that
$$\begin{array}{cc}& P_m^{},P_m^{}_{\underset{¯}{t},q}=\hfill \\ & \frac{1}{2}(1+k_1^2)(P_{m1}^+(;k_0,qk_1,u_0,u_1;q),P_{m1}^+(;k_0,qk_1,u_0,u_1;q))_{k_0,qk_1,u_0,u_1,q}.\hfill \end{array}$$
On the other hand, for $`m`$ we have $`P_m^{}=(1+k_1^2)C_{}P_m`$ (cf. the proof of proposition 5.3), so that
$$P_m,P_m^{}=\frac{(1+k_1^2)(1\gamma _m^2)}{(1+k_0k_1^1\gamma _m^1)(1k_0^1k_1^1\gamma _m^1)}P_m^{},P_m^{}$$
by proposition 5.3(ii). Here we have used that $`C_+f,g=f,C_{}^{}g`$ for all $`f,g𝒜`$ (see the proof of proposition 6.8) and that $`C_{}P_m^{}=P_m^{}`$. Similarly, we can relate $`P_m,P_m^{}`$ for $`m`$ to $`P_m^{},P_m^{}`$ using (the proof of) proposition 5.3(ii).
### 7.5.
The generalized Weyl character formula plays a crucial role in the study of shift operators for the symmetric Askey-Wilson polynomials. In turn, shift operators can be used to explicitly evaluate the quadratic norms $`(P_m^+,P_m^+)_{\underset{¯}{t},q}`$ ($`m_+`$). Combined with 7.4, this leads to explicit evaluations of all the diagonal terms of the bi-orthogonality relations in proposition 6.7 and proposition 6.8.
In section 11 we present another method for deriving explicit expressions of the diagonal terms, which uses the double affine Hecke algebra in an essential way. This method gives more insight in the particular structure of the diagonal terms. Namely, it shows that the diagonal terms can be naturally expressed in terms of the residue of the weight function in a certain specific simple pole, the constant term $`1,1`$, and the value of the Askey-Wilson polynomial at the point $`a^1=k_1^1u_1^1`$.
We return to shift operators in section 12 in order to evaluate the constant term $`1,1`$ (which is the well-known Askey-Wilson integral, see ).
## 8. The double affine Hecke algebra
### 8.1.
Recall from 2.21 that the double affine Hecke algebra $`(S;\underset{¯}{t};q)`$ is the subalgebra of $`\text{End}_{}\left(𝒜\right)`$ generated by $`\pi _{\underset{¯}{t},q}\left(H(R;k_0,k_1)\right)`$ and by $`𝒜`$, where $`𝒜`$ is regarded as subalgebra of $`\text{End}_{}(𝒜)`$ via its regular representation.
### 8.2.
We have observed in 2.22 that there is a unique surjective algebra homomorphism $`\varphi :(\underset{¯}{t};q)(S;\underset{¯}{t};q)`$ satisfying the conditions as stated in theorem 2.22. In particular, we have
$$z^{a_0^{}}T_0=T_0^1z^{a_0^{}}+u_0^1u_0,z^{a_1^{}}T_1^1=T_1z^{a_1^{}}+u_1u_1^1$$
in $`(S;\underset{¯}{t};q)`$. Combined with 2.15, we conclude that $`(S;\underset{¯}{t};q)`$ is spanned by the elements $`z^mY^n`$ and $`z^mT_1Y^n`$ where $`m,n`$. In fact, we have the following stronger result, see Sahi \[19, theorem 3.2\] for the higher rank setting.
###### Proposition.
The set $`\{z^mY^n,z^mT_1Y^n|m,n\}`$ is a linear basis of $`(S;\underset{¯}{t};q)`$.
###### Proof.
Assume that $`X=_{m,n}\left(c_{m,n}^1z^mY^n+c_{m,n}^2z^mT_1Y^n\right)=0`$ in $`\text{End}_{}\left(𝒜\right)`$ with only finitely many coefficients $`c_{m,n}^j`$ non-zero. Since $`z`$ is invertible in $`\text{End}_{}\left(𝒜\right)`$, we may assume without loss of generality that $`c_{m,n}^j=0`$ unless $`m_+`$. Suppose that not all coefficients $`c_{m,n}^j`$ are zero. Let $`m_0_+`$ be the largest positive integer such that $`c_{m_0,n}^j`$ is non-zero for some $`n`$ and some $`j\{1,2\}`$.
Let $`X`$ act on the non-symmetric Askey-Wilson polynomial $`P_l`$ ($`l`$), and consider the coefficient of $`x^{m_0+l}`$ in the resulting expression using proposition 3.5 and proposition 4.1. We obtain $`_{n_+}k_1c_{m_0,n}^2\gamma _l^n=0`$ for all $`l`$, hence $`c_{m_0,n}^2=0`$ for all $`n`$.
Let now $`X`$ act on $`P_l`$ with $`l_+`$, and again consider the coefficient of $`x^{m_0+l}`$ in the resulting expression. Then $`_{n_+}c_{m_0,n}^1\gamma _l^n=0`$ for all $`l_+`$, hence $`c_{m_0,n}^1=0`$ for all $`n`$. This gives the desired contradiction. ∎
### 8.3.
We can finish now the proof of theorem 2.22 using proposition 8.2. It suffices to show that the surjective algebra homomorphism $`\varphi :(\underset{¯}{t};q)(S;\underset{¯}{t};q)`$ defined in theorem 2.22 is injective. Set $`w:=V_1^1(V_1^{})^1=q^{1/2}V_0V_0^{}(\underset{¯}{t};q)`$, then $`\varphi (w)=z`$. Observe that $`V_0^{}=q^{1/2}w^1V_0+u_0u_0^1`$ and that $`V_1^{}=w^1V_1^1`$ in $`(\underset{¯}{t};q)`$, so that $`(\underset{¯}{t};q)`$ is generated by $`w^{\pm 1},V_0`$ and $`V_1`$ as an algebra. Furthermore,
$$V_0w=qw^1V_0+(k_0k_0^1)w+q^{1/2}(u_0u_0^1),V_1w=w^1V_1^1+u_1^1u_1$$
in $`(\underset{¯}{t};q)`$, so that any element in $`(\underset{¯}{t};q)`$ can be written as a finite linear combination of elements of the form $`f(w)X`$, where $`f(w)`$ is a Laurent polynomial in $`w`$ and $`X`$ is an element in the subalgebra of $`(\underset{¯}{t};q)`$ generated by $`V_0`$ and $`V_1`$. By 2.15 and by the relations (1) in theorem 2.22 for the generators $`V_0,V_1(\underset{¯}{t};q)`$ it follows that $`(\underset{¯}{t};q)`$ is spanned by $`\{w^mZ^n,w^mV_1Z^n|m,n\}`$, where $`Z:=V_1V_0`$. Since the image of these elements under $`\varphi `$ are linear independent by proposition 8.2, we conclude that $`\varphi `$ is injective.
### 8.4.
In the remainder of the paper we use the notations $`T_0^{}:=T_0^1z^{a_0^{}}(S;\underset{¯}{t},q)`$ and $`T_1^{}:=z^{a_1^{}}T_1^1(S;\underset{¯}{t};q)`$ for the images of $`V_0^{}`$ and $`V_1^{}`$ respectively under the algebra isomorphism $`\varphi :(\underset{¯}{t};q)(S;\underset{¯}{t};q)`$ (see theorem 2.22).
### 8.5.
We associate with the multiplicity function $`\underset{¯}{t}(k_0,k_1,u_0,u_1)`$ a dual multiplicity function $`\underset{¯}{\overset{~}{t}}`$ by interchanging $`k_0`$ and $`u_1`$, so $`\underset{¯}{\overset{~}{t}}(u_1,k_1,u_0,k_0)`$. We write $`\stackrel{~}{T}_0,\stackrel{~}{T}_1,\stackrel{~}{T}_0^{},\stackrel{~}{T}_1^{}`$ for the generators of $`(S;\underset{¯}{\overset{~}{t}};q)`$ (cf. 2.22 and 8.4), and we write $`\stackrel{~}{Y}=\stackrel{~}{T}_1\stackrel{~}{T}_0`$ for the associated Dunkl operator and $`\stackrel{~}{z}=\stackrel{~}{T}_1^1(\stackrel{~}{T}_1^{})^1=q^{1/2}\stackrel{~}{T}_0\stackrel{~}{T}_0^{}`$ for the corresponding “multiplication by $`x`$” operator. The first part of the following proposition is a special case of Sahi’s results in \[19, section 7\].
###### Proposition.
(i) The application $`T_0\stackrel{~}{T}_1^{}`$, $`T_1\stackrel{~}{T}_1`$, $`T_0^{}\stackrel{~}{T}_0^{}`$ and $`T_1^{}\stackrel{~}{T}_0`$ uniquely extend to an anti-algebra isomorphism $`\nu =\nu _{\underset{¯}{t},q}:(S;\underset{¯}{t};q)(S;\underset{¯}{\overset{~}{t}};q)`$. Furthermore, $`\nu _{\underset{¯}{t},q}^1=\nu _{\underset{¯}{\overset{~}{t}},q}`$ and $`\nu (z)=\stackrel{~}{Y}^1`$, $`\nu (Y)=\stackrel{~}{z}^1`$.
(ii) The application $`T_0\stackrel{~}{T}_1^1\stackrel{~}{T}_1^{}\stackrel{~}{T}_1`$, $`T_1\stackrel{~}{T}_1`$, $`T_0^{}\stackrel{~}{T}_0\stackrel{~}{T}_0^{}\stackrel{~}{T}_0^1`$ and $`T_1^{}\stackrel{~}{T}_0`$ uniquely extend to an algebra isomorphism $`\mu =\mu _{\underset{¯}{t},q}:(S;\underset{¯}{t};q)(S;\underset{¯}{\overset{~}{t}};q)`$. Furthermore, $`\mu (Y)=\stackrel{~}{z}^1`$.
###### Proof.
By theorem 2.22 it suffices to check that $`\mu `$ (respectively $`\nu `$) is compatible with the defining relations in $`(\underset{¯}{t};q)(S;\underset{¯}{t};q)`$. This can be done by direct computations. It is immediate that $`\nu _{\underset{¯}{\overset{~}{t}},q}`$ is the inverse of $`\nu _{\underset{¯}{t},q}`$.
Observe that the application $`\stackrel{~}{T}_0T_1^{}`$, $`\stackrel{~}{T}_1T_1`$, $`\stackrel{~}{T}_0^{}(T_1^{})^1T_0^{}T_1^{}`$ and $`\stackrel{~}{T}_1^{}T_1T_0T_1^1`$ uniquely extend to an algebra homomorphism from $`(S;\underset{¯}{\overset{~}{t}};q)`$ to $`(S;\underset{¯}{t};q)`$. It is immediate that this homomorphism is the inverse of $`\mu `$. ∎
### 8.6.
Following the terminology of Sahi \[19, section 7\], we call $`\nu =\nu _{\underset{¯}{t},q}`$ the duality anti-isomorphism. Furthermore, we call $`\mu =\mu _{\underset{¯}{t},q}`$ the duality isomorphism. These duality isomorphisms play a fundamental role in the theory of non-symmetric Askey-Wilson polynomials. In particular, the duality anti-isomorphism can be used to show that the geometric parameter $`x`$ and the spectral parameter $`\gamma `$ of the non-symmetric Askey-Wilson polynomial are in a sense interchangeable (see Sahi \[19, section 7\] or section 10). The duality isomorphism describes the intertwining properties of the action of the double affine Hecke algebra under the non-symmetric Askey-Wilson transform, see section 11.
### 8.7.
We write $`T_i^{}`$ and $`T_i^{}`$ ($`i=0,1`$) for the generators of $`(S;\underset{¯}{t}^1;q^1)`$, cf. 2.22, 6.6 and 8.4.
###### Proposition.
There exists a unique anti-algebra isomorphism $`{}_{}{}^{}:(S;\underset{¯}{t};q)(S;\underset{¯}{t}^1;q^1)`$ such that $`T_i^{}=(T_i^{})^1`$ and $`(T_i^{})^{}=(T_i^{})^1`$ for $`i=0,1`$. Furthermore, $`T^{}`$ coincides with the adjoint of $`T(S;\underset{¯}{t};q)`$ if the parameters satisfy the additional conditions 6.1.
###### Proof.
The first statement follows easily from theorem 2.22. For the second statement, it suffices to compute the adjoint of $`T_i^{}`$ ($`i=0,1`$) in view of proposition 6.6. Let $`z^{}=q^{1/2}T_0^{}T_0^{}`$ be the “multiplication by $`x`$” operator in $`(S;\underset{¯}{t}^1;q^1)`$. It is immediate that $`z^{}=(z^{})^1`$. Combined with proposition 6.6 we obtain $`(T_0^{})^{}=q^{1/2}(z^{})^1T_0^{}=(T_0^{})^1`$ and $`(T_1^{})^{}=T_1^{}z^{}=(T_1^{})^1`$. This gives the desired result. ∎
## 9. Intertwiners as creation operators
### 9.1.
In corollary 2.17 we have introduced the non-affine intertwiner $`S_1`$ and derived its basic property. The results of the previous section allow us to derive the following analogous result for the commutator $`[Y,T_1^{}](S;\underset{¯}{t};q)`$, see Sahi \[19, theorem 5.1\] for the result in the higher rank setting.
###### Corollary (The affine intertwiner).
Set $`S_0:=[Y,T_1^{}]=YT_1^{}T_1^{}Y(S;\underset{¯}{t};q)`$. Then $`g(Y)S_0=S_0g(q^1Y^1)`$ for all $`g(Y)[Y^{\pm 1}]`$.
###### Proof.
By 2.21 we have $`f(z)[T_0,z^1]=[T_0,z^1](s_0f)(z)`$ in $`(S;\underset{¯}{t};q)`$ for any Laurent polynomial $`f`$. Apply now the duality anti-isomorphism $`\nu _{\underset{¯}{t},q}`$ to this equality, and replace the parameters by dual parameters in the resulting identity. This gives $`S_0f(Y^1)=(s_0f)(Y^1)S_0`$ for any Laurent polynomial $`f`$. The corollary is now immediate. ∎
### 9.2.
Recall from corollary 4.3 that the action of $`S_1`$ on the non-symmetric Askey-Wilson polynomials is completely explicit. In the following lemma we give the analogous result for the action of the affine intertwiner $`S_0`$ on $`P_m`$ ($`m_+`$).
###### Lemma.
Let $`m_+`$, then $`S_0(P_m)=\left(\gamma _{m1}\gamma _m\right)k_1^1P_{m1}`$.
###### Proof.
Let $`m_+`$. By proposition 4.1 we have
$$S_0(P_m)=(Y\gamma _m)\left((\alpha _m+k_1^1k_1)z^1P_m+\beta _mz^1P_m\right).$$
It follows then from proposition 3.4 that the leading term of $`S_0(P_m)`$ with respect to the total order $``$ on the monomials equals
$$(\gamma _{m1}\gamma _m)\left((\alpha _m+k_1^1k_1)c_m+\beta _m\right)x^{m1},$$
where $`c_0:=1`$ and where $`c_m`$ ($`m`$) is the unique constant such that $`P_m(x)=x^m+c_mx^m+`$ lower order terms w.r.t $``$. By proposition 5.3(i) we have
$$c_m=1\frac{(1+k_0k_1^1\gamma _m)(1k_0^1k_1^1\gamma _m)}{(1\gamma _m^2)}=k_1^1\alpha _m.$$
Furthermore, recall from the proof of proposition 4.1 that $`\beta _m=k_1^1(k_1\alpha _m)(k_1^1+\alpha _m)`$. Hence the leading term of $`S_0(P_m)`$ reduces to $`(\gamma _{m1}\gamma _m)k_1^1x^{m1}`$. On the other hand, corollary 9.1 implies that $`S_0(P_m)=d_mP_{m1}`$ for some constant $`d_m`$. By the leading term considerations, we conclude that $`d_m=(\gamma _{m1}\gamma _m)k_1^1`$. ∎
### 9.3.
The intertwiners $`S_0`$ and $`S_1`$ can be used to create the non-symmetric Askey-Wilson polynomial $`P_m`$ ($`m`$) from the unit polynomial $`1𝒜`$ in the following way.
###### Proposition.
We have $`(S_1S_0)^m(1)=d_mP_m`$ for $`m_+`$ and $`\left(S_0(S_1S_0)^{m1}\right)(1)=d_mP_m`$ for $`m`$, with the constants $`d_m`$ ($`m`$) given by
$$\begin{array}{cc}\hfill d_m& =q^{(m+1)m}k_0^{2m}k_1^{2m}(qk_0^2k_1^2;q)_{2m},m_+,\hfill \\ \hfill d_m& =q^{m^2}k_0^{12m}k_1^{2m}(qk_0^2k_1^2;q)_{2m1},m.\hfill \end{array}$$
###### Proof.
By corollary 2.17 and corollary 9.1 we have $`g(Y)(S_1S_0)=(S_1S_0)g(qY)`$ for all $`g𝒜`$. It follows that $`F_m:=(S_1S_0)^m(1)𝒜`$ for $`m_+`$ satisfies $`g(Y)F_m=g(\gamma _m)F_m`$ for all $`m`$, so $`F_m=d_mP_m`$ for some constant $`d_m`$. Similarly, we obtain $`F_m:=S_0(S_1S_0)^{m1}(1)=d_mP_m`$ for some constant $`d_m`$ when $`m`$. By corollary 4.3 and lemma 9.2, we have the recurrence relations
$$\begin{array}{cc}\hfill d_m& =(\gamma _m\gamma _m)k_1d_m,m,\hfill \\ \hfill d_{m1}& =(\gamma _{m1}\gamma _m)k_1^1d_m,m_+.\hfill \end{array}$$
Together with the initial condition $`d_0=1`$, we obtain the explicit expressions for $`d_m`$ ($`m`$) by complete induction with respect to $`m`$. ∎
## 10. Evaluation formula and duality
### 10.1.
Let $`\text{Ev}=\text{Ev}_{\underset{¯}{t},q}:(S;\underset{¯}{t};q)`$ be the linear map defined by $`\text{Ev}(X):=\left(X(1)\right)(k_1^1u_1^1)`$, where $`1𝒜`$ is the Laurent polynomial identically equal to one. Observe that Ev satisfies
$$\text{Ev}\left(T_1^{\pm 1}X\right)=k_1^{\pm 1}\text{Ev}(X),X(S;\underset{¯}{t};q),$$
since $`(T_1f)(k_1^1u_1^1)=k_1f(k_1^1u_1^1)`$ for all $`f𝒜`$ by the explicit expression 2.19 for the difference-reflection operator $`T_1`$.
### 10.2.
Observe that we can evaluate $`\text{Ev}(P_m(z))=P_m(k_1^1u_1^1)`$ explicitly using proposition 5.10, since the two $`{}_{4}{}^{}\varphi _{3}^{}`$’s in the right hand side of the formula for $`P_m(x)`$ are equal to one when $`x=a^1=k_1^1u_1^1`$. We give here an alternative, inductive proof for the evaluation which only uses the Rodrigues type formula for the non-symmetric Askey-Wilson polynomials in terms of the intertwiners $`S_0`$ and $`S_1`$, see proposition 9.3. We abuse notation by writing $`\text{Ev}(f)=\text{Ev}(f(z))=f(k_1^1u_1^1)`$ for $`f𝒜`$.
###### Proposition.
(i) For $`m_+`$, we have
$$\begin{array}{cc}\hfill \text{Ev}\left(P_m\right)& =k_1^mu_1^m\frac{(qk_1^2,q^{1/2}k_0k_1u_0u_1,q^{1/2}k_0k_1u_0^1u_1;q)_m}{(q^{m+1}k_0^2k_1^2;q)_m}\hfill \\ & =a^m\frac{(qab,ac,ad;q)_m}{(q^mabcd;q)_m}.\hfill \end{array}$$
(ii) For $`m`$, we have
$$\begin{array}{cc}\hfill \text{Ev}\left(P_m\right)& =\frac{k_1^mu_1^m}{1+k_1^2}\frac{(k_1^2,q^{1/2}k_0k_1u_0u_1,q^{1/2}k_0k_1u_0^1u_1;q)_m}{(q^mk_0^2k_1^2;q)_m}\hfill \\ & =\frac{a^m}{(1a^1b^1)}\frac{(ab,ac,ad;q)_m}{(q^{m1}abcd;q)_m}.\hfill \end{array}$$
###### Proof.
We write $`F_m=(S_1S_0)^m(1)`$ for $`m_+`$ and $`F_m=(S_0(S_1S_0)^{m1})(1)`$ for $`m`$, so that $`F_m=d_mP_m`$ for $`m`$ with the specific constants $`d_m`$ as given in proposition 9.3. The proposition follows then from the explicit evaluation of the $`d_m`$, see proposition 9.3, and from the recurrence relations
$$\text{Ev}(F_m)=k_1^1\gamma _m(1k_0k_1\gamma _m)(1+k_0^1k_1\gamma _m)\text{Ev}(F_m),m,$$
respectively
$$\text{Ev}(F_m)=u_1^1\gamma _m(1u_0u_1q^{1/2}\gamma _m)(1+u_0^1u_1q^{1/2}\gamma _m)\text{Ev}(F_{m1}),m,$$
by complete induction with respect to $`m`$. Let $`m`$. For the first recurrence relation, observe that by formula 10.1 and by $`YF_m=\gamma _mF_m`$ we have
$$\text{Ev}(F_m)=\text{Ev}(S_1F_m)=\text{Ev}\left((k_1\gamma _mk_1T_0T_1)F_m\right).$$
To reduce the $`T_0T_1`$-term, we use the relation
$$T_0T_1=Y^1+(k_0k_0^1)T_1+(k_1k_1^1)T_1^1Y(k_0k_0^1)(k_1k_1^1)$$
in $``$ and formula 10.1, which yields
$$\text{Ev}(F_m)=k_1^1\gamma _m(1k_0k_1\gamma _m)(1+k_0^1k_1\gamma _m)\text{Ev}(F_m)$$
after a direct computation. For the second recurrence relation, observe that
$$\text{Ev}(F_m)=\text{Ev}(S_0F_{m1})=\text{Ev}\left((Yz^1T_1^1u_1\gamma _{m1})F_{m1}\right)$$
by formula 10.1, since $`YF_{m1}=\gamma _{m1}F_{m1}`$. To reduce the $`Yz^1T_1^1`$-term, we use the relation
$$Yz^1T_1^1=q^1z^1T_1^1Y^1+q^1(u_1^1u_1)Y^1+q^{1/2}(u_0^1u_0)$$
in $``$ and formula 10.1, which gives the desired recursion
$$\text{Ev}(F_m)=u_1^1\gamma _m(1u_0u_1q^{1/2}\gamma _m)(1+u_0^1u_1q^{1/2}\gamma _m)\text{Ev}(F_{m1})$$
after a direct computation. This completes the proof of the proposition. ∎
### 10.3.
The explicit expression for the (anti-)symmetric Askey-Wilson polynomial as linear combination of non-symmetric Askey-Wilson polynomials (see proposition 5.3) can be used to express $`\text{Ev}(P_m^\pm )`$ as a linear combination of $`\text{Ev}(P_m)`$ and $`\text{Ev}(P_m)`$. Combined with proposition 10.2, this leads to explicit evaluation formulas for the (anti-)symmetric Askey-Wilson polynomials $`P_m^\pm `$. In particular, it follows that
$$\text{Ev}(P_m^+)=P_m^+(a)=a^m\frac{(ab,ac,ad;q)_m}{(q^{m1}abcd;q)_m},m_+.$$
This result can also be obtained directly from the explicit expression of $`P_m^+`$ in terms of a terminating, balanced $`{}_{4}{}^{}\varphi _{3}^{}`$, see theorem 5.9.
### 10.4.
The evaluation mapping Ev and the duality anti-isomorphism $`\nu `$ are compatible in the following sense.
###### Lemma.
For all $`X(S;\underset{¯}{t};q)`$ we have $`\text{Ev}_{\underset{¯}{\overset{~}{t}},q}\left(\nu _{\underset{¯}{t},q}(X)\right)=\text{Ev}_{\underset{¯}{t},q}\left(X\right)`$.
###### Proof.
For $`X=f(z)g(Y)`$ with $`f`$ and $`g`$ Laurent polynomials, we have
$$\text{Ev}_{\underset{¯}{\overset{~}{t}},q}(\nu _{\underset{¯}{t},q}(X))=\left(g(\stackrel{~}{z}^1)f(\stackrel{~}{Y}^1)(1)\right)(k_1^1k_0^1)=f(k_1^1u_1^1)g(k_1k_0)=\text{Ev}_{\underset{¯}{t},q}(X),$$
and for $`X=f(z)T_1g(Y)`$ we have
$$\text{Ev}_{\underset{¯}{\overset{~}{t}},q}(\nu _{\underset{¯}{t},q}(X))=f(k_1^1u_1^1)k_1g(k_1k_0)=\text{Ev}_{\underset{¯}{t},q}(X).$$
Combined with proposition 8.2 we obtain the desired result. ∎
### 10.5.
We associate with the evaluation mappings $`\text{Ev}_{\underset{¯}{t},q}`$ and $`\text{Ev}_{\underset{¯}{\overset{~}{t}},q}`$ two bilinear forms
$$B:(S;\underset{¯}{t};q)\times (S;\underset{¯}{\overset{~}{t}};q),\stackrel{~}{B}:(S;\underset{¯}{\overset{~}{t}};q)\times (S;\underset{¯}{t};q),$$
which are defined by $`B(X,\stackrel{~}{X})=\text{Ev}_{\underset{¯}{t},q}\left(\nu _{\underset{¯}{\overset{~}{t}},q}(\stackrel{~}{X})X\right)`$ and $`\stackrel{~}{B}(\stackrel{~}{X},X)=\text{Ev}_{\underset{¯}{\overset{~}{t}},q}\left(\nu _{\underset{¯}{t},q}(X)\stackrel{~}{X}\right)`$ for $`X(S;\underset{¯}{t};q)`$ and $`\stackrel{~}{X}(S;\underset{¯}{\overset{~}{t}};q)`$.
###### Lemma.
Let $`X,X_1,X_2(S;\underset{¯}{t};q)`$ and $`\stackrel{~}{X},\stackrel{~}{X}_1,\stackrel{~}{X}_2(S;\underset{¯}{\overset{~}{t}};q)`$. Let $`f𝒜`$.
(i) $`B(X,\stackrel{~}{X})=\stackrel{~}{B}(\stackrel{~}{X},X)`$.
(ii) $`B(X_1X_2,\stackrel{~}{X})=B(X_2,\nu _{\underset{¯}{t},q}(X_1)\stackrel{~}{X})`$, and $`B(X,\stackrel{~}{X}_1\stackrel{~}{X}_2)=B(\nu _{\underset{¯}{\overset{~}{t}},q}(\stackrel{~}{X}_1)X,\stackrel{~}{X}_2)`$.
(iii) $`B((X(f))(z),\stackrel{~}{X})=B(X.f(z),\stackrel{~}{X})`$ and $`B(X,(\stackrel{~}{X}(f))(\stackrel{~}{z}))=B(X,\stackrel{~}{X}.f(\stackrel{~}{z}))`$.
(iv) $`B(XT_i,\stackrel{~}{X})=k_iB(X,\stackrel{~}{X})`$ for $`i=0,1`$.
###### Proof.
(i) This follows from lemma 10.4 and the fact that $`\nu _{\underset{¯}{t},q}`$ is an anti-algebra homomorphism with inverse $`\nu _{\underset{¯}{\overset{~}{t}},q}`$, see proposition 8.5.
(ii) This is an immediate consequence of proposition 8.5.
(iii) The first equality is a direct consequence of the identity $`(X(f))(z)(1)=X(f)=(X.f(z))(1)`$ in $`𝒜`$. The second identity follows from the first and from (i).
(iv) By the explicit expressions 2.19 for the difference-reflection operators $`T_i`$, we have $`T_i(1)=k_i1`$ for $`i=0`$ and $`i=1`$. The identities are now immediate. ∎
### 10.6.
We write $`x_m=k_1^{ϵ(m)}u_1^{ϵ(m)}q^m`$ ($`m`$) for the eigenvalues of the Cherednik-Dunkl operator $`\stackrel{~}{Y}(S;\underset{¯}{\overset{~}{t}};q)`$, see 3.5. We assume from now on that the parameters $`(\underset{¯}{t},q)`$ are such that $`P_m(x_0^1;\underset{¯}{t};q)=\text{Ev}_{\underset{¯}{t},q}(P_m(;\underset{¯}{t};q))0`$ and such that $`P_m(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)=\text{Ev}_{\underset{¯}{\overset{~}{t}},q}(P_m(;\underset{¯}{\overset{~}{t}};q))0`$ for all $`m`$, and similarly for $`P_m^+`$ ($`m_+`$). By proposition 10.2 and 10.3, the corresponding generic conditions on the parameters can be specified explicitly. We write $`s(\gamma ):=\gamma +\gamma ^1`$ for all $`\gamma ^{}`$.
###### Definition.
(i) The renormalized non-symmetric Askey-Wilson polynomials are defined by
$$E_{\gamma _m}(x;\underset{¯}{t};q):=\frac{P_m(x;\underset{¯}{t};q)}{P_m(x_0^1;\underset{¯}{t};q)},m.$$
In other words, the non-symmetric Askey-Wilson are normalized such that they take the value one at $`x=x_0^1`$.
(ii) The renormalized symmetric Askey-Wilson polynomials are defined by
$$E_{s(\gamma _m)}^+(x;\underset{¯}{t};q):=\frac{P_m^+(x;\underset{¯}{t};q)}{P_m^+(x_0;\underset{¯}{t};q)},m_+.$$
In other words, the symmetric Askey-Wilson polynomials are normalized such that they take the value one at $`x=x_0^{\pm 1}`$.
Observe that $`C_+E_{\gamma _m}=E_{s(\gamma _m)}^+`$ for $`m`$, where $`C_+=(1+k_1^2)^1(1+k_1T_1)`$ is the idempotent defined in 5.1, since $`(T_1f)(k_1^1u_1^1)=k_1`$ for all $`f𝒜`$.
### 10.7.
In the following theorem we prove the duality between the geometric parameter $`x=x_n`$ and the spectral parameter $`\gamma =\gamma _m`$ for the renormalized (non-)symmetric Askey-Wilson polynomials, see Sahi \[19, section 7\] for the result in the higher rank setting.
###### Theorem (Duality).
(i) For all $`m,n`$ and $`f𝒜`$, we have
$$f(\gamma _m^1)=\stackrel{~}{B}(f(\stackrel{~}{z}),E_{\gamma _m}(z;\underset{¯}{t};q)),f(x_n^1)=B(f(z),E_{x_n}(\stackrel{~}{z};\underset{¯}{\overset{~}{t}};q)).$$
In particular, $`E_{\gamma _m}(x_n^1;\underset{¯}{t};q)=E_{x_n}(\gamma _m^1;\underset{¯}{\overset{~}{t}};q)`$ for all $`m,n`$.
(ii) For all $`m,n`$ and $`f𝒜^W`$, we have
$$f(\gamma _m)=\stackrel{~}{B}(f(\stackrel{~}{z}),E_{s(\gamma _m)}^+(z;\underset{¯}{t};q)),f(x_n)=B(f(z),E_{s(x_n)}^+(\stackrel{~}{z};\underset{¯}{\overset{~}{t}};q)).$$
In particular, $`E_{s(\gamma _m)}^+(x_n;\underset{¯}{t};q)=E_{s(x_n)}^+(\gamma _m;\underset{¯}{\overset{~}{t}};q)`$ for all $`m,n`$.
###### Proof.
(i) The second statement follows from the first by taking $`f=E_{x_n}(;\underset{¯}{\overset{~}{t}};q)`$ in the first equality and $`f=E_{\gamma _m}(;\underset{¯}{t};q)`$ in the second equality and using lemma 10.5(i).
For the first equality, observe that
$$\begin{array}{cc}\hfill \stackrel{~}{B}(f(\stackrel{~}{z}),E_{\gamma _m}(z))& =\stackrel{~}{B}(1,f(Y^1)E_{\gamma _m}(z))=\stackrel{~}{B}(1,(f(Y^1)E_{\gamma _m})(z))\hfill \\ & =f(\gamma _m^1)\stackrel{~}{B}(1,E_{\gamma _m}(z))=f(\gamma _m^1)E_{\gamma _m}(x_0^1)=f(\gamma _m^1)\hfill \end{array}$$
by application of lemma 10.5(i), (ii) and (iii). The second equality is proved in a similar manner.
(ii) The proof is similar to the proof of (i), taking account of the fact that $`f(Y)E_{s(\gamma )}^+=f(\gamma )E_{s(\gamma )}^+`$ for all $`f(Y)[Y+Y^1]`$ by theorem 4.4 and proposition 5.3. ∎
### 10.8.
We write $`\sigma =\{\gamma _m\}_m`$ for the spectrum of the Dunkl operator $`Y(S;\underset{¯}{t};q)`$, see 3.5. We define an action of the affine Weyl group $`𝒲`$ on $`\sigma `$ by $`s_0(\gamma _m)=\gamma _{m1}`$ and $`s_1(\gamma _m)=\gamma _m`$ for all $`m`$.
The duality between the geometric and spectral parameter of the renormalized non-symmetric Askey-Wilson polynomials can be used to explicitly compute $`X(E_\gamma )`$ for $`X(S;\underset{¯}{t};q)`$ as linear combination of non-symmetric Askey-Wilson polynomials, cf. 4.1 and 4.2 for $`X=T_1`$. Observe that $`Y(E_\gamma )=\gamma E_\gamma `$ and that $`T_0^{}=q^{1/2}Y^1(T_1^{})^1`$ by theorem 2.22 and 8.4. Hence it suffices to expand $`T_1(E_\gamma )`$ and $`T_1^{}(E_\gamma )`$ as linear combination of renormalized non-symmetric Askey-Wilson polynomials. The result is as follows.
###### Proposition.
Let $`\gamma \sigma `$, then
$$\begin{array}{cc}\hfill T_1(E_\gamma )& =k_1E_\gamma +k_1^1\frac{(1k_0k_1\gamma ^1)(1+k_0^1k_1\gamma ^1)}{(1\gamma ^2)}\left(E_{s_1\gamma }E_\gamma \right),\hfill \\ \hfill T_1^{}(E_\gamma )& =u_1E_\gamma +u_1^1\frac{(1u_0u_1q^{1/2}\gamma )(1+u_0^1u_1q^{1/2}\gamma )}{(1q\gamma ^2)}\left(E_{s_0\gamma }E_\gamma \right).\hfill \end{array}$$
###### Proof.
The first formula is obviously correct for $`\gamma =\gamma _0`$. Let $`\gamma _0\gamma \sigma `$. By theorem 10.7 and lemma 10.5(ii) and (iii) we have
$$(T_1E_\gamma )(x_m^1)=B(E_\gamma (z),\stackrel{~}{T}_1.E_{x_m}(\stackrel{~}{z};\underset{¯}{\overset{~}{t}};q)).$$
By the commutation relation 2.21 between $`T_1`$ and $`f(z)`$ ($`f𝒜`$) and by the identity $`B(X,\stackrel{~}{X}\stackrel{~}{T}_1)=k_1B(X,\stackrel{~}{X})`$ (see lemma 10.5(i) and (iv)), we obtain
$$\begin{array}{cc}\hfill (& T_1E_\gamma )(x_m^1)=k_1B(E_\gamma (z),E_{x_m}(\stackrel{~}{z}^1))\hfill \\ & +\frac{(k_1k_1^1)+(k_0k_0^1)\gamma ^1}{\left(1\gamma ^2\right)}\left(B(E_\gamma (z),E_{x_m}(\stackrel{~}{z}))B(E_\gamma (z),E_{x_m}(\stackrel{~}{z}^1))\right).\hfill \end{array}$$
Here we have used lemma 10.5(ii), as well as the short-hand notation $`E_{x_m}(\stackrel{~}{z})=E_{x_m}(\stackrel{~}{z};\underset{¯}{\overset{~}{t}};q)`$. By lemma 10.5(ii), (iii) and theorem 10.7, we have
$$B(E_\gamma (z),E_{x_m}(\stackrel{~}{z}^1))=E_{x_m}(\gamma ;\underset{¯}{\overset{~}{t}};q)=E_{\gamma ^1}(x_m^1;\underset{¯}{t};q)$$
since $`\gamma \gamma _0`$. By theorem 10.7 it follows that the formula for $`(T_1E_\gamma )(x)`$ as stated in the proposition is correct for $`x=x_m^1`$ ($`m`$), hence it is correct as identity in $`𝒜`$.
For the second formula we proceed in a similar manner. First of all, observe that $`(T_1^{}E_\gamma )(x_m^1)=B(E_\gamma (z),\stackrel{~}{T}_0.E_{x_m}(\stackrel{~}{z};\underset{¯}{\overset{~}{t}};q))`$. We use now the commutation relation 2.21 between $`T_0`$ and $`f(z)`$ ($`f𝒜`$) and the identity $`B(X,\stackrel{~}{X}\stackrel{~}{T}_0)=u_1B(X,\stackrel{~}{X})`$, which follows from lemma 10.5(i) and (iv). Then
$$\begin{array}{cc}\hfill (& T_1^{}E_\gamma )(x_m^1)=u_1B(E_\gamma (z),E_{x_m}(q\stackrel{~}{z}^1))\hfill \\ & +\frac{(u_1u_1^1)+(u_0u_0^1)q^{1/2}\gamma }{(1q\gamma ^2)}(B(E_\gamma (z),E_{x_m}(\stackrel{~}{z}))B(E_\gamma (z),E_{x_m}(q\stackrel{~}{z}^1)).\hfill \end{array}$$
By lemma 9.2(ii), (iii) and theorem 10.7, we have
$$B(E_\gamma (z),E_{x_m}(q\stackrel{~}{z}^1))=E_{x_m}(q\gamma ;\underset{¯}{\overset{~}{t}};q)=E_{s_0\gamma }(x_m^1;\underset{¯}{t};q)$$
since $`q\gamma _n=\gamma _{n1}^1`$ for all $`n`$. It follows then by direct computations that the formula for $`(T_1^{}E_\gamma )(x)`$ as stated in the proposition is correct for $`x=x_m^1`$ ($`m)`$, hence it is correct as identity in $`𝒜`$. ∎
### 10.9.
The duality for the (non-)symmetric Askey-Wilson polynomials (see theorem 10.7) can also be used to derive recurrence relations for the (non-)symmetric Askey-Wilson polynomials from the difference(-reflection) equations $`LE_{s(\gamma )}^+=s(\gamma )E_{s(\gamma )}^+`$ (respectively $`YE_\gamma =\gamma E_\gamma `$). The Askey-Wilson $`q`$-difference equation $`LE_{s(\gamma )}=s(\gamma )E_\gamma `$ then gives the three term recurrence relation \[1, (1.24)–(1.27)\] for the symmetric Askey-Wilson polynomials (see van Diejen \[7, section 4\] for the argument in the higher rank setting).
## 11. The non-symmetric Askey-Wilson transform and its inverse
### 11.1.
We assume in this section that the parameters $`(\underset{¯}{t},q)`$ and $`(\underset{¯}{\overset{~}{t}},q)`$ satisfy the additional conditions 6.1.
Let $`\sigma ^{}`$ be the spectrum of $`Y^{}`$, so $`\sigma ^{}=\{\gamma _m^{}\}_m`$ with $`\gamma _m^{}=\gamma _m^1`$ for all $`m`$, see 3.5. Let $`F=F_{k_0,k_1,q}`$ be the functions $`g:\sigma ^{}`$ with finite support. By the non-degeneracy of the bilinear form $`..,._{\underset{¯}{t},q}`$ on $`𝒜`$ and by the bi-orthogonality relations 6.7 for the non-symmetric Askey-Wilson polynomials, we have a bijective linear map $`=_{\underset{¯}{t},q}:𝒜F`$ defined by
$$\left(_{\underset{¯}{t},q}(f)\right)(\gamma ):=f,E_\gamma ^{}()_{\underset{¯}{t},q},f𝒜,\gamma \sigma ^{},$$
where $`E_\gamma ^{}()=E_\gamma (;\underset{¯}{t}^1;q^1)`$ ($`\gamma \sigma ^{}`$) are the renormalized non-symmetric Askey-Wilson polynomials with respect to inverse parameters.
###### Definition.
The bijective map $`:𝒜F`$ is called the non-symmetric Askey-Wilson transform.
### 11.2.
Recall the action of $`𝒲`$ on $`\sigma ^{}`$ defined by $`s_0\gamma _m^{}=\gamma _{m1}^{}`$ and $`s_1\gamma _m^{}=\gamma _m^{}`$ for all $`m`$. This induces a left action of $`𝒲`$ on $`F`$ by $`(wg)(\gamma )=g(w^1\gamma )`$ for $`w𝒲`$, $`gF`$ and $`\gamma \sigma ^{}`$. Let $`\stackrel{~}{T}_i`$ ($`i=0,1`$) and $`\stackrel{~}{z}`$ be the linear endomorphisms of $`F`$ defined by $`(\stackrel{~}{z}g)(\gamma )=\gamma g(\gamma )`$,
$$\begin{array}{cc}\hfill \left(\stackrel{~}{T}_0g\right)(\gamma )& =u_1g(\gamma )+u_1^1\frac{(1u_0u_1q^{1/2}\gamma ^1)(1+u_0^1u_1q^{1/2}\gamma ^1)}{(1q\gamma ^2)}\left((s_0g)(\gamma )g(\gamma )\right),\hfill \\ \hfill \left(\stackrel{~}{T}_1g\right)(\gamma )& =k_1g(\gamma )+k_1^1\frac{(1k_0k_1\gamma )(1+k_0^1k_1\gamma )}{(1\gamma ^2)}\left((s_1g)(\gamma )g(\gamma )\right)\hfill \end{array}$$
for all $`gF`$ and all $`\gamma \sigma ^{}`$. Observe that these formulas can be obtained from the standard action of the generators $`\stackrel{~}{T}_i`$ ($`i=0,1`$) and $`\stackrel{~}{z}`$ of the double affine Hecke algebra $`(S;\underset{¯}{\overset{~}{t}};q)`$ on $`𝒜`$ (see 2.19 and 2.22) by formally replacing the $`𝒲`$-module $`𝒜`$ by the $`𝒲`$-module $`F`$.
###### Proposition.
There is a unique action of $`(S;\underset{¯}{\overset{~}{t}};q)`$ on $`F`$ such that the generators $`\stackrel{~}{z}`$ and $`\stackrel{~}{T}_i`$ ($`i=0,1`$) of $`(S;\underset{¯}{\overset{~}{t}};q)`$ act as the linear endomorphisms defined above. Furthermore,
$$(Xf)=\mu (X)(f),X(S;\underset{¯}{t};q),f𝒜,$$
where $`\mu `$ is the duality isomorphism defined in proposition 8.5.
###### Proof.
By proposition 8.7 we have
$$(Yf)(\gamma )=f,(Y^{})^1E_\gamma ^{}=\gamma ^1(f)(\gamma )=\left(\stackrel{~}{z}^1(f)\right)(\gamma )=\left(\mu (Y)(f)\right)(\gamma )$$
for all $`f𝒜`$ and all $`\gamma \sigma ^{}`$.
Again by proposition 8.7, we have $`(T_1f)(\gamma )=f,(T_1^{})^1E_\gamma ^{}`$. Combined with proposition 10.8, we derive that $`(T_1f)=\stackrel{~}{T}_1(f)=\mu (T_1)(f)`$ for all $`f𝒜`$.
In a similar manner, we derive from proposition 8.7 and proposition 10.8 that $`(T_1^{}f)=\stackrel{~}{T}_0(f)=\mu (T_1^{})(f)`$ for all $`f𝒜`$. The proposition now follows since $``$ is bijective and $`(S;\underset{¯}{t};q)`$ is generated as an algebra by $`Y`$, $`T_1`$ and $`T_1^{}`$. ∎
### 11.3.
Observe that the weight function $`\mathrm{\Delta }(\gamma ;\underset{¯}{\overset{~}{t}};q)`$ (see 6.2) has simple poles at $`\gamma \sigma ^{}`$. We define now a linear map $`𝒢=𝒢_{\underset{¯}{t},q}:F𝒜`$ by
$$𝒢_{\underset{¯}{t},q}(g)(x)=\underset{\gamma \sigma ^{}}{}g(\gamma )E_{\gamma ^1}(x;\underset{¯}{t};q)w(\gamma ;\underset{¯}{\overset{~}{t}};q),gF,$$
where $`w(\gamma )=w(\gamma ;\underset{¯}{\overset{~}{t}};q)`$ is defined by
$$w(\gamma ;\underset{¯}{\overset{~}{t}};q)=\underset{y=\gamma }{\text{Res}}\left(\frac{\mathrm{\Delta }(y;\underset{¯}{\overset{~}{t}};q)}{y}\right)\text{sgn}(\gamma ),\gamma \sigma ^{}$$
and $`\text{sgn}(\gamma _m^{})=ϵ(m)`$ for $`m`$ (see 3.3 for the definition of $`ϵ`$). Observe that $`w(\gamma ;\underset{¯}{\overset{~}{t}};q)=\alpha (\gamma ;k_1,k_0)w_+(\gamma ;\underset{¯}{\overset{~}{t}};q)`$ for all $`\gamma \sigma ^{}`$, where $`w_+(\gamma )=w_+(\gamma ;\underset{¯}{\overset{~}{t}};q)`$ is defined by
$$w_+(\gamma ;\underset{¯}{\overset{~}{t}};q)=\underset{y=\gamma }{\text{Res}}\left(\frac{\mathrm{\Delta }_+(y;\underset{¯}{\overset{~}{t}};q)}{y}\right)\text{sgn}(\gamma ),\gamma \sigma ^{},$$
see 6.2. The weight functions $`w(\gamma )`$ and $`w_+(\gamma )`$ can be written out explicitly in terms of $`q`$-shifted factorials, see \[1, section 2\] or \[8, section 7.5\] for $`w_+(\gamma )`$.
### 11.4.
In the following proposition we determine the intertwining properties of the $`(S;\underset{¯}{\overset{~}{t}};q)`$-action on $`F`$ under the linear map $`𝒢:F𝒜`$.
###### Proposition.
We have
$$𝒢(Xg)=\mu ^1(X)𝒢(g),X(S;\underset{¯}{\overset{~}{t}};q),gF,$$
where $`\mu `$ is the duality isomorphism defined in proposition 8.5.
###### Proof.
We write $`\stackrel{~}{T}_0=u_1+\stackrel{~}{\varphi }_0()(s_01)`$ and $`\stackrel{~}{T}_1=k_1+\stackrel{~}{\varphi }_1()(s_11)`$ with
$$\begin{array}{cc}\hfill \stackrel{~}{\varphi }_0(\gamma )& =u_1^1\frac{(1u_0u_1q^{1/2}\gamma ^1)(1+u_0^1u_1q^{1/2}\gamma ^1)}{(1q\gamma ^2)},\hfill \\ \hfill \stackrel{~}{\varphi }_1(\gamma )& =k_1^1\frac{(1k_0k_1\gamma )(1+k_0^1k_1\gamma )}{(1\gamma ^2)}.\hfill \end{array}$$
The weight function $`w(\gamma ;\underset{¯}{\overset{~}{t}};q)`$ ($`\gamma \sigma ^{}`$) satisfies the fundamental relations
$$\stackrel{~}{\varphi }_i(\gamma )w(\gamma ;\underset{¯}{\overset{~}{t}};q)=\stackrel{~}{\varphi }_i(s_i\gamma )w(s_i\gamma ;\underset{¯}{\overset{~}{t}};q),\gamma \sigma ^{},i=0,1.$$
This follows easily from the explicit expression for the weight function $`\mathrm{\Delta }`$, see 6.2 (compare also with the proof of proposition 6.6). It follows that
$$\begin{array}{cc}\hfill 𝒢(\stackrel{~}{T}_0g)(x)& =\underset{\gamma \sigma ^{}}{}(\stackrel{~}{T}_0g)(\gamma )E_{\gamma ^1}(x)w(\gamma )\hfill \\ & =\underset{\gamma \sigma }{}g(\gamma ^1)\left(u_1E_\gamma (x)+\stackrel{~}{\varphi }_0(\gamma ^1)(E_{s_0\gamma }(x)E_\gamma (x))\right)w(\gamma ^1)\hfill \\ & =\underset{\gamma \sigma }{}g(\gamma ^1)\left(T_1^{}E_\gamma \right)(x)w(\gamma ^1)=\left(T_1^{}𝒢(g)\right)(x)\hfill \end{array}$$
for all $`gF`$ by proposition 10.8. Similarly, we obtain $`𝒢(\stackrel{~}{T}_1g)(x)=\left(T_1𝒢(g)\right)(x)`$ for all $`gF`$ by proposition 10.8. Furthermore, it is immediate that $`𝒢(\stackrel{~}{z}g)(x)=(Y^1(𝒢(g))(x)`$ for all $`gF`$. We conclude that $`𝒢(Xg)=\mu ^1(X)𝒢(g)`$ for all $`gF`$ if $`X=\stackrel{~}{z}`$ or $`X=\stackrel{~}{T}_i`$ for $`i=0,1`$. The proposition follows, since these elements generate $`(S;\underset{¯}{\overset{~}{t}};q)`$ as an algebra. ∎
### 11.5.
Combining proposition 11.2 and proposition 11.4 leads to the following main result of this section.
###### Theorem.
(i) We have $`𝒢_{\underset{¯}{t},q}_{\underset{¯}{t},q}=c_{\underset{¯}{t},q}\text{Id}_𝒜`$ and $`_{\underset{¯}{t},q}𝒢_{\underset{¯}{t},q}=c_{\underset{¯}{t},q}\text{Id}_F`$ with the constant $`c_{\underset{¯}{t},q}`$ given by $`c_{\underset{¯}{t},q}=w(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)1,1_{\underset{¯}{t},q}`$.
(ii) For all $`\gamma \sigma `$ we have
$$\frac{E_\gamma ,E_{\gamma ^1}^{}_{\underset{¯}{t},q}}{1,1_{\underset{¯}{t},q}}=\frac{w(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)}{w(\gamma ^1;\underset{¯}{\overset{~}{t}};q)}.$$
###### Proof.
(i) Let $`f𝒜`$, then
$$𝒢\left(f\right)=𝒢\left((f(z)1)\right)=f(z)\left(𝒢((1))\right)$$
by proposition 11.2 and proposition 11.4, where $`1𝒜`$ is the function identically equal to one. By the definitions of $``$ and $`𝒢`$ and by the bi-orthogonality relations for the non-symmetric Askey-Wilson polynomials (see proposition 6.7) we have $`𝒢_{\underset{¯}{t},q}(_{\underset{¯}{t},q}(1))=c_{\underset{¯}{t},q}1`$ with the constant $`c_{\underset{¯}{t},q}`$ as given in the statement of the theorem. Hence $`𝒢=c\text{Id}_𝒜`$. The identity $`𝒢=c\text{Id}_F`$ follows then immediately from the fact that $`:𝒜F`$ is a bijection.
(ii) Let $`\gamma \sigma `$. By (i), we have $`𝒢((E_\gamma ))=cE_\gamma `$. On the other hand, by the explicit definitions of $``$ and $`𝒢`$ and by the bi-orthogonality relations for the non-symmetric Askey-Wilson polynomials (see proposition 6.7), we have
$$𝒢_{\underset{¯}{t},q}(_{\underset{¯}{t},q}(E_\gamma ))=w(\gamma ^1;\underset{¯}{\overset{~}{t}};q)E_\gamma ,E_{\gamma ^1}^{}_{\underset{¯}{t},q}E_\gamma .$$
Comparing coefficients of $`E_\gamma `$ leads to the desired result. ∎
### 11.6.
Let $`F^WF`$ be the $`W`$-invariant functions in $`F`$, i.e. the functions $`fF`$ satisfying $`f=s_1f`$. Equivalently, $`F^W`$ consists of the functions $`fF`$ satisfying $`\stackrel{~}{C}_+f=f`$, where $`\stackrel{~}{C}_+=(1+k_1^2)^1(1+k_1\stackrel{~}{T}_1)`$.
Let $`_+`$ be the restriction of the non-symmetric Askey-Wilson transform $``$ to $`𝒜^W𝒜`$ and let $`𝒢_+`$ be the restriction of $`𝒢`$ to $`F^WF`$.
###### Proposition.
$`_+`$ is a linear bijection from $`𝒜^W`$ to $`F^W`$ with inverse $`c^1𝒢_+`$, where $`c=c_{\underset{¯}{t},q}`$ is the constant defined in theorem 11.5. Furthermore, the constant $`c`$ can be rewritten as $`c=\frac{1}{2}(1+k_1^2)^2w_+(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)(1,1)_{\underset{¯}{t},q}`$, and
$$\begin{array}{cc}\hfill _+(f)(\gamma )& =\frac{1}{2}(1+k_1^2)(f,E_{s(\gamma )}^+(;\underset{¯}{t};q))_{\underset{¯}{t},q},\hfill \\ \hfill 𝒢_+(g)(x)& =(1+k_1^2)\underset{m_+}{}g(\gamma _m^{})E_{s(\gamma _m)}^+(x;\underset{¯}{t};q)w_+(\gamma _m^{};\underset{¯}{\overset{~}{t}};q)\hfill \end{array}$$
for all $`f𝒜^W`$ and all $`gF^W`$.
###### Proof.
By lemma 6.4 we have $`1,1_{\underset{¯}{t},q}=\frac{1}{2}(1+k_1^2)(1,1)_{\underset{¯}{t},q}`$. Furthermore, we have $`w(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)=(1+k_1^2)w_+(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)`$ since $`\alpha (\gamma _0^1;k_1,k_0)=1+k_1^2`$. This gives the alternative formula for the constant $`c_{\underset{¯}{t},q}`$.
Observe that $`E_{s(\gamma )}^+(x;\underset{¯}{t}^1;q^1)=E_{s(\gamma )}^+(x;\underset{¯}{t};q)`$ by proposition 6.9 and by the $`W`$-invariance of $`E_{s(\gamma )}^+`$. By lemma 6.4, 10.6 and the fact that $`C_+^{}=C_+^{}`$ (see the proof of proposition 6.8), we then derive for $`f𝒜^W`$ and $`\gamma \sigma ^{}`$ that
$$\begin{array}{cc}\hfill _+(f)(\gamma )=f,E_\gamma ^{}_{\underset{¯}{t},q}& =C_+f,E_\gamma ^{}_{\underset{¯}{t},q}=f,C_+^{}E_\gamma ^{}_{\underset{¯}{t},q}\hfill \\ & =f,E_{s(\gamma )}^+_{\underset{¯}{t},q}=\frac{1}{2}(1+k_1^2)(f,E_{s(\gamma )}^+)_{\underset{¯}{t},q}.\hfill \end{array}$$
In particular, we have $`(𝒜^W)F^W`$. For $`0m`$ we have
$$w(\gamma _m^{};\underset{¯}{\overset{~}{t}};q)+w(\gamma _m^{};\underset{¯}{\overset{~}{t}};q)=(1+k_1^2)w_+(\gamma _m^{};\underset{¯}{\overset{~}{t}};q).$$
Indeed, we use here that $`w_+(\gamma ;\underset{¯}{\overset{~}{t}};q)=w_+(\gamma ^1;\underset{¯}{\overset{~}{t}};q)`$ by the $`W`$-invariance of the weight function $`\mathrm{\Delta }_+(;\underset{¯}{\overset{~}{t}};q)`$, and that $`\alpha (\gamma ;k_1,k_0)+\alpha (\gamma ^1;k_1,k_0)=1+k_1^2`$, see 6.4. Hence we obtain for $`gF^W`$,
$$\begin{array}{cc}\hfill 𝒢_+(g)(x)=𝒢(\stackrel{~}{C}_+g)(x)& =\left(C_+𝒢(g)\right)(x)=\underset{\gamma \sigma ^{}}{}g(\gamma )E_{s(\gamma )}^+(x;\underset{¯}{t};q)w(\gamma ;\underset{¯}{\overset{~}{t}};q)\hfill \\ & =(1+k_1^2)\underset{m_+}{}g(\gamma _m^{})E_{s(\gamma _m)}^+(x;\underset{¯}{t};q)w_+(\gamma _m^{};\underset{¯}{\overset{~}{t}};q).\hfill \end{array}$$
In particular, $`𝒢(F^W)𝒜^W`$. Combined with proposition 11.5, this completes the proof of the proposition. ∎
###### Definition.
The bijection $`_+:𝒜^WF^W`$ is called the symmetric Askey-Wilson transform.
### 11.7.
We can repeat now the proof of theorem 11.5(ii) for the symmetric Askey-Wilson transform $`_+`$, using the alternative descriptions for $`_+`$ and $`𝒢_+`$ as given in proposition 11.6. This gives the following result on the quadratic norms of the symmetric Askey-Wilson polynomials.
###### Corollary.
For all $`\gamma \sigma `$ we have
$$\frac{(E_{s(\gamma )}^+,E_{s(\gamma )}^+)_{\underset{¯}{t},q}}{(1,1)_{\underset{¯}{t},q}}=\frac{w_+(\gamma _0^1;\underset{¯}{\overset{~}{t}};q)}{w_+(\gamma ^1;\underset{¯}{\overset{~}{t}};q)}.$$
## 12. The fundamental shift operator and the constant term
### 12.1.
In theorem 11.5 and corollary 11.7 we have obtained explicit expressions of $`E_\gamma ,E_{\gamma ^1}^{}_{\underset{¯}{t};q}`$ and of $`(E_{s(\gamma )}^+,E_{s(\gamma )}^+)_{\underset{¯}{t},q}`$ in terms of the constant term $`1,1_{\underset{¯}{t},q}=\frac{1}{2}(1+k_1^2)(1,1)_{\underset{¯}{t},q}`$ for all $`\gamma \sigma `$. The constant term $`(1,1)_{\underset{¯}{t},q}`$ is the well-known Askey-Wilson integral, which has been evaluated in many different ways, see for instance , , and . We give in this section yet another proof for the evaluation of $`(1,1)`$ using shift operators.
### 12.2.
In the following lemma we define explicit linear maps from symmetric Laurent polynomials to anti-symmetric Laurent polynomials and conversely in terms of the Cherednik-Dunkl operator $`Y`$. Recall the definition of the idempotents $`C_\pm H_0H=H(R;k_0,k_1)`$, see 5.1.
###### Lemma.
Let $`h_\pm (Y)=h_\pm (Y;k_0,k_1)H(R;k_0,k_1)`$ be defined by
$$h_\pm (Y)=Y^1(Y^{\pm 1}k_0k_1)(Y^{\pm 1}+k_0^1k_1).$$
(i) We have $`h_+(Y)C_+=C_{}h_+(Y)C_+`$ and $`h_{}(Y)C_{}=C_+h_{}(Y)C_{}`$ in $`H`$, i.e. $`h_\pm (Y)𝒜_\pm 𝒜_{}`$ under the action of the fundamental representation $`\pi _{\underset{¯}{t},q}`$.
(ii) We have $`C_\pm h_\pm (Y)C_{}=h_{}(Y)C_{}`$ in $`H`$.
###### Proof.
This follows by a straightforward computation from Lusztig’s formula 2.16, together with the fact that $`(T_1k_1^{\pm 1})C_\pm =0`$ in $`H_0H`$. ∎
### 12.3.
By lemma 12.2 and lemma 7.1 we have well-defined linear endomorphisms $`G_\pm (\underset{¯}{t};q):𝒜^W𝒜^W`$ defined by
$$\left(G_+f\right)(x)=\delta (x)^1\left(h_+(Y)f\right)(x),\left(G_{}f\right)=\left(h_{}(Y)(\delta .f)\right)(x),f𝒜^W,$$
where $`h_\pm (Y)`$ act under the fundamental representation $`\pi _{\underset{¯}{t},q}`$. Observe that $`G_{}`$ can be realized as the element $`h_{}(Y)\delta (z)`$ in $`(S;\underset{¯}{t};q)`$.
###### Proposition.
For $`m`$, we have
$$\begin{array}{cc}& G_+(\underset{¯}{t};q)P_m^+(;\underset{¯}{t};q)=h_+(\gamma _m;k_0,k_1)P_{m1}^+(x;k_0,qk_1,u_0,u_1;q),\hfill \\ & G_{}(\underset{¯}{t};q)P_{m1}^+(;k_0,qk_1,u_0,u_1;q)=h_{}(\gamma _m;k_0,k_1)P_m^+(x;\underset{¯}{t};q).\hfill \end{array}$$
###### Proof.
Let $`m`$. Since $`𝒜(m)=\text{span}\{P_m^+,P_m^{}\}`$ is an $`H`$-module with $`𝒜_{}(m)=𝒜(m)𝒜_{}=\text{span}\{P_m^{}\}`$, we have $`h_+(Y;k_0,k_1)P_m^+(;\underset{¯}{t};q)=c_+(m)P_m^{}(;\underset{¯}{t};q)`$ for some constant $`c_+(m)`$, see theorem 4.4, proposition 5.3 and lemma 12.2. Comparing leading terms using proposition 3.4, we see that $`c_+(m)=h_+(\gamma _m)`$. By the generalized Weyl character formula, see proposition 7.3, we obtain
$$G_+(\underset{¯}{t};q)P_m^+(;\underset{¯}{t};q)=h_+(\gamma _m)P_{m1}^+(;k_0,qk_1,u_0,u_1;q).$$
The shift property of $`G_{}`$ is proved in a similar manner. ∎
### 12.4.
In the remainder of this section we assume that the parameters satisfy the additional conditions 6.1. We write
$$\left(G_+(\underset{¯}{t}^1;q^1)f\right)(x)=\delta ^{}(x)^1\left(h_+(Y^{};k_0^1,k_1^1)f\right)(x),f𝒜$$
and $`G_{}(\underset{¯}{t}^1;q^1)=h_{}(Y^{};k_0^1,k_1^1)\delta ^{}(z^{})(S;\underset{¯}{t}^1;q^1)`$ for the shift-operators with respect to inverse parameters, where $`\delta ^{}(x)=\delta (x;k_0^1,k_1^1)`$ (see 7.2). The two shift operators $`G_+`$ and $`G_{}`$ are each-others adjoint in the following sense.
###### Proposition.
For all $`f,g𝒜^W`$, we have
$$\begin{array}{cc}\hfill (G_{}(\underset{¯}{t};q)f,g)_{\underset{¯}{t},q}& =(f,G_+(\underset{¯}{t}^1;q^1)g)_{k_0,qk_1,u_0,u_1,q},\hfill \\ \hfill (G_+(\underset{¯}{t};q)f,g)_{k_0,qk_1,u_0,u_1,q}& =k_1^4(f,G_{}(\underset{¯}{t}^1;q^1)g)_{\underset{¯}{t},q}.\hfill \end{array}$$
###### Proof.
Let $`f,g𝒜^W`$. By lemma 6.4, proposition 6.6 and lemma 7.1 we have
$$\begin{array}{cc}\hfill \frac{1}{2}(1+k_1^2)(G_{}(\underset{¯}{t};q)f,& g)_{\underset{¯}{t},q}=G_{}(\underset{¯}{t};q)f,g_{\underset{¯}{t},q}\hfill \\ & =\delta (z)f,h_{}((Y^{})^1)g_{\underset{¯}{t},q}=\delta (z)f,C_{}^{}h_{}((Y^{})^1)C_+^{}g_{\underset{¯}{t},q},\hfill \end{array}$$
where $`C_\pm ^{}(S;\underset{¯}{t}^1;q^1)`$ are the images of the primitive idempotents of $`H_0(k_1^1)`$ under $`\pi _{\underset{¯}{t}^1,q^1}`$, see 5.1 (compare with the proof of proposition 6.8). Now observe that $`h_{}((Y^{})^1;k_0,k_1)=k_1^2h_{}(Y^{};k_0^1,k_1^1)`$, hence
$$C_{}^{}h_{}((Y^{})^1)C_+^{}=k_1^2C_{}^{}h_{}(Y^{};k_0^1,k_1^1)C_+^{}=k_1^2h_+(Y^{};k_0^1,k_1^1)C_+^{}$$
by lemma 12.2. Consequently, we obtain
$$\begin{array}{cc}\hfill \frac{1}{2}(1+k_1^2)(G_{}(\underset{¯}{t};q)f,g)_{\underset{¯}{t},q}& =k_1^2\delta (z)f,h_+(Y^{};k_0^1,k_1^1)g_{\underset{¯}{t},q}\hfill \\ & =k_1^2\delta (z)f,\delta ^{}(z^{})G_+(\underset{¯}{t}^1;q^1)g_{\underset{¯}{t},q}\hfill \\ & =\frac{1}{2}(1+k_1^2)(f,G_+(\underset{¯}{t}^1;q^1)g)_{k_0,qk_1,u_0,u_1,q}\hfill \end{array}$$
where the last equality follows from lemma 7.2. The second formula is proved in a similar manner. ∎
### 12.5.
We write $`\nu (f)=\nu _{a,b,c,d}(f)=(f,f)_{\underset{¯}{t},q}`$ for the “quadratic norm” of $`f`$ with respect to the bilinear form $`(.,.)_{\underset{¯}{t},q}`$, where $`(a,b,c,d)`$ is the reparametrized multiplicity function, see 5.7. We use the short-hand notation
$$\nu _{a,b,c,d}(P_m^+)=\nu _{a,b,c,d}(P_m^+(;a,b,c,d)),m_+.$$
###### Corollary.
For $`m`$, we have
$$\nu _{a,b,c,d}(P_m^+)=\frac{(1q^m)(1q^{m1}cd)}{(1q^mab)(1q^{m1}abcd)}\nu _{qa,qb,c,d}(P_{m1}^+).$$
###### Proof.
This is an immediate consequence of proposition 12.3, proposition 12.4 and proposition 6.9. ∎
### 12.6.
Observe that the Askey-Wilson second order $`q`$-difference operator $`L`$ (see 5.8) and its eigenvalues $`\gamma _m+\gamma _m^1`$ ($`m_+`$) are symmetric in $`a,b,c,d`$. It follows that the symmetric Askey-Wilson polynomials $`P_m^+(x;a,b,c,d)`$ ($`m_+`$) are symmetric in the four parameters $`a,b,c,d`$. Hence corollary 12.5 can be reformulated with the special role of $`(a,b)`$ replaced by an arbitrary pair of the four parameters $`a`$, $`b`$, $`c`$, $`d`$. This leads to the following result.
###### Corollary.
Let $`k,l,m,n_+`$ and set $`t=k+l+m+n`$. Then
$$\begin{array}{cc}& \frac{\nu _{a,b,c,d}(P_t^+)}{\nu _{q^{2k}a,q^{2l}b,q^{2m}c,q^{2n}d}(1)}=\hfill \\ & \frac{(q,q^{2t1}abcd,q^{2k+2l}ab,q^{2k+2m}ac,q^{2k+2n}ad,q^{2l+2m}bc,q^{2l+2n}bd,q^{2m+2n}cd;q)_{\mathrm{}}}{(q^{t+1},q^{t1}abcd,q^tab,q^tac,q^tad,q^tbc,q^tbd,q^tcd;q)_{\mathrm{}}}.\hfill \end{array}$$
###### Proof.
We write
$$\frac{\nu _{a,b,c,d}(P_k^+)}{\nu _{q^2a,b,c,d}(P_{k1}^+)}=\frac{\nu _{a,b,c,d,q}(P_k^+)}{\nu _{qa,qb,c,d}(P_{k1}^+)}\frac{\nu _{qa,qb,c,d}(P_{k1}^+)}{\nu _{qa,b,q^1c,d}(P_k^+)}\frac{\nu _{qa,b,q^1c,d}(P_k^+)}{\nu _{q^2a,b,c,d}(P_{k1}^+)}$$
for $`k`$ and use the symmetry in the parameters $`a,b,c,d`$ and corollary 12.5 to obtain
$$\frac{\nu _{a,b,c,d}(P_k^+)}{\nu _{q^2a,b,c,d}(P_{k1}^+)}=\frac{(1q^k)(1q^{k1}bc)(1q^{k1}bd)(1q^{k1}cd)}{(1q^{k1}abcd)(1q^kab)(1q^kac)(1q^kad)}$$
for $`k`$. Now use again the symmetry in the parameters $`a,b,c,d`$ and complete induction with respect to $`k,l,m`$ and $`n`$ to obtain the desired result. ∎
### 12.7.
Corollary 12.6 relates the quadratic norm $`\nu (P_m^+)`$ to the constant term $`\nu (1)`$, but it can also be used to evaluate the Askey-Wilson integral $`\nu (1)`$ itself. The Askey-Wilson integral was evaluated for the first time by Askey and Wilson \[1, theorem 2.1\] (see also e.g. , and for alternative proofs).
###### Theorem (Constant term evaluation).
We have
$$\nu _{a,b,c,d}(1)=\frac{2(abcd;q)_{\mathrm{}}}{(q,ab,ac,ad,bc,bd,cd;q)_{\mathrm{}}}.$$
###### Proof.
Let $`(a,b,c,d)=(1,1,q^{\frac{1}{2}},q^{\frac{1}{2}})`$, then we have $`\mathrm{\Delta }_+(x)1`$ for the corresponding weight function of the bilinear form $`(.,.)`$. Hence $`P_k^+(x)=x^k+x^k`$ ($`k`$) for the corresponding symmetric Askey-Wilson polynomial. Furthermore,
$$\nu _{1,1,q^{\frac{1}{2}},q^{\frac{1}{2}}}(P_k)=2,k.$$
Combined with corollary 12.6 this implies that the theorem is correct for all parameter values $`(a,b,c,d)=(q^{2k},q^{2l},q^{\frac{1}{2}+2m},q^{\frac{1}{2}+2n})`$ with $`k,l,m,n_+`$ and $`k+l+m+n`$ (we use here that the formula in corollary 12.6 extends to this particular choice of parameter values by continuity). The proof is now completed by analytic continuation. ∎
### 12.8.
The constant term evaluation (theorem 12.7) and corollary 12.5 yield an explicit evaluation of $`\nu _{a,b,c,d}(P_m^+)`$ for $`m_+`$ which is in accordance with Askey and Wilson’s result \[1, theorem 2.2\]. As remarked in 7.5, this then yields explicit evaluations for all the diagonal terms in the bi-orthogonality relations of proposition 6.7 and proposition 6.8.
Another way to obtain the diagonal terms explicitly is by using corollary 11.7 and 10.3 (respectively theorem 11.5 and proposition 10.2) to reduce the diagonal terms for the (non-)symmetric Askey-Wilson polynomials to the constant term evaluation (theorem 12.7). Explicitly, we obtain the following formulas for the diagonal terms:
$$(P_m^+,P_m^+)=\frac{2(q^{2m1}abcd,q^{2m}abcd;q)_{\mathrm{}}}{(q^{m+1},q^mab,q^mac,q^mad,q^mbc,q^mbd,q^mcd,q^{m1}abcd;q)_{\mathrm{}}}$$
for $`m_+`$,
$$P_m,P_m^{}=\frac{(q^{2m}abcd,q^{2m}abcd;q)_{\mathrm{}}}{(q^{m+1},q^{m+1}ab,q^mac,q^mad,q^mbc,q^mbd,q^mcd,q^mabcd;q)_{\mathrm{}}}$$
for $`m_+`$,
$$P_m,P_m^{}=\frac{(q^{2m1}abcd,q^{2m1}abcd;q)_{\mathrm{}}}{(q^m,q^mab,q^mac,q^mad,q^mbc,q^mbd,q^{m1}cd,q^{m1}abcd;q)_{\mathrm{}}}$$
for $`m`$ and finally
$$P_m^{},P_m^{}=\frac{ab1}{ab}\frac{(q^{2m1}abcd,q^{2m}abcd;q)_{\mathrm{}}}{(q^m,q^{m+1}ab,q^mac,q^mad,q^mbc,q^mbd,q^{m1}cd,q^mabcd;q)_{\mathrm{}}}$$
for $`m`$.
### 12.9.
There is yet another way to relate the diagonal terms of the non-symmetric Askey-Wilson polynomials to the constant term $`1,1`$. This method is based on the Rodrigues type formula for the non-symmetric Askey-Wilson polynomial (see proposition 9.3), which allows us to compute the diagonal terms by induction with respect to the degree of the non-symmetric Askey-Wilson polynomial. For the induction step, one needs the following two additional properties of the intertwiners. The first property is that $`S_0^{}=q^1S_0^{}`$ and $`S_1^{}=S_1^{}`$, where $`S_0^{},S_1^{}(S;\underset{¯}{t}^1;q^1)`$ are the intertwiners with respect to inverse parameters, cf. proposition 8.7. The second property is
$$\begin{array}{cc}\hfill S_0^2& =q^1u_1^2\underset{\xi =\pm 1}{}(1u_0^1u_1^1q^{\xi /2}Y^\xi )(1+u_0u_1^1q^{\xi /2}Y^\xi ),\hfill \\ \hfill S_1^2& =k_1^2\underset{\xi =\pm 1}{}(1k_0^1k_1^1Y^\xi )(1+k_0k_1^1Y^\xi )\hfill \end{array}$$
which are most easily proved in the image of the duality isomorphism $`\mu `$, see Sahi \[19, corollary 5.2\] in the higher rank setting. We leave the details to the reader. |
warning/0001/hep-ph0001197.html | ar5iv | text | # 1 Overview
## 1 Overview
In this talk I want to consider the emergence of ’classical’ field configurations - topological defects - after a phase transition, and the extent to which thermal fluctuations can inhibit this process. This is of particular interest in the early universe, for which we expect a sequence of transitions from a very symmetric initial state, and in which the presence of classical defects can have important astrophysical consequences.
The relevance of topological defects is that when symmetry breaks, it does not do so uniformly. At the very least, the field is uncorrelated on the scale of the causal horizon at any time. Since a broken symmetry is, necessarily, characterised by degenerate vacua, the choice of different vacua in domains in which the fields are uncorrelated will lead naturally to topological defects between them as the field does its best to order on large scales. The nature of the defects depends on the relevant homotopy group of the ground-state manifold. The most acceptable defect on cosmological grounds is the ’cosmic string’ - a generalised field vortex - which may have played in role in structure formation.
Given a theory that permits vortices, at some time after the transition we expect to find a network of them, behaving essentially classically as Nambu-Goto strings, intersecting, chopping off loops which decay, and straightening segments to reduce field gradients. Simple calculations suggest late-time scaling solutions, with similarity on a wide range of scales.
The details do not concern us here. What interests us is how this collection of essentially classical objects, which can be observed directly, in principle, came into existence. The simplest question, that we shall address here, is what is the density of cosmic strings (or other defects) at the time of their appearance?
The early universe is very hot, but such a problem requires us to go beyond equilibrium Thermal Field Theory. In practice, we often know remarkably little about the dynamics of thermal systems. For simplicity, I shall assume scalar field order parameters, with continuous transitions. In principle, the field correlation length diverges at a continuous transition. In practice, it does not since there is not enough time. One possibility is that the separation of ’defects’ is characterised by the correlation length when it checks its growth. If this were simply so, a measurement of defect densities would be a measurement of correlation lengths. Estimates of this early field ordering and its contingent defects in the early universe have been made by Kibble, using simple thermal arguments or causal arguments different from the one above (although that is also due to Kibble).
There are great difficulties in converting such predictions for the early universe into experimental observations since, but for a possible stray monopole, we have no direct evidence for them having existed <sup>2</sup><sup>2</sup>2Although this does not impede our ability to make predictions for defect-driven fluctuations in the CMB, for example. Zurek suggested that similar arguments to those in were applicable to condensed matter systems for which direct experiments on defect densities could be performed. This has lead to considerable activity from theorists working on the boundary between QFT and Condensed Matter theory and from condensed matter experimentalists. To date almost all experiments have involved superfluids, for which vortices can be produced readily. All but one experiment is in agreement with these simple causal predictions and we shall pay particular attention to this one failure of prediction. In this talk I shall
* review the Kibble/Zurek causality predictions for initial correlation lengths and defect densities.
* summarise the results of the condensed matter experiments and present an alternative picture for the onset of defect production for condensed matter systems. I shall then show how this alternative picture gives essentially the same results as the Zurek picture for those condensed matter systems for which there is experimental agreement.
* provide an explanation for why some condensed matter experiments will be in disagreement with Zurek’s predictions, including the experiment in question. We shall suggest that the prediction fails, in part, because of the presence of thermal noise.
* use these ideas to address the more complicated problem of the appearance of ’classical’ defect configurations in QFT in the light of Kibble’s predictions, and the role of thermal noise in them.
## 2 When symmetry breaks, how big are the smallest identifiable pieces?
Defects in the large-scale ordering of the field can only appear once the transition has taken place. If it is the case that defect density can be identified simply from the field correlation length, the maximum density (an experimental observable in condensed matter systems, although not for the early universe) will be associated with the smallest identifiable correlation length in the broken phase once the transition has been effected. This provides the initial condition for the evolution of field ordering. From this viewpoint, we can observe the defects by default merely by determining the correlation length at that time.
In order to see how to identify these ’smallest pieces’<sup>3</sup><sup>3</sup>3The title of this section is essentially that posed in recent papers by Zurek. it is sufficient to consider the simplest theory with vortices, that of a single relativistic complex scalar field in three spatial dimensions, undergoing a temperature quench. In the first instance we assume that the qualitative dynamics of the transition are conditioned by the field’s equilibrium free energy, of the form
$$F(T)=d^3x\left(|\varphi |^2+m^2(T)|\varphi |^2+\lambda |\varphi |^4\right)$$
(2.1)
Prior to the transition, at temperature $`T>T_c`$, the critical temperature, $`m(T)>0`$ plays the role of an effective ’plasma’ mass due to the interactions of $`\varphi `$ with the heat bath, which includes its own particles. After the transition, when $`T`$ is effectively zero, $`m^2(0)=M^2<0`$ enforces the $`U(1)`$ symmetry-breaking, with field expectation values $`|\varphi |=\eta `$, $`\eta ^2=M^2/\lambda `$. The change in temperature that leads to the change in the sign of $`m^2`$ is most simply understood as a consequence of the system expanding. Thus, in the early universe, a weakly interacting relativistic plasma at temperature $`TM`$ has an entropy density $`sT^3`$. As long as thermal equilibrium can be maintained, constant entropy $`S`$ per comoving volume, $`Ssa(t)^3`$, gives $`Ta(t)^1`$ and falling, for increasing scale factor $`a(t)`$. Models that attempt to take inflation into account, however, lead to ’preheating’ that is not Boltzmannian. Nonetheless, even in such cases it is possible to isolate an effective temperature for long-wavelength modes. This is all that is necessary, but is too sophisticated for the simple scenarios that we shall present here. We shall not even include a metric in Eq.2.1.
The minima of the final potential of Eq.2.1 now constitute the circle $`\varphi =\eta e^{i\alpha }`$. When the transition starts $`\varphi `$ begins to fall into the valley of the potential, choosing a random phase. This randomly chosen phase can vary from point to point subject to continuity. At late times the failure of the field to be uniform in phase on large scales will lead to it twisting around classical ’defects’ - solutions to $`\delta F/\delta \varphi =0`$ that locally minimise the energy stored in field gradients and potentials. Those of interest to us are vortices, tubes of ’false’ vacuum $`\varphi 0`$, around which the field phase changes by $`\pm 2\pi `$. In an early universe context these are the simplest possible ’cosmic strings’.
How this collapse takes place determines the size of the first identifiable domains. It was suggested by Kibble and Zurek that this size is essentially the equilibrium field correlation length $`\xi _{eq}`$ at some appropriate temperature close to the transition. I shall argue later that this is too simple but, nonetheless, it is a plausible starting point. Two very different mechanisms have been proposed for estimating this size.
### 2.1 Thermal activation
In the early work on the cosmic string scenario an alternative possibility to simple causality was to assume that initial domain size was fixed in the Ginzburg regime by the correlation length at that time, rather than the causal radius. By this we mean the following. Once we are below $`T_c`$, and the central hump in $`V(\varphi )=m^2(T)|\varphi |^2+\lambda |\varphi |^4`$ is forming, $`T_G`$ signals the temperature above which there is a significant probability for thermal fluctuations over the central hump on the scale of the correlation length. Most simply, it is determined by the condition
$$\mathrm{\Delta }V(T_G)\xi _{eq}^3(T_G)T_G$$
(2.2)
where $`\mathrm{\Delta }V(T)`$ is the difference between the central maximum and the minima of $`V(\varphi ,T)`$. We find $`|1T_G/T_c|=O(\lambda )`$.
Whereas, above $`T_G`$ there will be a population of ’domains’, fluctuating in and out of existence, at temperatures below $`T_G`$ fluctuations from one minimum to the other become increasingly unlikely. When this happens the correlation length is
$$\xi _{eq}(T_G)=O\left(\frac{\xi _0}{\sqrt{1T_G/T_c}}\right),$$
(2.3)
where $`\xi _0=M^1`$ is the natural unit of length, the Compton wavelength of the $`\varphi `$ particles.
It is tempting to identify $`\xi _{eq}(T_G)`$ with the scale at which stable domains begin to form. We shall see later that this is incorrect, for quenches that are not too slow. However, some care is needed if (as can happen in condensed matter physics) we never leave the Ginzburg regime.
The formation of large domains is an issue that requires more than equilibrium physics. The most simple dynamical arguments can be understood in terms of causality.
### 2.2 Causality
We have already mentioned that causality puts an upper bound on domain size. Specifically, if $`G(r,t)`$ is the two-field correlation function at time $`t`$ for separation $`r`$, then $`G`$ vanishes for $`r2t`$ approximately. This was used by Kibble to put an upper bound on monopole density in the early universe. If this causal bound and the Ginzburg criteria attempt to set scales once the critical temperature has been passed, the causal arguments considered now attempt to set scales before it is reached.
Here we attempt a lower bound on domain size, an upper bound on defect density. Suppose the temperature $`T(t)`$ varies sufficiently slowly with time $`t`$ that it makes sense to replace $`V(\varphi ,T)`$ by $`V(\varphi ,T(t))`$. With $`m^2(T(t))`$ vanishing at $`T=T_c`$, which we suppose happens at $`t=0`$, the equilibrium correlation length of the field fluctuations $`\xi _{eq}(T(t))=|m^1(T(t))|`$ diverges at $`T(t)=T_c`$. It is sufficient to adopt a mean-field approximation in which $`m^2(T)(TT_c)`$. The true correlation length $`\xi (t)`$ cannot diverge like $`\xi _{eq}(T(t))`$, since it can only grow so far in a finite time.
Initially, for $`t<0`$, when we are far from the transition, we again assume effective equilibrium, and the field correlation length $`\xi (t)`$ tracks $`\xi _{eq}(T(t))`$ approximately. However, as we get closer to the transition $`\xi _{eq}(T(t))`$ begins to increase arbitrarily fast. As a crude upper bound, the true correlation length fails to keep up with $`\xi _{eq}(T(t))`$ by the time $`\overline{t}`$ at which $`\xi _{eq}`$ is growing at the speed of light, $`d\xi _{eq}(T(\overline{t}))/dt=1`$. It was suggested by Kibble that, once we have reached this time $`\xi (t)`$ freezes in, remaining approximately constant until the time $`t+\overline{t}`$ after the transition when is once again becomes comparable to the now decreasing value of $`\xi _{eq}`$. The correlation length $`\xi _{eq}(\overline{t})=\xi _{eq}(\overline{t})`$ is argued to provide the scale for the minimum domain size after the transition.
Specifically, if we assume a time-dependence $`m^2(t)=M^2t/t_Q`$ in the vicinity of $`t=0`$, when the transition begins to be effected, then the causality condition gives $`t_C=t_Q^{1/3}(2M)^{2/3}`$. As a result,
$$M\xi _{eq}(\overline{t})=(M\tau _0)^{1/3},$$
(2.4)
which, with condensed matter in mind, we write as
$$\overline{\xi }=\xi _{eq}(\overline{t})=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^{1/3}$$
(2.5)
where $`\tau _0=\xi _0=M^1`$ are the natural time and distance scales. In contrast to Eq.2.3, Eq.2.5 depends explicitly on the quench rate, as we would expect.
### 2.3 QFT or Condensed Matter
This approach of Kibble was one of the motivations for a similar analysis by Zurek of transitions in condensed matter. Qualitatively, neither the Ginzburg thermal fluctuations, with fluctuation length Eq.2.3, nor the simple causal argument above depend critically on the fact that the free energy Eq.2.1 is originally assumed to be derived from a relativistic quantum field theory. After rescaling, $`F`$ could equally well be the Ginzburg-Landau free energy for the complex order-parameter field whose magnitude determines the superfluid density. That is,
$$F(T)=d^3x\left(\frac{\mathrm{}^2}{2m}|\varphi |^2+\alpha (T)|\varphi |^2+\frac{1}{4}\beta |\varphi |^4\right),$$
(2.6)
in which $`\alpha (T)m^2(T)`$ vanishes at the critical temperature $`T_c`$. The only difference is that, in the causal argument, the speed of light should be replaced by the speed of (second) sound, with different critical index.
Explicitly, let us assume the mean-field result $`\alpha (T)=\alpha _0ϵ(T)`$, where $`ϵ=(T/T_c1)`$, remains valid as $`T/T_c`$ varies with time $`t`$. In particular, we first take $`\alpha (t)=\alpha (T(t))=\alpha _0t/\tau _Q`$ in the vicinity of $`T_c`$. The fundamental length scale $`\xi _0`$ is given from Eq.2.6 as $`\xi _0^2=\mathrm{}^2/2m\alpha _0`$. The Gross-Pitaevski theory suggests a natural time-scale $`\tau _0=\mathrm{}/\alpha _0`$. When, later, we adopt the time-dependent Landau-Ginzburg (TDLG) theory we find this still to be true, empirically, at order-of-magnitude level, and we keep it.
It follows that the equilibrium correlation length $`\xi _{eq}(t)`$ and the relaxation time $`\tau (t)`$ diverge when $`t`$ vanishes as
$$\xi _{eq}(t)=\xi _0\left|\frac{t}{\tau _Q}\right|^{1/2},\tau (t)=\tau _0\left|\frac{t}{\tau _Q}\right|^1.$$
(2.7)
The speed of sound is $`c(t)=\xi _{eq}(t)/\tau (t)`$, slowing down as we approach the transition as $`|t|^{1/2}`$. The causal counterpart to $`d\xi _{eq}(t)/dt=1`$ for the relativistic field is $`d\xi _{eq}(t)/dt=c(t)`$. This is satisfied at $`t=\overline{t}`$, where $`\overline{t}=\sqrt{\tau _Q\tau _0}`$, with corresponding correlation length
$$\overline{\xi }=\xi _{eq}(\overline{t})=\xi _{eq}(\overline{t})=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^{1/4}.$$
(2.8)
(cf. Eq.2.5). A variant of this argument that gives essentially the same results is obtained by comparing the quench rate directly to the relaxation rate of the field fluctuations. We stress that, yet again, the assumption is that the length scale that determines the initial correlation length of the field freezes in before the transition begins.
## 3 Experiments
The end result of the simple causality arguments is that, both for QFT and condensed matter, when the field begins to order itself its correlation length has the form
$$\overline{\xi }=\xi _{eq}(\overline{t})=\xi _0\left(\frac{\tau _Q}{\tau _0}\right)^\gamma .$$
(3.9)
for appropriate $`\gamma `$. <sup>4</sup><sup>4</sup>4In fact, the powers of Eq.2.5 and Eq.2.8 are mean-field results, changed on implementing the renormalisation group.
The jump that Kibble made in QFT was to assume that the correlation length Eq.3.9 also sets the scale for the typical minimum intervortex distance. That is, the initial vortex density $`n_{def}`$ is<sup>5</sup><sup>5</sup>5Equivalently, the length of vortices in a box volume $`v`$ is $`O(n_{def}v)`$. is assumed to be
$$n_{def}=\frac{1}{f^2}\frac{1}{\overline{\xi }^2}=\frac{1}{f^2\xi _0^2}\left(\frac{\tau _0}{\tau _Q}\right)^{2\gamma }.$$
(3.10)
for $`\gamma =1/3`$ and $`f=O(1)`$. We stress that this assumption is independent of the argument that lead to Eq.2.5. Since $`\xi _0`$ also measures cold vortex thickness, $`\tau _Q\tau _0`$ corresponds to a measurably large number of widely separated vortices.
Even if cosmic strings were produced in so simple a way in the very early universe it is not possible to compare the density Eq.3.10 with experiment, in large part because of our uncertainty as to what is the appropriate theory. It was Zurek who first suggested that this causal argument for defect density be tested in condensed matter systems.
### 3.1 Superfluid helium
Vortex lines in both superfluid $`{}_{}{}^{4}He`$ and $`{}_{}{}^{3}He`$ are good analogues of global cosmic strings. In $`{}_{}{}^{4}He`$ the bose superfluid is characterised by a complex field $`\varphi `$, whose squared modulus $`|\varphi |^2`$ is the superfluid density. The Landau-Ginzburg theory for $`{}_{}{}^{4}He`$ has, as its free energy, $`F(T)`$ of Eq.2.6. The static classical field equation $`\delta F/\delta \varphi =0`$ has vortex solutions as before. Specifically, a simple (winding number unity) static straight vortex along the z-axis has the form
$$\mathrm{\Phi }(𝐱)=h(r)e^{\pm i\theta },$$
(3.11)
where $`\theta =\mathrm{arctan}(y/x)`$ and $`r^2=x^2+y^2`$. For small $`r`$, $`h(r)=O(r)`$, and for large $`r`$, $`h(r)=\eta (1O(\xi _0^2/r^2))`$, with effective width $`\xi _0`$.
The situation is more complicated, but more interesting, for $`{}_{}{}^{3}He`$. The reason is that $`{}_{}{}^{3}He`$ is a fermion. Thus the mechanism for superfluidity is very different from that of $`{}_{}{}^{4}He`$. Somewhat as in a BCS superconductor, these fermions form the counterpart to Cooper pairs. However, whereas the (electron) Cooper pairs in a superconductor form a $`{}_{}{}^{1}S`$ state, the $`{}_{}{}^{3}He`$ pairs form a $`{}_{}{}^{3}P`$ state. The order parameter $`A_{\alpha i}`$ is a complex $`3\times 3`$ matrix $`A_{\alpha i}`$. There are two distinct superfluid phases, depending on how the $`SO(3)\times SO(3)\times U(1)`$ symmetry is broken. If the normal fluid is cooled at low pressures, it makes a transition to the $`{}_{}{}^{3}HeB`$ phase.
The Landau-Ginzburg free energy is, necessarily, more complicated, permitting many types of vortex, but the effective potential $`V(A_{\alpha i},T)`$ has the diagonal form $`V(A,T)=\alpha (T)|A_{ai}|^2+O(A^4)`$ for small fluctuations, and this is all that we need for the production of vortices at very early times. Thus the Zurek analysis leads to the prediction Eq.3.10, as before, for appropriate $`\gamma `$.
### 3.2 Counting vortices
Although $`{}_{}{}^{3}He`$ is more complicated to work with, the experiments to check Eq.3.10 are cleaner, since even individual vortices can be detected by magnetic resonance. Second, because vortex width is many atomic spacings the Landau-Ginzburg theory is good ($`\gamma =1/4`$).
So far, experiments have been of two types. In the Helsinki experiment superfluid $`{}_{}{}^{3}He`$ in a rotating cryostat is bombarded by slow neutrons. Each neutron entering the chamber releases 760 keV, via the reaction $`n+^3Hep+^3He+760keV`$. The energy goes into the kinetic energy of the proton and triton, and is dissipated by ionisation, heating a region of the sample above its transition temperature. The heated region then cools back through the transition temperature, creating vortices. Vortices above a critical size grow and migrate to the centre of the apparatus, where they are counted by an NMR absorption measurement. The quench is very fast, with $`\tau _Q/\tau _0=O(10^3)`$. Agreement with Eq.3.10 and Eq.2.8 is good, at the level of an order of magnitude. This is even though it is now argued that the Helsinki experiment should not show agreement because of the geometry of the heating event.
The second type of experiment has been performed at Grenoble and Lancaster. Rather than count individual vortices, the experiment detects the total energy going into vortex formation. As before, $`{}_{}{}^{3}He`$ is irradiated by neutrons. After each absorption the energy released in the form of quasiparticles is measured, and found to be less than the total 760 keV. This missing energy is assumed to have been expended on vortex production. Again, agreement with Zurek’s prediction Eq.3.10 and Eq.2.8 is good.
The experiments in $`{}_{}{}^{4}He`$, conducted at Lancaster, follow Zurek’s original suggestion. The idea is to expand a sample of normal fluid helium so that it becomes superfluid at essentially constant temperature. That is, we change $`1T/T_c`$ from negative to positive by reducing the pressure and increasing $`T_c`$. As the system goes into the superfluid phase a tangle of vortices is formed, because of the random distribution of field phases. The vortices are detected by scattering second sound off them, and its attenuation gives a good measure of vortex density. A mechanical quench is slow, with $`\tau _Q`$ some tens of milliseconds, and $`\tau _Q/\tau _0=O(10^{10})`$<sup>6</sup><sup>6</sup>6For $`{}_{}{}^{4}He`$ mean-field theory is poor, and a better value for $`\gamma `$ is $`\gamma =1/3`$.. Two experiments have been performed. In the first fair agreement was found with the prediction Eq.3.10, but the second experiment has failed to see any vortices whatsoever.
There is certainly no agreement, in this or any other experiment on $`{}_{}{}^{3}He`$, with the thermal fluctuation density that would be based on Eq.2.3.
## 4 The Kibble-Zurek picture for the value of $`\overline{\xi }`$ is correct.
To do better than the simple causality arguments we need a concrete model for the dynamics.
### 4.1 Condensed matter: the TDLG equation at early times
We assume that the dynamics of the transition can be derived from the explicitly time-dependent Landau-Ginzburg free energy
$$F(t)=d^3x\left(\frac{\mathrm{}^2}{2m}|\varphi |^2+\alpha (t)|\varphi |^2+\frac{1}{4}\beta |\varphi |^4\right).$$
(4.12)
obtained from Eq.2.6 on identifying $`\alpha (t)=\alpha (T(t))=\alpha _0ϵ(t)`$, where $`ϵ=(T/T_c1)`$. In a quench in which $`T_c`$ or $`T`$ changes it is convenient to shift the origin in time, to write $`ϵ(t)`$ as
$$ϵ(t)=ϵ_0\frac{t}{\tau _Q}\theta (t)$$
(4.13)
for $`\mathrm{}<t<\tau _Q(1+ϵ_0)`$, after which $`ϵ(t)=1`$. $`ϵ_0`$ measures the original relative temperature and $`\tau _Q`$ defines the quench rate. The quench begins at time $`t=0`$ and the transition from the normal to the superfluid phase begins at time $`t=ϵ_0\tau _Q`$. Times subsequent to that are defined by $`\mathrm{\Delta }t=tt_0`$.
Motivated by Zurek’s later numerical simulations, we adopt the time-dependent Landau-Ginzburg (TDLG) equation for $`F`$, on expressing $`\varphi `$ as $`\varphi =(\varphi _1+i\varphi _2)/\sqrt{2}`$, that
$$\frac{1}{\mathrm{\Gamma }}\frac{\varphi _a}{t}=\frac{\delta F}{\delta \varphi _a}+\eta _a,$$
(4.14)
where $`\eta _a`$ is Gaussian thermal noise, satisfying
$$\eta _a(𝐱,t)\eta _b(𝐲^{},t^{})=2\delta _{ab}T(t)\mathrm{\Gamma }\delta (𝐱𝐲)\delta (tt^{}).$$
(4.15)
This is a crude approximation for $`{}_{}{}^{4}He`$, and a simplified form of a realistic description of $`{}_{}{}^{3}He`$ but it is not a useful description of QFT, as it stands.
It is relatively simple to determine the validity of Zurek’s argument since it assumes that freezing in of field fluctuations occurs just before symmetry breaking begins. At that time the effective potential $`V(\varphi ,T)`$ is still roughly quadratic and we can see later that, for the relevant time-interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$ the self-interaction term can be neglected ($`\beta =0`$).
In space, time and temperature units in which $`\xi _0=\tau _0=k_B=1`$, Eq.4.14 then becomes
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ(t)]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t).$$
(4.16)
where $`\overline{\eta }`$ is the renormalised noise. The solution of the ’free’-field linear equation is straightforward, giving a Gaussian equal-time correlation function
$$\varphi _a(𝐫,t)\varphi _b(\mathrm{𝟎},t)=\delta _{ab}G(𝐫,t)$$
(4.17)
where
$$G(r,t)=_0^{\mathrm{}}𝑑\tau \overline{T}(t\tau /2)\left(\frac{1}{4\pi \tau }\right)^{3/2}e^{r^2/4\tau }e^{_0^\tau 𝑑sϵ(ts/2)}.$$
(4.18)
and $`\overline{T}`$ is the renormalised temperature. At time $`t_0=ϵ_0\tau _0`$, when the transition begins, a saddle-point calculation shows that, provided the quench is not too fast,
$$G(r,t_0)\frac{T_c}{4\pi r}e^{a(r/\overline{\xi })^{4/3}},$$
(4.19)
where $`a=O(1)`$, confirming Zurek’s result.
Zurek’s prediction is robust, since further calculation shows that $`\xi (t)`$ does not vary strongly in the interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$, where $`\mathrm{\Delta }t=tt_0`$.
### 4.2 QFT: Closed time-path ensemble averaging at early times
For QFT the situation is rather different. In the previous section, instead of working with the TDLG equation, we could have worked with the equivalent Fokker-Planck equation for the probability $`p_t^{FP}[\mathrm{\Phi }]`$ that, at time $`t>0`$, the measurement of $`\varphi `$ will give the function $`\mathrm{\Phi }(𝐱)`$. Thus $`G(r,t)`$ of Eq.4.17 can be written as
$$\delta _{ab}G(𝐫,t)=\varphi _a(𝐫,t)\varphi _b(\mathrm{𝟎},t)=𝒟\mathrm{\Phi }p_t^{FP}[\mathrm{\Phi }]\mathrm{\Phi }_a(𝐫)\mathrm{\Phi }_b(\mathrm{𝟎}).$$
(4.20)
When solving the dynamical equations for a hot quantum field it is convenient to work with probabilities from the start. Taking $`t=0`$ as our starting time for the evolution of the complex field $`\varphi `$ suppose that, at this time, the system is in a pure state, in which the measurement of $`\varphi `$ would give $`\mathrm{\Phi }_0(𝐱)`$. That is:-
$$\widehat{\varphi }(t=0,𝐱)|\mathrm{\Phi }_0,t=0=\mathrm{\Phi }_0|\mathrm{\Phi }_0,t=0.$$
(4.21)
The probability $`p_t[\mathrm{\Phi }]`$ that, at time $`t>0`$, the measurement of $`\varphi `$ will give the function $`\mathrm{\Phi }(𝐱)`$, is the double path integral
$$p_t[\mathrm{\Phi }]=_{\varphi _\pm (0)=\mathrm{\Phi }_0}^{\varphi _\pm (t)=\mathrm{\Phi }}𝒟\varphi _+𝒟\varphi _{}\mathrm{exp}\left\{i\left(S_t[\varphi _+]S_t[\varphi _{}]\right)\right\},$$
(4.22)
where $`𝒟\varphi _\pm =𝒟\varphi _{\pm ,1}𝒟\varphi _{\pm ,2})`$ and $`S_t[\varphi ]`$ is the (time-dependent) action obtained from Eq.2.1, on substituting $`m(t)=m(T(t))`$ for $`m(T)`$.
$`p_t[\mathrm{\Phi }]`$ can be written in the closed time-path form in which, instead of separately integrating $`\varphi _\pm `$ along the time paths $`0tt_f`$, the integral can be interpreted as time-ordering of a field $`\varphi `$ along the closed path $`C_+C_{}`$ of Fig.1, where $`\varphi =\varphi _+`$ on $`C_+`$ and $`\varphi =\varphi _{}`$ on $`C_{}`$. When we extend the contour from $`t_f`$ to $`t=\mathrm{}`$ either $`\varphi _+`$ or $`\varphi _{}`$ is an equally good candidate for the physical field, but we choose $`\varphi _+`$.
The choice of a pure state at time $`t=0`$ is too simple to be of any use. For simplicity, we assume that $`\mathrm{\Phi }`$ is Boltzmann distributed at time $`t=0`$ at an effective temperature of $`T_0=\beta _0^1`$ according to the Hamiltonian $`H[\mathrm{\Phi }]`$ corresponding to the free-field action $`S[\varphi ]`$, obtained by setting $`\lambda =0`$ in Eq.2.1, in which $`\varphi `$ is taken to be periodic in imaginary time with period $`\beta _0`$. We now have the explicit form for $`p_t[\mathrm{\Phi }]`$,
$$p_t[\mathrm{\Phi }]=_B𝒟\varphi e^{iS_C[\varphi ]}\delta [\varphi _+(t_f)\mathrm{\Phi }],$$
(4.23)
written as the time ordering of a single field along the contour $`C=C_+C_{}C_3`$, extended to include a third imaginary leg, where $`\varphi `$ takes the values $`\varphi _+`$, $`\varphi _{}`$ and $`\varphi _3`$ on $`C_+`$, $`C_{}`$ and $`C_3`$ respectively, for which $`S_C`$ is $`S[\varphi _+]`$, $`S[\varphi _{}]`$ and $`S_0[\varphi _3]`$.
Just as we had no need to calculate $`p^{FP}[\mathrm{\Phi }]_t`$ explicitly in condensed matter we can average in QFT without having to calculate $`p_t[\mathrm{\Phi }]`$ explicitly. Specifically,
$$G_{ab}(r,t)=\mathrm{\Phi }_a(𝐫)\mathrm{\Phi }_b(\mathrm{𝟎})_t=𝒟\mathrm{\Phi }p_t[\mathrm{\Phi }]\mathrm{\Phi }_a(𝐫)\mathrm{\Phi }_b(\mathrm{𝟎})$$
(4.24)
is given by
$$G_{ab}(r,t)=\varphi _a(𝐫,t)\varphi _b(\mathrm{𝟎},t),$$
(4.25)
the equal-time thermal Wightman function with the given thermal boundary conditions.
Fortunately, as for the condensed matter case, the interval $`\overline{t}\mathrm{\Delta }t\overline{t}`$ occurs in the linear regime, when the self-interactions are unimportant. The relevant equation for constructing the correlation functions of this one-field system is now the second-order equation
$$\frac{^2\varphi _a}{t^2}=\frac{\delta F}{\delta \varphi _a},$$
(4.26)
for $`F`$ of Eq.2.1. This is solvable in terms of the mode functions $`\chi _k^\pm (t)`$, identical for $`a=1,2`$, satisfying
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)\right]\chi _k^\pm (t)=0,$$
(4.27)
subject to $`\chi _k^\pm (t)=e^{\pm i\omega _{in}t}`$ at $`t0`$, for incident frequency $`\omega _{in}=\sqrt{𝐤^2+ϵ_0M^2}`$ and $`m^2(t)=ϵ(t)M^2`$, where $`ϵ(t)`$ is parameterised as for the TDLG equation above. This corresponds to a temperature quench from an initial state of thermal equilibrium at temperature $`T_0>T_c`$, where $`(T_0/T_c1)=ϵ_0`$. The diagonal correlation function $`G(r,t)`$ of Eq.4.17 is given as the equal-time propagator
$`G(r,t)`$ $`=`$ $`{\displaystyle d/^3ke^{i𝐤.𝐱}\chi _k^+(t)\chi _k^{}(t)C(k)}`$
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle 𝑑kk^2\frac{\mathrm{sin}kr}{kr}\chi _k^+(t)\chi _k^{}(t)C(k)},`$
where $`C(k)=\mathrm{coth}(\omega _{in}(k)/2T_0)/2\omega _{in}(k)`$ encodes the initial conditions.
An exact solution can be given in terms of Airy functions. Dimensional analysis shows that, on ignoring the k-dependence of $`C(k)`$, appropriate for large $`r`$ (or small $`k`$), $`\xi _{eq}(\overline{t})`$ of Eq.2.5 again sets the scale of the equal-time correlation function. Specifically,
$$G(r,t_0)𝑑\kappa \frac{\mathrm{sin}\kappa (r/\overline{\xi })}{\kappa (r/\overline{\xi })}F(\kappa ),$$
(4.29)
where $`F(0)=1`$ and $`F(\kappa )\kappa ^3`$ for large $`\kappa `$. Kibble’s insight is correct, at least for this case of a single (uncoupled) field.
## 5 Vortex densities do not determine correlation lengths directly
We have seen that there is no reason to disbelieve the causal arguments of Kibble for QFT and Zurek for condensed matter as to the field correlation length at the time of the transition. The excellent agreement with the $`{}_{}{}^{3}He`$ experiments also shows that, despite the very interesting simulations of Kopnin et al., this length does, indeed, characterise vortex separation for condensed matter at the time when the defects form.
However, the recent Lancaster experiment shows that this cannot always be the case. Significantly, for $`{}_{}{}^{3}He`$ the Ginzburg regime is extremely narrow, whereas for $`{}_{}{}^{4}He`$ it is very broad. In fact, the $`{}_{}{}^{4}He`$ experiments begin and end in the Ginzburg regime, where thermal fluctuations dominate. The causality arguments are too simple to accommodate these facts.
If these differences are to be visible in the formalism, it can only be through the way in which we relate vortex density to correlation length. We have already observed that the TDLG equation can be recast as the Fokker-Planck equation, whereby the ensemble averages can be understood as averaging with respect to the probability $`p_t[\mathrm{\Phi }(𝐱)]`$ that, at time $`t`$, the field takes value $`\mathrm{\Phi }(𝐱)`$. We can use these probabilities, implicit in the correlation functions, to estimate defect densities.
### 5.1 Classical defects in condensed matter
It would be foolish to estimate the probability of finding profiles like $`\mathrm{\Phi }(𝐱)`$ of Eq.3.11 directly. One way is to work through line zeroes, since vortices have line zeroes of the complex field $`\varphi `$ at the centre of their cores. The converse is not true since zeroes occur on all scales. However, a starting-point for counting vortices in superfluids is to count line zeroes of an appropriately coarse-grained field, in which structure on a scale smaller than $`\xi _0`$, the classical vortex size, is not present. That is, we do not want to entertain vortices within vortices. For the moment, we put in a cutoff $`l=O(\xi _0)`$ by hand into the Fourier transform $`G(k,t)`$ of $`G(r,t)`$, as
$$G_l(r,t)=d/^3ke^{i𝐤.𝐫}G(k,t)e^{k^2l^2}.$$
(5.30)
We stress that the long-distance correlation length $`\xi _{eq}(\overline{t})`$ depends essentially on the position of the nearest singularity of $`G(k,t)`$ in the complex k-plane, independent of $`l`$.
This is not the case for the line-zero density $`n_{zero}`$. For example, in our Gaussian approximation of the previous section $`n_{zero}`$ can be calculated exactly from the two-point correlation function $`G(r,t)`$ with $`p_t[\mathrm{\Phi }]`$ implicit. It can be shown, quite easily that it depends on the short-distance behaviour of $`G_l(r,t)`$ as
$$n_{zero}(t)=\frac{1}{2\pi }\frac{G_l^{\prime \prime }(0,t)}{G_l(0,t)},$$
(5.31)
the ratio of fourth to second moments of $`G(k,t)e^{k^2l^2}`$.
There are several prerequisites before line zeroes can be identified with vortex cores, and $`n_{zero}(t)`$ with $`n_{def}(t)`$.
* The field, on average must have achieved its symmetry-broken ground-state equilibrium value
$$|\varphi |^2=\alpha _0/\beta .$$
(5.32)
This, in itself, is sufficient to show that the causal time $`\overline{t}`$ is not the time to begin looking for defects, since $`|\varphi |^2`$ is small at this time. This, in turn, requires that $`G(k,t)`$ be non-perturbatively (in $`\beta `$) large.
* Only when $`n_{zero}/l`$ is small in comparison to $`n_{zero}/l`$ at $`l=\xi _0`$ will the line-zeroes have the non-fractal nature of classical defects on small-scales, although vortices may behave like random walks on larger scales. As the power in the long wavelength modes increases the ’Bragg’ peak develops in $`k^2G(k,t)`$, moving in towards $`k=0`$. This condition then becomes the condition that the peak dominates its tail.
* The energy in field gradients should be commensurate with the energy in classical vortices with same density as that of line zeroes.
We stress that these are necessary, but not sufficient, conditions for classical vortices. In particular, although they can be satisfied in the self-consistent linear approximation that will be outlined below, only the full nonlinearity of the system can establish classical profiles. We will term such zeroes as satisfy these conditions proto-vortices. It has to be said that most (but not all) numerical lattice simulations cannot distinguish between proto-vortices and classical vortices.
### 5.2 TDLG condensed matter
We begin with condensed matter, which we will find to be easier. As the system evolves away from the transition time, the free equation Eq.4.16 ceases to be valid, to be replaced by the full equation
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ(t)+\overline{\beta }|\varphi (𝐱,t)|^2]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t),$$
(5.33)
where $`\overline{\beta }`$ is the rescaled coupling.
In order to retain some analytic understanding of the way that the density is such an ideal quantity to make predictions for, we adopt the approximation of preserving Gaussian fluctuations by linearising the self-interaction as
$$\dot{\varphi }_a(𝐱,t)=[^2+ϵ_{eff}(t)]\varphi _a(𝐱,t)+\overline{\eta }_a(𝐱,t),$$
(5.34)
where $`ϵ_{eff}`$ contains a (self-consistent) term $`O(\overline{\beta }|\varphi |^2)`$. Additive renormalisation is necessary, so that $`ϵ_{eff}ϵ`$, as given earlier, for $`tt_0`$.
Self-consistent linearisation is the standard approximation in non-equilibrium QFT, but is not strictly necessary here, since numerical simulations that identify line zeroes of the field can be made that use the full self-interaction. However, there are none that address our particular problems exactly. Given the similarities with the QFT case, for which it is difficult to do much better than a Gaussian, there are virtues in comparing the Gaussian approximation for the two cases.
The solution for $`G(r,t)`$ is a straightforward generalisation of Eq.4.18,
$$G(r,t)=_0^{\mathrm{}}𝑑\tau \overline{T}(t\tau /2)\left(\frac{1}{4\pi \tau }\right)^{3/2}e^{r^2/4\tau }e^{_0^\tau 𝑑sϵ_{eff}(ts/2)},$$
(5.35)
where $`\overline{T}`$ is the rescaled temperature, as before.
Assuming a single zero of $`ϵ_{eff}(t)`$ at $`t=t_0`$, at $`r=0`$ the exponential in the integrand peaks at $`\tau =\overline{\tau }=2(tt_0)`$, the counterpart of the Bragg peak in proper-time. Expanding about $`\overline{\tau }`$ to quadratic order gives
$$G_l(0,t)\overline{T}_ce^{2_{t_0}^t𝑑u|ϵ_{eff}(u)|}_0^{\mathrm{}}\frac{d\tau e^{(\tau 2(tt_0))^2|ϵ^{}(t_0)|/4}}{[4\pi (\tau +\overline{l}^2)]^{3/2}},$$
(5.36)
where we have put in the momentum cutoff $`k^1>l=\overline{l}\xi _0=O(\xi _0)`$of Eq.5.30 by hand. For times $`t>ϵ_0\tau _Q`$ we see that, as the unfreezing occurs, long wavelength modes with $`k^2<t/\tau _Qϵ_0`$ grow exponentially.
The effect of the back-reaction is to stop the growth of $`G_l(0,t)G_l(0,t_0)=|\varphi |^2_t|\varphi |^2_0`$ at its symmetry-broken value $`\overline{\beta }^1`$ in our dimensionless units. A necessary condition for this is $`lim_u\mathrm{}ϵ_{eff}(u)=0`$. That is, we must choose $`ϵ_{eff}(t)=ϵ(t)+\overline{\beta }(G_l(0,t)G_l(0,t_0))`$, thereby preserving Goldstone’s theorem.
At $`t=t_0`$, when the approximation Eq.5.36 is good, both numerator and denominator are dominated by the short wavelength fluctuations at small $`\tau `$. Even though the field is correlated over a distance $`\overline{\xi }l`$ the density of line zeroes $`n_{zero}=O(l^2)`$ depends entirely on the scale at which we look. In no way would we wish to identify these line zeroes with prototype vortices. However, as time passes the peak of the exponential grows and $`n_{zero}`$ becomes increasingly insensitive to $`l`$. How much time we have depends on the magnitude of $`\overline{\beta }`$, since once $`G(0,t)`$ has reached this value it stops growing. The time $`t^{}`$ at which this happens can be estimated by substituting $`ϵ(u)`$ for $`ϵ_{eff}(u)`$ in the expression for $`G_l(0,t)`$ above.
For $`t>t^{}`$ the equation for $`n_{zero}(t)`$ is not so simple since the estimate for $`G_l(0,t)`$ above, based on a single isolated zero of $`ϵ_{eff}(t)`$, breaks down because of the approximate vanishing of $`ϵ_{eff}(t)`$ for $`t>t^{}`$. A more careful analysis shows that $`G_l(0,t)`$ can be written as
$$G_l(0,t)_0^{\mathrm{}}\frac{d\tau \overline{T}(t\tau /2)}{[4\pi (\tau +\overline{l}^2)]^{3/2}}\overline{G}(\tau ,t),$$
(5.37)
where $`\overline{G}(\tau ,t)`$ has the same peak as before at $`\tau =2(tt_0)`$, in magnitude and position, but $`\overline{G}(\tau ,t)1`$ for $`\tau <2(tt^{})`$. Thus, for $`\tau _Q\tau _0`$, $`G_l(0,t)`$ can be approximately separated as $`G_l(0,t)G_l^{UV}(t)+G^{IR}(t)`$, where
$$G_l^{UV}(t)=\overline{T}(t)_0^{\mathrm{}}𝑑\tau /[4\pi (\tau +\overline{l}^2)]^{3/2},$$
(5.38)
the counterpart of the Bragg ’tail’, describes the scale-dependent short wavelength thermal noise, proportional to temperature, and
$$G^{IR}(t)=\frac{\overline{T}_c}{(8\pi (tt_0))^{3/2}}_{\mathrm{}}^{\mathrm{}}𝑑\tau \overline{G}(\tau ,t)$$
(5.39)
describes the scale-independent, temperature independent, long wavelength fluctuations. A similar decomposition $`G_l(0,t)G_l^{UV}(t)+G^{IR}(t)`$ can be performed. In particular, $`G^{IR}(t)/G^{IR}(t)=O(t^1)`$.
Firstly, suppose that, for $`tt^{}`$, $`G^{IR}(t)G_l^{UV}(t)`$ and $`G^{IR}(t)G_l^{UV}(t)`$, as would be the case for a temperature quench $`\overline{T}(t)0`$. Then, with little thermal noise, we have widely separated line zeroes, with density $`n_{zero}(t)G^{IR}(t)/2\pi G^{IR}(t)`$ and $`n_{zero}/l`$ is small in comparison to $`n_{zero}/l`$ at $`l=\xi _0`$. Further, once the line zeroes have straightened on small scales at $`t>t^{}`$, the Gaussian field energy, largely in field gradients, is
$$\overline{F}_Vd^3x\frac{1}{2}(\varphi _a)^2=VG^{\prime \prime }(0,t),$$
(5.40)
where $`V`$ is the spatial volume. This matches the energy
$$\overline{E}Vn_{def}(t)(2\pi G(0,t))=VG^{\prime \prime }(0,t)$$
(5.41)
possessed by a network of classical global strings with density $`n_{zero}`$, in the same approximation of cutting off their logarithmic tails.
From our comments above, we identify such essentially non-fractal line-zeroes with prototype vortices, and $`n_{zero}`$ with $`n_{def}`$. Of course, we require non-Gaussianity to create true classical energy profiles. Nonetheless, the Halperin-Mazenko result may be well approximated for a while even when the fluctuations are no longer Gaussian.
For times $`t>t^{}`$
$$n_{zero}(t)\frac{\overline{t}}{8\pi (tt_0)}\frac{1}{\xi _0^2}\sqrt{\frac{\tau _0}{\tau _Q}},$$
(5.42)
the solution to Vinen’s equation
$$\frac{n_{zero}}{t}=\chi _2\frac{\mathrm{}}{m}n_{def}^2,$$
(5.43)
where $`\chi _2=4\pi `$ in our approximation<sup>7</sup><sup>7</sup>7Calculations for $`\chi _2`$ for realistic values of $`\xi _0`$ and $`\tau _0`$ give $`\chi _2>4\pi `$ for both $`{}_{}{}^{4}He`$ and $`{}_{}{}^{3}He`$. What is remarkable in this approximation is that the density of line zeroes uses no property of the self-mass contribution to $`ϵ_{eff}(t)`$, self-consistent or otherwise.
This decay law is assumed in the analysis of the Lancaster experiments. The empirical value of $`\chi _2`$ used in them is not taken from quenches, but turbulent flow experiments. It is suggested that $`\chi _20.005`$, a good three orders of magnitude smaller than our prediction above. Although the TDLG theory is not very reliable for $`{}_{}{}^{4}He`$, if our estimate is sensible it does imply that vortices produced in a temperature quench decay much faster than those produced in turbulence.
Equally importantly, we shall see that, for early time at least, thermal fluctuations are large in the Lancaster experiments. However, for $`{}_{}{}^{3}He`$, with negligable UV contributions, we estimate the primordial density of vortices as
$$n_{zero}(t^{})\frac{\overline{t}}{8\pi (t^{}t_0)}\frac{1}{\xi _0^2}\sqrt{\frac{\tau _0}{\tau _Q}},$$
(5.44)
in accord with the original prediction of Zurek. Because of the rapid growth of $`G(0,t)`$, $`(t^{}t_0)/\overline{t}=p>1=O(1)`$. We note that the factor<sup>8</sup><sup>8</sup>8An errant factor of 3 appeared in the result of of $`f^2=8\pi p`$ gives a value of $`f=O(10)`$, in agreement with the empirical results of and the numerical results of <sup>9</sup><sup>9</sup>9The temperature quench of the latter is somewhat different from that considered here, but should still give the same results in this case.
Whereas Eq.5.44 is appropriate for $`{}_{}{}^{3}He`$, the situation for the Lancaster $`{}_{}{}^{4}He`$ experiments is complex, since they are pressure quenches for which the temperature $`T`$ is almost constant at $`TT_c`$. Unlike temperature quenches, thermal fluctuations here remain at full strength<sup>10</sup><sup>10</sup>10Even for $`{}_{}{}^{3}He`$, $`T/T_c`$ never gets very small, and henceforth we take $`T=T_c`$ in $`G_l(0,t)`$ above. The necessary time-independence of $`G^{IR}(t)`$ for $`t>t^{}`$ is achieved by taking $`ϵ_{eff}(u)=O(u^1)`$. In consequence, as $`t`$ increases beyond $`t^{}`$ the relative magnitude of the UV and IR contributions to $`G_l(0,t)`$ remains approximately constant at its value at $`t=t^{}`$.
Nonetheless, as long as the UV fluctuations are insignificant at $`t=t^{}`$ the density of line zeroes will remain largely independent of scale. This follows if $`G^{IR}(t^{})G_l^{UV}(t^{})`$, since $`G_l(0,t)`$ becomes scale-independent later than $`G_l(0,t)`$. In we showed that this is true provided
$$(\tau _Q/\tau _0)(1T_G/T_c)<C\pi ^4,$$
(5.45)
where $`C=O(1)`$ and $`T_G`$ is the Ginzburg temperature. With $`\tau _Q/\tau _0=O(10^3)`$ and $`(1T_G/T_c)=O(10^{12})`$ this inequality is well satisfied for a linearised TDLG theory for $`{}_{}{}^{3}He`$ derived<sup>11</sup><sup>11</sup>11Ignoring the position-dependent temperature of from the full TDGL theory, but there is no way that it can be satisfied for $`{}_{}{}^{4}He`$, when subjected to a slow mechanical quench, as in the Lancaster experiment, for which $`\tau _Q/\tau _0=O(10^{10})`$, since the Ginzburg regime is so large that $`(1T_G/T_c)=O(1)`$. As far as the left hand side of Eq.5.45 is concerned, the $`{}_{}{}^{4}He`$ quench is nineteen orders of magnitude slower than its $`{}_{}{}^{3}He`$ counterpart.
When the inequality is badly violated, as with $`{}_{}{}^{4}He`$ for slow pressure quenches, then the density of zeroes $`n_{def}=O(l^2)`$ after $`t^{}`$ again depends explicitly on the scale $`l`$ at which we look and they are not candidates for vortices. Since the whole of the quench takes place within the Ginzburg regime this is not implausible. However, it is possible that, even though the thermal noise never switches off, there is no more than a postponement of vortex production, since our approximations must break down at some stage. The best outcome is to assume that the effect of the thermal fluctuations on fractal behaviour is diminished, only leading to a delay in the time at which vortices finally appear. Even if we suppose that $`n_{def}`$ above is a starting point for calculating the density at later times, albeit with a different $`t_0`$, thereby preserving Vinen’s law, we then have the earlier problem of the large $`\chi _2=O(f^2)`$, which would make it almost impossible to see vortices.
For all that, a numerical simulation that goes beyond the Gaussian approximation specifically tailored to the Lancaster parameters is crucial if we are to understand what is really happening. We hope to pursue this elsewhere.
## 6 The Appearance of Structure in QFT
When, in Section 5.2 we set up the closed time-path formalism for the field probabilities $`p_t[\mathrm{\Phi }]`$, our aim was the limited one of establishing the role of Kibble’s causal correlation length $`\overline{\xi }`$ in Eq.4.29. We now appreciate, from condensed matter theory, that this does not, of itself, imply vortices at that separation.
### 6.1 Proto-vortices in QFT
To establish a link between the correlation function $`G(r,t)`$ and vortices is even more problematic in QFT than for condensed matter systems. Yet again, we attempt to count vortices by counting line zeroes. In the Gaussian approximations that we shall continue to adopt the expression Eq.5.31 for $`n_{zero}`$ is equally applicable to QFT. This counting of zeroes is the basis of numerous numerical simulations of cosmic string networks built from Gaussian fluctuations.
The prerequisites for line zeroes in condensed matter that we posed after Eq.5.31 still stand for QFT (except that $`|\varphi ^2|=M^2/\lambda `$), but there are further complications peculiar to QFT. In particular, in QFT we need to consider the whole density matrix $`\mathrm{\Phi }^{}|\rho (t)|\mathrm{\Phi }`$ rather than just the diagonal elements $`p_t[\mathrm{\Phi }]=\mathrm{\Phi }|\rho (t)|\mathrm{\Phi }`$. Classicality is understood in terms of ’decoherence’, manifest most simply by the approximate diagonalisation of the reduced density matrix on coarse-graining. By this we mean the separation of the whole into the ’system’, and its ’environment’ whose degrees of freedom are integrated over, to give a reduced density matrix. The environment can be either other fields with which our scalar is interacting or even the short wavelength modes of the scalar field itself . When interactions are taken into account this leads to quantum noise and dissipation.
In the Gaussian approximations that we shall adopt here, with $`\mathrm{\Phi }=0`$, integrating out short wavelengths with $`k>l^1`$ is just equivalent to a momentum cut-off at the same value. This gives neither noise nor dissipation and diagonalisation does not occur. Nonetheless, from our viewpoint of counting line-zeroes, fluctuations are still present when $`l=O(M^1)`$ that can prevent us from identifying line-zeroes with proto-vortices, if the quenches are too slow.
For all these caveats, there are other symptoms of classical behaviour once $`G_l(0;t)`$ is non-perturbatively large. Instead of a field basis, we can work in a particle basis and measure the particle production as the transition proceeds. The presence of a non-perturbatively large peak in $`k^2G(k;t)`$ at $`k=k_0`$ signals a non-perturbatively large occupation number $`N_{k_0}1/\lambda `$ of particles at the same wavenumber $`k_0`$. With $`n_{zero}`$ of (5.31) of order $`k_0^2`$ this shows that the long wavelength modes can now begin to be treated classically. From a slightly different viewpoint, the Wigner functional only peaks about the classical phase-space trajectory once the power is non-perturbatively large. More crudely, the diagonal density matrix elements are only then significantly non-zero for non-perturbatively large field configurations $`\varphi \lambda ^{1/2}`$, like vortices.
### 6.2 Mode growth v. fluctuations
For early times we revert to the mode decomposition of Eq.4.27. The term $`coth(\omega _{in}/2T_0)`$ appearing in it can be approximated by $`2T_0/\sqrt{ϵ_0}M`$. Even though this is a temperature quench, it shows strong similarities to the pressure quench of condensed matter, since both the long and short wavelength contributions to $`G(r,t)`$ are scaled by the same temperature and we cannot switch off the latter.
The field becomes ordered, as before, because of the exponential growth of long-wavelength modes, which stop growing once the field has sampled the groundstates. What matters is the relative weight of these modes (the ’Bragg’ peak) to the fluctuating short wavelength modes in the decomposition Eq.4.2 at this time, since the contribution of these latter is very sensitive to the cutoff $`l`$. Only if their contribution to Eq.3.10 is small when field growth stops can a network of line-zeroes be well-defined at early times, let alone have the predicted density. Since the peak is non-perturbatively large this requires small coupling, which we assume.
Consider a quench with $`ϵ(t)`$ as in Eq.4.13, in which the symmetry-breaking begins at relative time $`\mathrm{\Delta }t=tt_0=0`$. For a free roll, the exponentially growing modes that appear when $`\mathrm{\Delta }t>t_k^{}=t_Qk^2/M^2`$ lead to the approximate WKB solution
$$G(r;\mathrm{\Delta }t)\frac{T}{M|m(\mathrm{\Delta }t)|}\left(\frac{M}{\sqrt{\mathrm{\Delta }tt_Q}}\right)^{3/2}e^{\frac{4M\mathrm{\Delta }t^{3/2}}{3\sqrt{t_Q}}}e^{r^2/\xi ^2(\mathrm{\Delta }t)}$$
(6.46)
where $`\xi ^2(\mathrm{\Delta }t)=2\sqrt{\mathrm{\Delta }tt_Q}/M`$. The provisional freeze-in time $`t_{}`$ when $`|\varphi ^2|=M^2/\lambda `$ is then, for $`Mt_Q<(1/\lambda )`$,
$$M\mathrm{\Delta }t_{}(Mt_Q)^{1/3}(\mathrm{ln}(1/\lambda ))^{2/3}M\overline{t}(\mathrm{ln}(1/\lambda ))^{2/3},$$
(6.47)
where $`\mathrm{\Delta }t_{}=t_{}t_0`$. This is greater than $`M\overline{t}`$, but not by a large multiple. Comparison with condensed matter, for which the ratio is a few ($`35`$) suggests that we don’t need a superweak theory.
At this qualitative level the correlation length at $`t_{}`$ is
$$M^2\xi ^2(t_{})(Mt_Q)^{2/3}(\mathrm{ln}(1/\lambda ))^{1/3}.$$
(6.48)
The effect of the other modes is larger than for the instantaneous quench, giving, at $`t=t_{}`$
$$n_{zero}=\frac{M^2}{\pi (M\tau _Q)^{2/3}}(\mathrm{ln}(1/\lambda ))^{1/3}[1+E].$$
(6.49)
The error term $`E=O(\lambda ^{1/2}(Mt_Q)^{4/3}(\mathrm{ln}(1/\lambda ))^{1/3})`$ is due to oscillatory modes, sensitive to the cutoff. In mimicry of Eq.3.10 it is helpful to rewrite Eq.6.49 as
$$n_{zero}=\left[\frac{1}{\pi \xi _0^2}\left(\frac{\tau _0}{\tau _Q}\right)^{2/3}\right](\mathrm{ln}(1/\lambda ))^{1/3}[1+E].$$
(6.50)
in terms of the scales $`\tau _0=\xi _0=M^1`$. The first term in Eq.6.50 is the Kibble estimate of Eq.3.10, the second is the small multiplying factor, that yet again shows that estimate can be correct, but for completely different reasons. The third term shows when it can be correct, since $`E`$ is also a measure of the sensitivity of $`n_{def}`$ to the scale at which it is measured. The condition $`E^21`$, necessary for a proto-vortex network to be defined, is then guaranteed if
$$(\tau _Q/\tau _0)^2(1T_G/T_c)<C,$$
(6.51)
$`C=O(1)`$, on using the relation $`(1T_G/T_c)=O(\lambda )`$. This is the QFT counterpart to Eq.5.45.
For example, suppose that this approach is relevant to the local strings of a strong Type-II $`U(1)`$ theory for the early universe, in which the time-temperature relationship $`tT^2=\mathrm{\Gamma }M_{pl}`$ is valid, where we take $`\mathrm{\Gamma }=O(10^1)`$ in the GUT era. If $`G`$ is Newton’s constant and $`\mu `$ the classical string tension then, following , $`Mt_Q10^1\lambda ^{1/2}(G\mu )^{1/2}`$. The dimensionless quantity $`G\mu 10^610^7`$ is the small parameter of cosmic string theory. A value $`\lambda 10^2`$ gives $`Mt_Q(Mt_{})^a,a2`$, once factors of $`\pi `$, etc.are taken into account, rather than $`Mt_Q1/\lambda `$, and the density of Eq.6.50 may be relevant.
### 6.3 Backreaction in QFT
To improve upon the free-roll result more honestly, but retain the Gaussian approximation for the field correlation functions, the best we can do is adopt a mean-field approximation along the lines of , as we did for the CM systems earlier. As there, it does have the correct behaviour of stopping domain growth as the field spreads to the potential minima. As before, only the large-$`N`$ expansion preserves Goldstone’s theorem.
$`G(r;t)`$ still has the mode decomposition of Eq.4.2, but the modes $`\chi _k^\pm `$ now satisfy the equation
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)+\lambda \mathrm{\Phi }^2(\mathrm{𝟎})_t\right]\chi _k^\pm (t)=0,$$
(6.52)
where we have taken $`N=2`$. Because $`\lambda \varphi ^4`$ theory is not asymptotically free, particularly in the Hartree approximation, the renormalised $`\lambda `$ coupling shows a Landau ghost. This means that the theory can only be taken as a low energy effective theory.
The end result is, on making a single subtraction at $`t=0`$, is
$$\left[\frac{d^2}{dt^2}+𝐤^2+m^2(t)+\lambda d/^3pC(p)[\chi _p^+(t)\chi _p^{}(t)1]\right]\chi _k^\pm (t)=0.$$
(6.53)
which we write as
$$\left[\frac{d^2}{dt^2}+𝐤^2\mu ^2(t)\right]\chi _k(t)=0.$$
(6.54)
On keeping just the unstable modes in $`\mathrm{\Phi }^2(\mathrm{𝟎})_t`$ then, as it grows, its contribution to (6.53) weakens the instabilities, so that only longer wavelengths become unstable. At $`t^{}`$ the instabilities shut off, by definition, and oscillatory behaviour ensues. Since the mode with wavenumber $`k>0`$ stops growing at time $`t_k^+<t^{}`$, where $`\mu ^2(t_k^+)=𝐤^2`$, the free-roll density at $`t^{}`$ must be an overestimate.
An approximation that improves upon the WKB approximation is
$$\chi _k(t)\left(\frac{\pi M}{2\mathrm{\Omega }_k(\eta )}\right)^{1/2}\mathrm{exp}\left(_0^t𝑑t\mathrm{\Omega }(t)\right)$$
(6.55)
when $`\eta =M(t_k^+t)>0`$ is large, and $`\mathrm{\Omega }_k^2(t)=\mu ^2(t)𝐤^2`$. On expanding the exponent in powers of $`k`$ and retaining only the quadratic terms we recover the WKB approximation when $`\mu (t)`$ is non-zero.
The result is that the effect of the back-reaction is to give a time-delay $`\mathrm{\Delta }t`$ to $`t^{}`$, corresponding to a decrease in the value $`k_0(t)`$ at which the power peaks of order
$$\frac{\mathrm{\Delta }t}{t^{}}=O\left(\frac{1}{ln(1/\lambda )}\right).$$
(6.56)
The backreaction has little effect for times $`t<t^{}`$. For $`t>t^{}`$ oscillatory modes take over the correlation function and we expect oscillations in $`G(k;t)`$.
In practice the backreaction rapidly forces $`\mu ^2(t)`$ towards zero if the coupling is not too small. For couplings that are not too weak, this requires that we graft purely oscillatory long wavelength behaviour onto the non-perturbatively large exponential mode
$$\chi _k^+(t^{})\alpha _k\mathrm{exp}\left(_0^t^{}𝑑t^{}\mu (t^{})\right)\mathrm{exp}\left(\frac{\sqrt{\tau _Qt^{}}}{M}k^2\right)$$
(6.57)
The end result is a new power spectrum, obtained by superimposing oscillatory behaviour onto the old spectrum, frozen at time $`t^{}`$. As a gross oversimplification, the contribution from the earlier exponential modes alone can only be to contribute terms something like
$`G(r;t)`$ $``$ $`{\displaystyle \frac{T}{M|m(t^{})|}}e^{4M(t^{})^{3/2}/3\sqrt{\tau _Q}}{\displaystyle _{|𝐤|<M}}d/^3ke^{i𝐤.𝐱}e^{2\sqrt{t^{}\tau _Q}k^2/M}`$ (6.58)
$`\times `$ $`\left[\mathrm{cos}k(tt^{})+{\displaystyle \frac{\mathrm{\Omega }(k)W^{}(k)}{k}}\mathrm{sin}k(tt^{})\right]^2`$
to $`G`$, where $`\mathrm{\Omega }=M(t^{}t_k)^{1/2}/\tau _Q^{1/2}`$ and $`W^{}=1/4(t^{}t_k)`$. The details are almost irrelevant, since the density of line zeroes is independent of the normalisation, and only weakly dependent on the power spectrum.
The $`k=0`$ mode of Eq.6.58 encodes the simple solution $`\chi _{k=0}(t)=a+bt`$ when $`\mu ^2=0`$. As observed by Boyanovsky et al. this has built into it the basic causality discussed by Kibble. Specifically, for $`r,t\mathrm{}`$, but $`r/2t`$ constant $`(1)`$,
$$G(r,t)\frac{C}{r}\mathrm{\Theta }(2t/r1).$$
(6.59)
It follows directly that this causality, engendered by the Goldstone particles of the self-consistent theory, has little effect on the density of line-zeroes that we expect to mature into fully classical vortices, since that is determined by the behaviour at $`r=0`$.
Further, for large $`t`$ the power spectrum effectively has a $`k^2`$ behaviour for small $`k`$, unlike the white noise that would follow from Eq.6.46. It has been suggested that, for such a spectrum, most, if not all, of the vortices are in loops, with little or no self-avoiding ’infinite’ string (but see ). If there was no infinite string the evolution of the network could be very different from that of white noise, where approximately $`75\%`$ of the string is ’infinite’. Although causality due to massless Goldstone modes is unrealistic, the linking of causal behaviour to the long wavelength spectrum is general. It has to be said that this approximation should not be taken very seriously for large $`t`$ on different grounds, since we would expect rescattering to take place at times $`\mathrm{\Delta }t=O(1/\lambda )`$ in a way that is precluded by the Gaussian approximation.
Returning to our original concerns, if Eq.6.51 is not satisfied, it is difficult to imagine how clean vortices, or proto-vortices, can appear later without some additional ingredient.
## 7 Conclusions
We examined the Kibble /Zurek predictions for the onset of phase transitions and the appearance of defects (in particular, vortices or global cosmic strings) as a signal of the symmetry breaking. Our results are in agreement with their prediction Eq.3.9 as to the magnitude of the correlation length at the time the transition truly begins, equally true for condensed matter and QFT.
However, this is not simply a measure of the separation of defects at the time of their appearance. The time $`\overline{t}`$ is too early for the field to have found the true groundstates of the theory. We believe that time, essentially the spinodal time, is the time at which proto-vortices can appear, which can later evolve into the standard classical vortices of the theory.
Even then, they may not appear because of thermal field fluctuations. In TDLG condensed matter thermal noise is proportional to temperature. If temperature is fixed, but not otherwise, as in the pressure quenches of $`{}_{}{}^{4}He`$, this noise can inhibit the production of vortices, although there are other factors to be taken into account (such as their decay rate). On the other hand, on quenching from a high temperature in QFT there are always thermal fluctuations, and these can also disturb the appearance of vortices. The condition that thermal fluctuations are ignorable at the time that the field has achieved the true groundstates can be written
$$(\tau _Q/\tau _0)^\gamma (1T_G/T_c)<C,$$
(7.60)
where $`\gamma =1`$ for condensed matter and $`\gamma =2`$ for QFT. $`C=O(1)`$.
This restores the role of the Ginzburg temperature $`T_G`$ that the simple causal arguments overlooked, but does not restore thermal fluctuations as the exclusive agent for vortex production, as happened in early arguments. Quenches in $`{}_{}{}^{4}He`$ provide the major example for which Eq.7.60 is not satisfied.
What happens at late time is unclear, although for TDLG numerical simulations can be performed (but have yet to address this problem exactly). On the other hand, not only is the case of a single self-interacting quantum scalar field in flat space-time a caricature of the early universe, but it is extremely difficult to go beyond the Gaussian approximation. To do better requires that we do differently. There are several possible approaches. One step is to take the FRW metric of the early universe seriously, whereby the dissipation due to the expansion of the universe can change the situation dramatically. Other approaches are more explicit in their attempts to trigger decoherence explicitly, as we mentioned earlier. Most simply, the short wavelength parts of the field can be treated as an environment to be integrated over, to give a coarse-grained theory of long-wavelength modes acting classically in the presence of noise. However, such noise is more complicated than in TDLG theory, being multiplicative as well as additive, and coloured. This is an area to be pursued elsewhere.
I particularly thank Glykeria Karra and Eleftheria Kavoussanaki, with whom much of this work was done. I also would like to thank Tom Kibble and many of my colleagues working in this area for fruitful discussions. This is an area with a substantial, if scattered, literature and I have aimed to be exemplary, rather than inclusive. I apologise to any authors who I have not cited explicitly. This work is the result of a network supported by the European Science Foundation. |
warning/0001/math0001003.html | ar5iv | text | # \quad1.1. Definition
NEW MODULI SPACES OF POINTED CURVES
AND PENCILS OF FLAT CONNECTIONS
A. Losev<sup>1</sup>, Yu. Manin<sup>2</sup>
<sup>1</sup>Institute of Theoretical and Experimental Physics, Moscow, Russia
<sup>2</sup>Max–Planck–Institut für Mathematik, Bonn, Germany
Dedicated to William Fulton on the occasion of his 60th birthday
Abstract. It is well known that formal solutions to the Associativity Equations are the same as cyclic algebras over the homology operad $`(H_{}(\overline{M}_{0,n+1}))`$ of the moduli spaces of $`n`$–pointed stable curves of genus zero. In this paper we establish a similar relationship between the pencils of formal flat connections (or solutions to the Commutativity Equations) and homology of a new series $`\overline{L}_n`$ of pointed stable curves of genus zero. Whereas $`\overline{M}_{0,n+1}`$ parametrizes trees of $`^1`$’s with pairwise distinct nonsingular marked points, $`\overline{L}_n`$ parametrizes strings of $`^1`$’s stabilized by marked points of two types. The union of all $`\overline{L}_n`$’s forms a semigroup rather than operad, and the role of operadic algebras is taken over by the representations of the appropriately twisted homology algebra of this union.
0. Introduction and plan of the paper
One of the remarkable basic results in the theory of the Associativity Equations (or Frobenius manifolds) is the fact that their formal solutions are the same as cyclic algebras over the homology operad $`(H_{}(\overline{M}_{0,n+1}))`$ of the moduli spaces of $`n`$–pointed stable curves of genus zero. This connection was discovered by physicists, who observed that the data of both types come from models of topological string theories. Precise mathematical treatment was given in \[KM\] and \[KMK\].
In this paper we establish a similar relationship between the pencils of formal flat connections (or solutions to the Commutativity Equations: see 3.1–3.2 below) and homology of a new series $`\overline{L}_n`$ of pointed stable curves of genus zero. Whereas $`\overline{M}_{0,n+1}`$ parametrizes trees of $`^1`$’s with pairwise distinct nonsingular marked points, $`\overline{L}_n`$ parametrizes strings of $`^1`$’s, and all marked points with exception of two are allowed to coincide (see the precise definitions in 1.1 and 2.1). Moreover, the union of all $`\overline{L}_n`$’s forms a semigroup rather than operad, and the role of operadic algebras is taken over by the representations of the appropriately twisted homology algebra of this union: see precise definitions in 3.3.
This relationship was discovered on a physical level in \[Lo1\], \[Lo2\]. Here we give a mathematical treatment of some of the main issues raised in these papers.
This paper is structured as follows.
In §1 we introduce the notion of $`(A,B)`$–pointed curves whose combinatorial structure generalizes that of strings of projective lines described above. We then describe a construction of “adjoining a generic black point” which allows us to produce families of such curves and their moduli stacks inductively. This is a simple variation of one of the arguments due to F. Knudsen .
In §2 we define and study the spaces $`\overline{L}_n`$ for which we give two complementary constructions. The first one identifies $`\overline{L}_n`$ with one of the moduli spaces of pointed curves. The second one exhibits $`\overline{L}_n`$ as a well–known toric manifold associated with the polytope called permutohedron in \[Ka2\]. These constructions put $`\overline{L}_n`$ into two quite different contexts and suggest generalizations in different directions.
As moduli spaces, $`\overline{L}_n`$ become components of the extended modular operad which we define and briefly discuss in §4. We expect that there exists an appropriate extension of the Gromov–Witten invariants producing algebras over extended operads involving gravitational descendants.
As toric varieties, $`(\overline{L}_n)`$ form one of the several series related to the generalized flag spaces of classical groups: see \[GeSe\]. It would be interesting to generalize to other series our constructions.
In this paper we use the toric description in order to prove for $`\overline{L}_n`$’s an analog of Keel’s theorem (Theorem 2.7.1) and its extension (Theorem 2.9), crucial for studying representations of the twisted homology algebra.
This twisted homology algebra $`H_{}T`$ and its relationship with pencils of formal flat connections are discussed in §3, which contains the main result of this paper: Theorem 3.3.1.
Acknowledgement. Yu. Manin is grateful to M. Kapranov who, after having seen the formula $`\chi (\overline{L}_n)=n!`$, suggested that $`\overline{L}_n`$ must be the toric variety associated with the permutohedron.
§1. $`(A,B)`$–pointed curves
###### \quad1.1. Definition
Let $`A,B`$ be two finite disjoint sets, $`S`$ a scheme, $`g0`$. An $`(A,B)`$–pointed curve of genus $`g`$ over $`S`$ consists of the data
$$(\pi :CS;x_i:SC,iA;x_j:SC,jB)$$
$`(1.1)`$
where
(i) $`\pi `$ is a flat proper morphism whose geometric fibres $`C_s`$ are reduced and connected curves, with at most ordinary double points as singularities, and $`g=H^1(C_s,𝒪_{C_s}).`$
(ii) $`x_i,iAB,`$ are sections of $`\pi `$ not containing singular points of geometric fibres.
(iii) $`x_ix_j=\mathrm{}`$ if $`iA,jAB,ij.`$
Such a curve $`(1.1)`$ is called stable, if the normalization of any irreducible component $`C^{}`$ of a geometric fibre carries $`3`$ pairwise different special points if $`C^{}`$ is of genus $`0`$ and $`1`$ special points if $`C^{}`$ is of genus $`1.`$ Special points are inverse images of singular points and of the structure sections $`x_i`$.
1.2. Remarks. a) If we put in this definition $`B=\mathrm{}`$, we will get the usual notion of an $`A`$–pointed (pre)stable curve whose structure sections are not allowed to intersect pairwise. Now we divide the sections into two groups: “white” sections $`x_i,iA`$ are not allowed to intersect any other section, whereas “black” sections $`x_j,jB`$ cannot intersect white ones, but otherwise are free and can even pairwise coincide. (However, both types of sections are not allowed to intersect singularities of fibres).
If we take in this definition one–element set $`B=\{\}`$, we will get a natural bijection between $`(A,\{\})`$–pointed curves and $`(A\{\},\mathrm{})`$–pointed curves. If $`\mathrm{card}B2,`$ the two notions become essentially different.
b) The dual modular graph of a geometric fibre is defined in the same way as in the usual case (for the conventions we use see \[Ma\], III.2). Tails now can be of two types, and we may refer to them and their marks as “black” and “white” ones as well. Combinatorial type of a geometric fiber is, by definition, the isomorphism class of the respective modular graph with $`(A,B)`$–marking of its tails.
c) Let $`TS`$ be an arbitrary base change. It produces from any $`(A,B)`$–pointed (stable) curve (1.2) over $`S`$ another $`(A,B)`$–pointed (stable) curve over $`T`$: $`(C_T;x_{i,T}).`$
1.3. A construction. In this subsection, we start with an $`(A,B)`$–pointed curve (1.1) and produce from it another $`(A,B^{})`$–pointed curve:
$$(\pi ^{}:C^{}S^{};x_i^{},iAB^{}).$$
$`(1.2)`$
The base of the new curve will be $`S^{}:=C.`$ There will be one extra black mark, say, $``$, so that $`B^{}=B\{\}`$. The new curve and sections will be produced in two steps. At the first step we make the base change $`CS`$ as in 1.2 c), obtaining an $`(A,B)`$–pointed curve $`X:=C\times _SC`$, with sections $`x_{i,C}.`$ We then add the extra section $`\mathrm{\Delta }:CC\times _SC`$ which is the relative diagonal, and mark it by $``$. We did not yet produce an $`(A,B^{})`$–pointed curve over $`S^{}=C`$, because the extra black section can (and generally will) intersect both singular points of the fibres and white sections as well.
At the second step of the construction, we remedy this by birationally modifying $`C\times _SCC`$ as in \[Kn1\], Definition 2.3. More precisely, we define $`C^{}:=\mathrm{Proj}\mathrm{Sym}𝒦`$ as the relative projective spectrum of the symmetric algebra of the sheaf $`𝒦`$ on $`X=C\times _SC`$ defined as the cokernel of the map
$$\delta :𝒪_X𝒥_\mathrm{\Delta }\stackrel{ˇ}{}𝒪_X(\underset{iA}{}x_{i,C}),\delta (t)=(t,t).$$
$`(1.3)`$
Here $`𝒥_\mathrm{\Delta }`$ is the $`𝒪_X`$–ideal of $`\mathrm{\Delta }`$, and $`𝒥_\mathrm{\Delta }\stackrel{ˇ}{}`$ is its dual sheaf considered as a subsheaf of meromorphic functions, as in \[Kn1\], Lemma 2.2 and Appendix.
We claim now that we get an $`(A,B^{})`$–pointed curve, because Knudsen’s treatment of his modification can be directly extended to our case. In fact, the modification we described is nontrivial only in a neighbourhood of those points, where $`\mathrm{\Delta }`$ intersects either singular points of the fibres, or $`A`$–sections. The $`B`$–sections do not intersect these neighborhoods, if they are small enough, and do not influence the local analysis due to Knudsen (\[Kn1\], pp. 176–178).
1.3.1. Remark. We can try to modify this construction in order to be able to add an extra white point, instead of a black one. However, for $`\mathrm{card}B2,`$ we will not be able then to avoid the local analysis of the situation by referring to \[Kn1\]. In fact, points where $`\mathrm{\Delta }`$ intersects at least two $`B`$–sections simultaneously, will have to be treated anew.
§2. Spaces $`\overline{L}_n`$
2.1. Spaces $`\overline{L}_n`$. In this subsection we will inductively define for any $`n1`$ the $`(\{0,\mathrm{}\},\{1,\mathrm{},n\})`$–pointed stable curve of genus zero
$$(\pi _n:C_n\overline{L}_n;x_0^{(n)},x_{\mathrm{}}^{(n)};x_1^{(n)},\mathrm{},x_n^{(n)}).$$
$`(2.1)`$
Namely, put
$$C_1:=^1,\overline{L}_1=\mathrm{a}\mathrm{point},$$
and choose for $`x_0^{(1)},x_{\mathrm{}}^{(1)},x_1^{(1)}`$ arbitrary pairwise distinct points.
If (2.1) is already constructed, we define the next family $`(C_{n+1}\overline{L}_{n+1},\mathrm{})`$ as the result of the application of the construction 1.3 to $`C_n/\overline{L}_n.`$ In particular, we have a canonical isomorphism $`C_n=\overline{L}_{n+1}.`$
###### \quad2.2. Theorem
a) $`\overline{L}_n`$ is a smooth separated irreducible proper manifold of dimension $`n1.`$ It represents the functor which associates with every scheme $`T`$ the set of the isomorphism classes of $`(\{0,\mathrm{}\},\{1,\mathrm{},n\})`$–pointed stable curves of genus zero over $`T`$ whose geometric fibers have combinatorial types described below.
The symmetric group $`𝕊_n`$ renumbering the structure sections acts naturally and compatibly on $`\overline{L}_n`$ and the universal curve. In particular, we can define the spaces $`\overline{L}_B,C_B`$ for any finite set $`B`$, functorial with respect to the bijections of the sets.
b) Combinatorial types of geometric fibres of $`C_n\overline{L}_n`$ are in a natural bijection with ordered partitions
$$\{1,\mathrm{},n\}=\sigma _1\mathrm{}\sigma _l,1ln,\sigma _i\mathrm{}.$$
$`(2.2)`$
Partition (2.2) corresponds to the linear graph with vertices $`(v_1,\mathrm{},v_l)`$ of genus zero, edges joining $`(v_i,v_{i+1}),1il1`$, $`A`$–tail $`0`$ at the vertex $`v_1`$, $`A`$–tail $`\mathrm{}`$ at the vertex $`v_l`$, and $`B`$–tails marked by the elements of $`\sigma _i`$ at the vertex $`v_i`$.
We will call $`l=l(\sigma )`$ the length of the partition $`\sigma `$ as in (2.2).
c) Denote by $`L_\sigma `$ the set of all points of $`\overline{L}_n`$ corresponding to the curves of the combinatorial type $`\sigma `$, and by $`\overline{L}_\sigma `$ its Zariski closure. Then $`L_\sigma `$ are locally closed subsets, and we have
$$\overline{L}_\sigma =\underset{\tau \sigma }{}L_\tau $$
$`(2.3)`$
where $`\tau \sigma `$ means that $`\tau `$ is obtained from $`\sigma `$ by replacing each $`\sigma _i`$ by an ordered partition of $`\sigma _i`$ into non–empty subsets.
d) For every $`\sigma `$, there exists a natural isomorphism
$$L_{|\sigma _1|}\times \mathrm{}\times L_{|\sigma _l|}L_\sigma $$
$`(2.4)`$
such that the pointed curve induced by this isomorphism over $`L_{|\sigma _1|}\times \mathrm{}\times L_{|\sigma _l|}`$ can be obtained by clutching the curves $`C_{|\sigma _i|}/L_{|\sigma _i|}`$ in an obvious linear order ($`\mathrm{}`$–section of the $`i`$–th curve is identified with the $`0`$–section of the $`(i+1)`$–th curve, see \[Kn1\], Theorem 3.4), and subsequent remarking of the $`B`$–sections.
In particular, $`L_\sigma `$ is a smooth irreducible submanifold of codimension $`l(\sigma )1.`$
The similar statements hold for the closed strata $`\overline{L}_\sigma .`$
Proof. Properness and smoothness follow by induction and Knudsen’s local analysis which we already invoked.
The statement about the combinatorial types is proved by induction as well. In fact, if everything is already proved for $`C_n`$, then we must look at a geometric fibre $`C_{n,s}`$ of $`C_n`$ and see what happens to it after the blow up described in 1.3. If $`\mathrm{\Delta }`$ intersects a smooth point of $`C_{n,s}`$, not coinciding with $`x_{0,s},x_{\mathrm{},s}`$, nothing happens, except that we get a new black point on this fibre, and a new tail at the respective vertex of the dual graph. If $`\mathrm{\Delta }`$ intersects an intersection point of two neighboring components of $`C_{n,s}`$, then after blowing up these two components become disjoint, and we get a new component intersecting both of them, with a new black point on it. The linear structure of the graph is preserved. Finally, if $`\mathrm{\Delta }`$ intersects $`C_{n,s}`$ at $`x_{0,s}`$ or $`x_{\mathrm{},s}`$, then after blowing up we will get a new end component, with $`x_{0,s}`$, resp. $`x_{\mathrm{},s}`$ and the new black point on it. Thus the new combinatorial types will be linear and indexed by partitions of $`(n+1)`$. To check that all partitions are obtained in this way, it suffices to remark that $`\mathrm{\Delta }`$, being the relative diagonal, can intersect the fibre of a given type at any point.
In order to check the statement about the functor represented by $`\overline{L}_n`$ we apply the following inductive reasoning. For $`n=1`$ the statement is almost obvious. In fact, let $`\pi :CS`$ be a $`(\{0,\mathrm{}\},\{1\})`$–pointed stable curve of genus zero over $`T`$. From the stability it follows that all geometric fibres are projective lines. Since the three structure sections pairwise do not intersect, the family can be identified with $`^1\times T`$ endowed with three constant sections. This means that it is induced by the trivial morphism $`T\overline{L}_1`$.
Assume that the statement is true for $`n`$. In order to prove it for $`n+1`$, consider a $`(\{0,\mathrm{}\},\{1,\mathrm{},n+1\})`$–pointed stable curve of genus zero $`\pi :CT.`$ First of all, one can produce from it a $`(\{0,\mathrm{}\},\{1,\mathrm{},n\})`$–pointed stable curve of genus zero $`\pi :C^{}T`$ obtained by forgetting $`x_{n+1}`$ and subsequent stabilization. The respective map $`C^{}C`$ is given by the relative projective spectrum of the algebra $`_{k=0}^{\mathrm{}}\pi _{}(𝒦^k)`$ where $`𝒦:=\omega _{C/T}(x_0+x_1+\mathrm{}+x_n+x_{\mathrm{}})`$. By induction, $`C^{}`$ is induced by a morphism $`p:T\overline{L}_n`$. Addition of an extra black section to $`C^{}`$ and subsequent stabilization boils down exactly to the construction 1.3 applied to $`C^{}/T`$ which allows us to lift $`p`$ to a unique morphism $`q:T\overline{L}_{n+1}.`$
Separatedness is checked by the standard deformation arguments.
The statement about renumbering follows from the description of the functor.
A similar adaptation of Knudsen’s arguments allows us to prove the remaining statements, and we leave them to the reader.
Notice that below we will give another direct description of the spaces $`\overline{L}_B`$ and all the structure morphisms connecting them in terms of toric geometry. This will provide easy alternate proofs of their properties. Except for §4, we can restrict ourselves to this alternate description.
2.2.1. Remark. Dual graphs of the degenerate fibers of $`C_n`$ over $`\overline{L}_n`$ come with a natural orientation from $`x_0`$ to $`x_{\mathrm{}}`$. We could have allowed ourselves not to distinguish between the two white points, interchanging them by isomorphisms, but this would produce several upleasant consequences. First, our manifolds would become actual stacks, starting already with $`\overline{L}_1`$. Second, we would have lost the toric interpretation of these spaces. Third, and most important, we would meet an ambiguity in the definition of the multiplication between the homology spaces: see (3.5) below. With our choice, we can simply introduce the involution permuting $`x_0`$ and $`x_{\mathrm{}}`$ as a part of the structure and look how it interacts with other parts.
###### \bf\quad2.3. Theorem
$`\overline{L}_n`$ has no odd cohomology. Let
$$p_n(q):=\underset{i=0}{\overset{n1}{}}\mathrm{dim}H^{2i}(\overline{L}_n)q^i$$
$`(2.5)`$
be the Poincaré polynomial of $`\overline{L}_n`$. Then we have
$$1+\underset{n=1}{\overset{\mathrm{}}{}}\frac{p_n(q)}{n!}y^n=\frac{q1}{qe^{(q1)y}}[q][[y]].$$
$`(2.6)`$
Letting here $`q1`$ we get $`{\displaystyle \frac{1}{1y}}`$ so that $`\chi (\overline{L})=n!.`$
Proof. Since $`\overline{L}_n`$ are defined over $``$, we can apply the classical Weil’s technique of counting points over $`𝔽_q`$ (thus treating $`q`$ not as a formal variable but as a power of prime). After the counting is done, we will see that $`\mathrm{card}\overline{L}_n(𝔽_q)`$ is a polynomial in $`q`$ with positive integer coefficients, so that we can right away identify it with $`p_n`$:
$$p_n(q)=\mathrm{card}\overline{L}_n(𝔽_q)$$
$`(2.7)`$
The latter number can be calculated by directly applying (2.3) to the one–element partition $`\sigma `$, so that we get
$$\frac{\frac{p_n(q)}{n!}=\underset{l=1}{\overset{n}{}}}{(s_1,\mathrm{},s_l)s_1+\mathrm{}+s_l=ns_i1\frac{(q1)^{s_11}}{s_1!}\mathrm{}\frac{(q1)^{s_l1}}{s_l!}}$$
$$\underset{l=1}{\overset{n}{}}[\mathrm{coeff}.\mathrm{of}x^{nl}\mathrm{in}\left(\frac{e^x1}{x}\right)^l](q1)^{nl}.$$
Inserting this in the left hand side of (2.6) and summing over $`n`$ first, we obtain
$$\underset{n=1}{\overset{\mathrm{}}{}}\frac{p_n(q)}{n!}y^n=\underset{l=1}{\overset{\mathrm{}}{}}\underset{n=l}{\overset{\mathrm{}}{}}[\mathrm{coeff}.\mathrm{of}x^n\mathrm{in}(e^x1)^l](q1)^n$$
$$=\underset{l=1}{\overset{\mathrm{}}{}}\frac{1}{(q1)^l}(e^{(q1)y}1)^l$$
which gives (2.6).
2.3.1. Special cases. Here is a list of the Poincaré polynomials for small values of $`n`$:
$$p_1=1,p_2=q+1,p_3=q^2+4q+1,p_4=q^3+11q^2+11q+1,$$
$$p_5=q^4+26q^3+66q^2+26q+1,p_6=q^6+57q^5+302q^4+302q^2+57q+1.$$
The rank of $`H^2(L_n)`$ is $`2^nn1`$. Individual coefficients of of $`p_n(q)`$ are well known in combinatorics. They are called Euler numbers:
$$a_{n,i}=\mathrm{dim}H^{2i}(\overline{L}_n).$$
2.4. $`\overline{L}_n`$ and toric actions. Let $`\epsilon `$ be the trivial partition of $`B`$ of length one. The “big cell” $`L_\epsilon `$ of $`\overline{L}_B`$ (see 2.2 c)) has a canonical structure of the torsor (principal homogeneous space) over the torus $`T_B:=𝔾_m^B/𝔾_m`$ (where the subgroup $`𝔾_m`$ is embedded diagonally). In fact, $`^1\{x_0,x_{\mathrm{}}\}`$ is a $`𝔾_m`$–torsor, and the respective action of $`𝔾_m^B`$ on $`L_\epsilon `$, moving $`x_i,iB`$ via the $`i`$–th factor, produces an isomorphic marked curve exactly via the action of the diagonal.
Similarly, every stratum $`L_\sigma `$ is a torsor over $`T_\sigma :=_iT_{\sigma _i}`$ (see (2.4)), and there is a canonical surjective morphism $`T_BT_\sigma `$ so that $`L_B`$ is a union of $`T_B`$–orbits. In order to show that $`L_B`$ is a toric variety, it remains to show that these actions are compatible. This again can be done using the explicit construction of $`\overline{L}_n`$ and induction. For a change, we will provide a direct toric construction. We start with a more systematic treatment of the combinatorics involved.
2.4.1. Partitions of finite sets. For any finite set $`B`$, we call a partition $`\sigma `$ of $`B`$ a totally ordered set of non–empty subsets of $`B`$ whose union is $`B`$ and whose pairwise intersectons are empty. If a partition consists of $`N`$ subsets, it is called $`N`$–partition. If its components are denoted $`\sigma _1,\mathrm{},\sigma _N`$, or otherwise listed, this means that they are listed in their structure order. Another partition can be denoted $`\tau `$, $`\sigma ^{(1)}`$ etc. Notice that no particular ordering of $`B`$ is a part of the structure. This is why we replaced $`\{1,\mathrm{},n\}`$ here by an unstructured set $`B.`$
Let $`\sigma `$ be a partition of $`B`$, $`i,jB.`$ We say that $`\sigma `$ separates $`i`$ and $`j`$ if they belong to different components of $`\sigma `$. We then write $`i\sigma j`$ in order to indicate that the component containing $`i`$ comes earlier that the one containing $`j`$ in the structure order.
Let $`\tau `$ be an $`N+1`$–partition of $`B`$. If $`N1,`$ it determines a well ordered family of $`N`$ 2–partitions $`\sigma ^{(a)}`$:
$$\sigma _1^{(a)}:=\tau _1\mathrm{}\tau _a,\sigma _2^{(a)}:=\tau _{a+1}\mathrm{}\tau _N,a=1,\mathrm{},N.$$
$`(2.8)`$
In reverse direction, call a family of 2–partitions $`(\sigma ^{(i)})`$ good if for any $`ij`$ we have $`\sigma ^{(i)}\sigma ^{(j)}`$ and either $`\sigma _1^{(i)}\sigma _1^{(j)},`$ or $`\sigma _1^{(j)}\sigma _1^{(i)}.`$ Any good family is naturally well–ordered by the relation $`\sigma _1^{(i)}\sigma _1^{(j)}`$, and we will consider this ordering as a part of the structure. If a good family of 2–partitions consists of $`N`$ members, we will usually choose superscripts $`1,\mathrm{},N`$ to number these partitions in such a way that $`\sigma _1^{(i)}\sigma _1^{(j)}`$ for $`i<j.`$
Such a good family produces one $`(N+1)`$–partition $`\tau `$:
$$\tau _1:=\sigma _1^{(1)},\tau _2:=\sigma _1^{(2)}\sigma _1^{(1)},\mathrm{},\tau _N:=\sigma _1^{(N)}\sigma _1^{(N1)},\tau _{N+1}=\sigma _2^{(N)}.$$
$`(2.9)`$
This correspondence between good $`N`$–element families of 2–partitions and $`(N+1)`$–partitions is one–to–one, because clearly $`\sigma _1^{(i)}=\tau _1\mathrm{}\tau _i`$ for $`1iN.`$
Consider the case when $`\tau ^{(1)}=\sigma `$ is a 2–partition, and $`\tau ^{(2)}=\tau `$ is an $`N`$–partition, $`N2`$. Their union is good, iff there exists $`aN`$ and a 2–partition $`\alpha =(\tau _{a1},\tau _{a2})`$ of $`\tau _a`$ such that
$$\sigma =(\tau _1\mathrm{}\tau _{a1}\tau _{a1},\tau _{a2}\tau _{a+1}\mathrm{}\tau _N).$$
$`(2.10)`$
In this case we denote
$$\sigma \tau =\tau (\alpha ):=(\tau _1,\mathrm{},\tau _{a1},\tau _{a1},\tau _{a2},\tau _{a+1},\mathrm{},\tau _N).$$
$`(2.11)`$
###### \quad2.4.2. Lemma
Let $`\tau `$ be a partition of $`B`$ of length $`1,`$ and $`\sigma `$ a 2–partition. Then one of the three mutually exclusive cases occurs:
(i) $`\sigma `$ coincides with one of the partitions $`\sigma ^{(a)}`$ in (2.8). In this case we will say that $`\sigma `$ breaks $`\tau `$ between $`\tau _a`$ and $`\tau _{a+1}.`$
(ii) $`\sigma `$ coincides with one of the partitions (2.10). In this case we will say that $`\sigma `$ breaks $`\tau `$ at $`\tau _a`$.
(iii) None of the above. In this case we will say that $`\sigma `$ does not break $`\tau `$. This happens exactly when there is a neighboring pair $`(\tau _b,\tau _{b+1})`$ of elements of $`\tau `$ with the following property:
$$\tau _b\sigma _1\mathrm{},\tau _{b+1}\sigma _1\mathrm{}.$$
$`(2.12)`$
We will call $`(\tau _b,\tau _{b+1})`$ a bad pair (for $`\sigma `$).
Proof. Consider the sequence of sets
$$\sigma _1\tau _1,\sigma _1\tau _2,\mathrm{},\sigma _1\tau _N.$$
Produce from it a sequence of numbers 0, 1, 2 by the following rule: replace $`\sigma _1\tau _b`$ by 2, if it coincides with $`\tau _b`$, by 0 if it is empty, and by 1 otherwise. Cases (i) and (ii) above together will furnish all sequences of the form $`(2\mathrm{}20\mathrm{}0)`$, $`(2\mathrm{}210\mathrm{}0)`$, $`(10\mathrm{}0)`$. Each remaining admissible sequence will contain at least one pair of neighbors from the list 01, 02, 11, 12. For the respective pair of sets, (2.12) will hold.
2.5. Fan $`F_B`$. In this subsection we will describe a fan $`F_B`$ in the space $`N_B`$, where $`N_B:=\mathrm{Hom}(𝔾_m,T_B)`$, $`T_B:=𝔾_m^B/𝔾_m`$ as in the beginning of 2.4. Up to notation, we use \[Fu\] as the basic reference on fans and toric varieties.
Clearly, $`N_B`$ can be canonically identified with $`^B/`$, the latter subgroup being embedded diagonally. Similarly, $`N_B=^B/`$. We will write the vectors of this space (resp. lattice) as functions $`B`$ (resp. $`B`$) considered modulo constant functions. For a subset $`\beta B`$, let $`\chi _\beta `$ be the function equal 1 on $`\beta `$ and 0 elsewhere.
###### \quad2.5.1. Definition
The fan $`F_B`$ consists of the following $`l`$–dimensional cones $`C(\tau )`$ labeled by $`(l+1)`$–partitions $`\tau `$ of $`B`$.
If $`\tau `$ is the trivial 1–partition, $`C(\tau )=\{0\}`$.
If $`\sigma `$ is a 2–partition, $`C(\sigma )`$ is generated by $`\chi _{\sigma _1}`$, or, equivalently, $`\chi _{\sigma _2}`$, modulo constants.
Generally, let $`\tau `$ be an $`(l+1)`$–partition, and $`\sigma ^{(i)},i=1,\mathrm{},l`$, the respective good family of 2–partitions (2.9). Then $`C(\tau )`$ as a cone is generated by all $`C(\sigma ^{(i)})`$.
It is not quite obvious that $`F_B`$ is well defined. We sketch the relevant arguments.
First, all cones $`C(\tau )`$ are strongly convex. In fact, according to \[Fu\], p. 14, it suffices to check that $`C(\tau )(C(\tau ))=0`$. But $`C(\tau )`$ consists of classes of linear combinations with non–negative coefficients of functions
$$\chi _{\tau _1},\chi _{\tau _1}+\chi _{\tau _2},\mathrm{},\chi _{\tau _1}+\mathrm{}+\chi _{\tau _l}$$
if $`\tau `$ has length $`l+1`$. Non–vanishing function of this type cannot be constant.
Second, the same argument shows that $`C(\tau )`$ is actually $`l`$–dimensional.
Third, since the cone $`C(\tau )`$ is simplicial, one sees that $`(l1)`$–faces of $`C(\tau )`$ are exactly $`C(\tau ^{(i)})`$ where $`\tau ^{(i)}`$ is obtained from $`\tau `$ by uniting $`\tau _i`$ with $`\tau _{i+1}`$, which is equivalent to omitting $`C(\sigma ^{(i)})`$ from the list of generators. More generally, $`C(\tau ^{})`$ is a face of $`C(\tau )`$ iff $`\tau \tau ^{}`$ as in (2.3), that is, if $`\tau `$ is a refinement of $`\tau ^{}`$.
Fourth, let $`C(\tau ^{(i)}),i=1,2,`$ be two cones. We have to check that their intersection is a cone of the same type. An obvious candidate is $`C(\tau )`$ where $`\tau `$ is the crudest common refinement of $`\tau ^{(1)}`$ and $`\tau ^{(2)}`$. This is the correct answer.
In order to see this, let us a give a different description of $`F_B`$ which will simultaneously show that the support of $`F_B`$ is the whole space. Let $`\chi :𝔹`$ represent an element $`\overline{\chi }N_B.`$ It defines a unique partition $`\tau `$ of $`B`$ consisting of the level sets of $`\chi `$ ordered in such a way that the values of $`\chi `$ decrease. Clearly, $`\tau `$ depends only on $`\overline{\chi }`$, and $`\chi `$ modulo constants can be expressed as a linear combination of $`\chi _{\tau _1}+\mathrm{}+\chi _{\tau _i}`$, $`1il`$ with positive coefficients. In other words, $`\chi `$ belongs to the interior part of $`C(\tau )`$. On the boundary, some of the strict inequalities between the consecurive values of $`\chi `$ become equalities. This proves the last assertion.
We see now that $`F_B`$ satisfies the definition of \[Fu\], p. 20, and so is a fan.
2.6. Toric varieties $`\overline{}_B`$. We now define $`\overline{}_B`$ (later to be identified with $`\overline{L}_B`$) as the toric variety associated with the fan $`F_B`$.
To check that it is smooth, it suffices to show that each $`C(\tau )`$ is generated by a part of a basis of $`N_B`$ (see \[Fu\], p. 29). In fact, let us choose a total ordering of $`B`$ such that if $`i\tau _k,j\tau _l`$ and $`k<l`$, then $`i<j`$. Let $`B_kB`$ consist of the first $`k`$ elements of $`B`$ in this ordering. Then the classes of the characteristic functions of $`B_1,B_2,\mathrm{},B_{n1}`$, $`n=\mathrm{card}B`$, form a basis of $`N_B`$, and $`\{\chi _{\sigma ^{(i)}}\}`$ is a part of it.
To check that $`\overline{}_B`$ is proper, we have to show that the support of $`F_B`$ is the total space. We have already proved this.
As any toric variety, $`\overline{}_B`$ carries a family of subvarieties which are the closures of the orbits of $`T_B`$ and which are in a natural bijection with the cones $`C(\tau )`$ in $`F_B`$. We denote them $`\overline{}_\tau `$. They are smooth. The respective orbit which is an open subset of $`\overline{}_\tau `$ is denoted $`_\tau `$.
2.6.1. Forgetful morphisms and a family of pointed curves over $`\overline{}_B`$. Assume that $`BB^{}`$. Then we have the projection morphism $`^B^{}^B`$ which induces the morphism $`f^{B^{},B}:N_B^{}N_B.`$ It satisfies the property stated in the last lines of \[Fu\], p. 22: for each cone $`C(\tau ^{})F_B^{}`$, there exists a cone $`C(\tau )F_B`$ such that $`f^{B^{},B}(C(\tau ^{}))C(\tau ).`$ In fact, $`\tau `$ is obtained from $`\tau ^{}`$ by deleting elements of $`B^{}B`$ and then deleting the empty subsets of the resulting partition of $`B`$.
Therefore, we have a morphism $`f_{}^{B^{},B}:\overline{}_B^{}\overline{}_B`$ (\[Fu\], p. 23) which we will call forgetful one (it forgets elements of $`B^{}B`$).
###### \quad2.6.2. Proposition
If $`B^{}B`$ consists of one element, then the forgetful morphism $`\overline{}_B^{}\overline{}_B`$ has a natural structure of a stable $`(\{0,\mathrm{}\},B)`$–pointed curve of genus zero.
Proof. Let us first study the fibers of the forgetful morphism. Let $`\tau `$ be a partition of $`B`$ of length $`l+1`$ and $`_\tau `$ the respective orbit in $`\overline{}_B.`$ Its inverse image in $`\overline{}_B^{}`$ is contained in the union $`\overline{}_\tau ^{}`$ where $`\tau ^{}`$ runs over partitions of $`B^{}`$ obtained by adding the forgotten point either to one of the parts $`\tau _i`$, or inserting it in between $`\tau _i`$ and $`\tau _{i+1}`$, or else putting it at the very beginning or at the very end as a separate part.
The inverse image of any point $`x_\tau `$ is acted upon by the multiplicative group $`𝔾_m=\mathrm{Ker}(T_B^{}T_B)`$. This action breaks the fiber into a finite number of orbits which coincide with the intersections of this fiber with various $`_\tau ^{}`$ described above. When $`\tau ^{}`$ is obtained by adding the forgotten point to one of the parts, this intersection is a torsor over the kernel, otherwise it is a point. As a result, we get that the fiber is a chain of $`^1`$’s, whose components are labeled by the components of $`\tau `$ and singular points by the neighboring pairs of components.
The forgetful morphism is flat, because locally in toric coordinates it is described as adjoining a variable and localization.
In order to describe the two white sections of the forgetful morphism, consider two partitions $`(B^{}B,B)`$ and $`(B,B^{}B)`$ of $`B^{}`$ and the respective closed strata. It is easily seen that the forgetful morphism restricted to these strata identifies them with $`\overline{}_B`$. We will call them $`x_0`$ and $`x_{\mathrm{}}`$ respectively.
Finally, to define the $`j`$–th black section, $`jB`$, consider the morphism of lattices $`s_j:N_BN_B^{}`$ which extends a function $`\chi `$ on $`B`$ to the function $`s_j(\chi )`$ on $`B^{}`$ taking the value $`\chi (j)`$ at the forgotten point. This morphism satisfies the condition of \[Fu\], p. 22: each cone $`C(\tau )`$ from $`F_B`$ lands in an appropriate cone $`C(\tau ^{})`$ from $`F_B^{}`$. This must be quite clear from the description at the end of 2.5.1: $`\tau ^{}`$ is obtained from $`\tau `$ by adding the forgotten point to the same part to which $`j`$ belongs. Hence we have the induced morphisms $`s_j:\overline{}_B\overline{}_B^{}`$ which obviously are sections. Moreover, they do not intersect $`x_0`$ and $`x_{\mathrm{}}`$, and they are distributed among the components of the reducible fibers exactly as expected.
###### \quad2.6.3. Theorem
The morphism $`\overline{}_B\overline{L}_B`$ inducing the family described in the Proposition 2.6.2 is an isomorphism.
This can be proved by induction on $`\mathrm{card}B`$ with the help of the more detailed analysis of the forgetful morphism, as above. We omit the details because they are not instructive.
An important corollary of this Theorem is the existence of a surjective birational morphism $`\overline{M}_{0,n+2}\overline{L}_n`$ corresponding to any choice of two different labels $`i,j`$ in $`(1,\mathrm{},n+2)`$. In terms of the of the respective functors, this morphism blows down all the components of a stable $`(n+2)`$–labeled curve except for those that belong to the single path from the component containing the $`i`$–th point to the one containing the $`j`$–th point.
In fact, M. Kapranov has shown the existence of such a morphism for $`\overline{}_n`$ in place of $`\overline{L}_n`$ (see \[Ka2\], p. 102). He used a different description of $`\overline{}_n`$ in terms of the defining polyhedron, which he identified with the so called permutohedron, the convex hull of the $`𝕊_n`$–orbit of $`(1,2,\mathrm{},n)`$. He has also proved that $`\overline{}_n`$ can be identified with the closure of the generic orbit of the torus in the space of complete flags in an $`n`$–dimensional vector space.
2.7. Combinatorial model of $`H^{}(\overline{}_B)`$. We will denote by $`[\overline{}_\sigma ]_{}`$ (resp. $`[\overline{}_\sigma ]^{}`$) the homology (resp. the dual cohomology) class of $`\overline{}_\sigma `$.
The remaining parts of this section (and the Appendix) are dedicated to the study of linear and non–linear relations between these classes, in the spirit of \[KM\] and \[KMK\], but with the help of the standard toric techniques.
Consider a family of pairwise commuting independent variables $`l_\sigma `$ numbered by 2–partitions of $`B`$ and introduce the ring
$$H_B^{}:=_B/I_B$$
$`(2.13)`$
where $`_B`$ is freely generated by $`l_\sigma `$ (over an arbitrary coefficient ring $`k`$), and the ideal $`I_B`$ is generated by the following elements indexed by pairs $`i,jB`$:
$$r_{ij}^{(1)}:=\underset{\sigma :i\sigma j}{}l_\sigma \underset{\tau :j\tau i}{}l_\tau ,$$
$`(2.14)`$
$$r^{(2)}(\sigma ,\tau ):=l_\sigma l_\tau \mathrm{if}i\sigma j\mathrm{and}j\tau i\mathrm{for}\mathrm{some}i,j.$$
$`(2.15)`$
###### \quad2.7.1. Theorem
a) There is a well defined ring isomorphism $`_B/I_BA^{}(\overline{}_B,k)`$ such that $`l_\sigma \mathrm{mod}I_B[\overline{}_\sigma ]^{}`$. The Chow ring $`A^{}(\overline{}_B,k)`$ and the cohomology ring $`H^{}(\overline{}_B,k)`$ are canonically isomorphic.
b) The boundary divisors (strata corresponding to 2–partitions) intersect transversally.
Proof. We must check that the ideal of relations between $`2^n2`$ dual classes of the boundary divisors $`[\overline{}_\sigma ]^{}`$ contains and is generated by the following relations:
$$R_{ij}^{(1)}:\underset{\sigma :i\sigma j}{}[\overline{}_\sigma ]^{}\underset{\tau :j\tau i}{}[\overline{}_\tau ]^{}=0.$$
$`(2.16)`$
If $`i\sigma j`$ and $`j\tau i`$, then
$$R^{(2)}(\sigma ,\tau ):[\overline{}_\sigma ]^{}[\overline{}_\tau ]^{}=0.$$
$`(2.17)`$
We refer to the Proposition on p. 106 of \[Fu\] which gives a system of generators for this ideal for any smooth proper toric variety (Fulton additionally assumes projectivity which we did not check, but see \[Da\], Theorem 10.8 for the general proper case).
In our notation, these generators look as follows.
To get the complete system of linear relations, we must choose some elements $`m`$ in the dual lattice of $`N_B`$ spanning this lattice and form the sums $`_\sigma m(\chi _{\sigma _1})[\overline{}_\sigma ]^{}`$, where $`\sigma `$ runs over all 2–partitions. In our case, the dual lattice is spanned by the linear functionals $`m_{ij}:\chi \chi (i)\chi (j)`$ for all pairs $`i,jB.`$ Writing the respective relation, we get (2.16).
The complete system of nonlinear relations is given by the monomials $`l_{\sigma ^{(1)}}\mathrm{}l_{\sigma ^{(k)}}`$ such that $`(C(\sigma ^{(1)}),\mathrm{},C(\sigma ^{(k)}))`$ do not span a cone in $`F_B`$. This means that some pair $`(C(\sigma ^{(a)}),C(\sigma ^{(b)}))`$ already does not span a cone, because otherwise the respective 2–partitions would form a good family (cf. 2.4.1). And in view of Lemma 2.4.2 (iii), we can find $`i,jB`$ such that $`i\sigma ^{(a)}j`$ and $`j\sigma ^{(b)}i`$. Hence (2.16) and (2.17) together constitute a generating system of relations.
The remaining statements are true for all smooth complete toric varieties defined by simplicial fans.
2.8. Combinatorial structure of the cohomology ring. In the remaining part of this section we fix a finite set $`B`$ and study $`H_B^{}`$ as an abstract ring.
For an $`(N+1)`$–partition $`\tau `$ define the respective good monomial $`m(\tau )`$ by the formula
$$m(\tau )=l_{\sigma ^{(1)}}\mathrm{}l_{\sigma ^{(N)}}_B.$$
If $`\tau `$ is the trivial 1–partition, we put $`m(\tau ):=1.`$ In view of the Theorem 2.7.1, $`m(\tau )`$ represents the cohomology class of $`\overline{}_\tau `$.
Notice that if we have two good families of 2–partitions whose union is also good, then the product of the respective good monomials is a good monomial. This defines a partial operation $``$ on pairs of partitions
$$m(\tau ^{(1)})m(\tau ^{(2)})=m(\tau ^{(1)}\tau ^{(2)}).$$
###### \quad2.8.1. Proposition
Good monomials and $`I_B`$ span $`_B.`$ Therefore, images of good monomials span $`H_B^{}.`$
Proof. We make induction on the degree. In degrees zero and one the statement is clear because $`l_\sigma `$ are good. If it is proved in degree $`N`$, it suffices to check that for any 2–partition $`\sigma `$ and any nontrivial partition $`\tau `$, $`l_\sigma m(\tau )`$ is a linear combination of good monomials modulo $`I_B.`$ We will consider the three cases of Lemma 2.4.2 in turn.
(i) $`\sigma `$ breaks $`\tau `$ between $`\tau _a`$ and $`\tau _{a+1}`$.
This means that $`l_\sigma `$ divides $`m(\tau )`$.
Choose $`i\tau _a,j\tau _{a+1}.`$ In view of (2.14), we have
$$\left(\underset{\rho :i\rho j}{}l_\rho \underset{\rho :j\rho i}{}l_\rho \right)m(\tau )I_B.$$
$`(2.18)`$
But if $`j\rho i`$, then $`l_\rho m(\tau )I_B`$ because of (2.15). Among the terms with $`i\rho j`$ there is one $`ł_\sigma .`$ For all other $`\rho `$’s, $`l_\rho `$ cannot divide $`m(\tau )`$ since other divisors put $`i`$ and $`j`$ in the same part of the respective partition. Therefore, $`l_\rho m(\tau )`$ either belongs to $`I_B`$, or is good. So finally (2.18) allows us to express $`l_\sigma m(\tau )`$ as a sum of good monomials and an element of $`I_B:`$
$$l_\sigma m(\tau )=\underset{\rho \sigma ,i\rho j}{}m(\rho \tau )\mathrm{mod}I_B$$
where the terms for which $`\rho \tau `$ is not defined must be interpreted as zero. More precisely, there are two types of non–vanishing terms. One corresponds to all 2–partitions $`\alpha `$ of $`\tau _a`$ such that $`i\tau _{a1}`$ which we will write as $`i\alpha `$. Another corresponds to 2–partitions $`\beta `$ of $`\tau _{a+1}`$ with $`j`$ belonging to the second part, $`\beta j`$:
$$l_\sigma m(\tau )=\underset{\alpha :i\alpha }{}m(\tau (\alpha ))\underset{\beta :\beta j}{}m(\tau (\beta ))\mathrm{mod}I_B.$$
$`(2.19)`$
Notice that there are several ways to write the right hand side, depending on the choice of $`i,j.`$ Hence good monomials are not linearly independent modulo $`I_B.`$
(ii) $`\sigma `$ breaks $`\tau `$ at $`\tau _a`$.
According to the analysis above, this means that
$$l_\sigma m(\tau )=m(\sigma \tau )=m(\tau (\alpha ))$$
$`(2.20)`$
for an appropriate partition $`\alpha `$ of $`\tau _a`$.
(iii) $`\sigma `$ does not break $`\tau `$.
In this case, let $`(\tau _b,\tau _{b+1})`$ be a bad pair for $`\sigma `$. Then from (2.12) it follows that there exist $`i,jB`$ such that $`i\sigma j`$ and $`j\sigma ^{(a)}i`$. Hence $`l_\sigma m(\tau )`$ is divisible by $`r^{(2)}(\sigma ,\sigma ^{(a)})`$ and
$$l_\sigma m(\tau )=0\mathrm{mod}I_B.$$
2.8.2. Linear combinations of good monomials belonging to $`I_B`$. Let $`\tau =(\tau _1,\mathrm{},\tau _N)`$ be a partition of $`B`$. Choose $`aN`$ such that $`|\tau _a|2`$, and two elements $`i,j\tau _a,ij.`$ For any ordered 2–partition $`\alpha =(\tau _{a1},\tau _{a2})`$ of $`\tau _a`$, denote by $`\tau (\alpha )`$ the induced $`N+1`$–partition of $`B`$ as above:
$$(\tau _1,\mathrm{},\tau _{a1},\tau _{a1},\tau _{a2},\tau _{a+1},\mathrm{},\tau _N).$$
Finally, put
$$r_{ij}^{(1)}(\tau ,a):=\underset{\alpha :i\alpha j}{}m(\tau (\alpha ))\underset{\alpha :j\alpha i}{}m(\tau (\alpha )).$$
$`(2.21)`$
Choosing for $`\tau `$ the trivial 1–partition, we get (2.14) so that these elements span the intersection of $`I_B`$ with the space of good monomials of degree one.
Generally, all $`r_{ij}^{(1)}(\tau ,a)`$ belong to $`I_B.`$ In fact, keeping the notations above, consider
$$r_{ij}^{(1)}m(\tau )=\left(\underset{\rho :i\rho j}{}l_\rho \underset{\rho :j\rho i}{}l_\rho \right)m(\tau )I_B.$$
$`(2.22)`$
Arguing as above, we see that the summand corresponding to $`\rho `$ in (2.18) either belongs to $`I_B`$, or is a good monomial, and the latter happens exactly for those partitions $`\rho `$ that are of the type $`\tau (\alpha )`$ with either $`i\alpha j`$, or $`j\alpha i.`$ Hence (2.21) lies in $`I_B.`$ This proves our claim.
###### \quad2.9. Theorem
Elements (2.21) span the intersection of $`I_B`$ with the space generated by good monomials.
Proof. Define the linear space $`H_B`$ generated by the symbols $`\mu (\tau )`$ for all partitions of $`B`$ as above which satisfy analogs of the linear relations (2.21): for all $`(\tau ,\tau _a,i,j)`$ as above we have
$$\underset{\alpha :i\alpha j}{}\mu (\tau (\alpha ))\underset{\alpha :j\alpha i}{}\mu (\tau (\alpha ))=0.$$
$`(2.23)`$
###### \quad2.9.1. Technical Lemma
There exists an (obviously unique) structure of the $`H_B^{}`$–module on $`H_B`$ with the following multiplication table.
(i) If $`\sigma `$ breaks $`\tau `$ between $`\tau _a`$ and $`\tau _{a+1}`$, then for any choice of $`i\tau _a,j\tau _{a+1}`$
$$l_\sigma \mu (\tau )=\underset{\alpha :i\alpha }{}\mu (\tau (\alpha ))\underset{\beta :\beta j}{}\mu (\tau (\beta )).$$
$`(2.24)`$
(cf. (2.19)).
(ii) If $`\sigma `$ breaks $`\tau `$ at $`\tau _a`$, then
$$l_\sigma \mu (\tau )=\mu (\sigma \tau ).$$
$`(2.25)`$
(cf. (2.20)).
(iii) If $`\sigma `$ does not break $`\tau `$, then
$$l_\sigma \mu (\tau )=0.$$
$`(2.26)`$
Our proof of the Technical Lemma consists in the direct verification that the prescriptions (2.24)–(2.23) are compatible with all relations that we have postulated. Unfortunately, such strategy requires the painstaking case–by–case treatment of a long list of combinatorially distinct situations, and we relegate it to the Appendix.
2.9.2. Deduction of Theorem 2.9 from the Technical Lemma. Since elements (2.21) belong to $`I_B,`$ there exists a surjective linear map $`s:H_BH_B^{}`$, $`\mu (\tau )m(\tau ).`$ Now denote by $`\mathrm{𝟙}`$ the element $`\mu (\epsilon )`$ where $`\epsilon `$ is the 1–partition. Then $`t:m(\sigma )m(\sigma )\mathrm{𝟙}`$ is a linear map $`H_B^{}H_B`$. From (2.25) one easily deduces that $`m(\tau )\mathrm{𝟙}=\mu (\tau )`$ so that $`s`$ and $`t`$ are mutually inverse. Therefore, (2.22) span the linear relations between the images of good monomials in $`H_B^{}.`$
According to the Theorem 2.4.1, $`H_B`$, together with its structure of $`H_B^{}`$–module, is a combinatorial model of the homology module $`H_{}(\overline{}_B,k)`$. The generators $`\mu (\tau )`$ correspond to $`[\overline{}_\tau ]_{}`$.
§3. Pencils of flat connections
and the Commutativity Equations
3.1. Notation. Let $`M`$ be a (super)manifold over a field $`k`$ of characteristic zero in one of the standard categories (smooth, complex analytic, schemes, formal $`\mathrm{}`$). We use the conventions spelled out in \[Ma\], I.1. In particular, differentials in the de Rham complex $`(\mathrm{\Omega }_M^{},d)`$ and connections are odd. This determines our sign rules; parity of an object $`x`$ is denoted $`\stackrel{~}{x}.`$
Let $``$ be a locally free sheaf (of sections of a vector bundle) on $`M,`$ $`_0`$ a connection on $``$, that is an odd $`k`$–linear operator $`\mathrm{\Omega }_M^1`$ satisfying the Leibniz identity
$$_0(\phi f)=d\phi f+(1)^{\stackrel{~}{\phi }}\phi _0f,\phi 𝒪_M,f.$$
$`(3.1)`$
This operator extends to a unique operator on the $`\mathrm{\Omega }_M^{}`$–module $`\mathrm{\Omega }_M^{}`$ denoted again $`_0`$ and satisfying the same identity (3.1) for any $`\phi \mathrm{\Omega }_M`$. Any other connection differential $``$ restricted to $``$ has the form $`_0+𝒜`$ where $`𝒜:\mathrm{\Omega }_M^1`$ is an odd $`𝒪_M`$–linear operator: $`𝒜(\phi f)=(1)^{\stackrel{~}{\phi }}\phi 𝒜(f).`$ Any connection naturally extends to the whole tensor algebra generated by $`,`$ in particular, to $`nd.`$
The connection $`_0`$ is called flat, iff $`_0^2=0.`$ A pencil of flat connections is a line in the space of connections $`_\lambda :=_0+\lambda 𝒜`$ such that $`_\lambda ^2=0`$ ($`\lambda `$ is an even parameter). In the smooth, analytic or formal category, $`_0`$ is flat iff $``$ locally admits a basis of flat sections $`f,_0f=0`$.
###### \quad3.2. Proposition
$`_0+\lambda 𝒜`$ is a pencil of flat connections iff the following two conditions are satisfied:
(i) Everywhere locally on $`M`$, we have
$$𝒜=_0$$
$`(3.2)`$
for some $`nd`$.
(ii) Such an operator $``$ satisfies the quadratic differential equation
$$_0_0=0.$$
$`(3.3)`$
Proof. Calculating the coefficient of $`\lambda `$ in $`_\lambda ^2=0`$ we get $`_0𝒜=0`$. But the complex $`\mathrm{\Omega }_M^{}`$ is the resolution of the sheaf of flat sections $`\mathrm{Ker}_0.`$ This furnishes (i); (ii) means the vanishing of the coefficient at $`\lambda ^2.`$
3.2.1. Remarks. a) Write $``$ as a matrix in a basis of $`_0`$–flat sections of $``$, whose entries are local functions on $`M.`$ Then (3.3) becomes
$$dd=0.$$
$`(3.4)`$
These equations written in local coordinates $`(t^i)`$ on $`M`$ were called “$`t`$–part of the $`t`$$`t^{}`$ equations” by S. Cecotti and C. Vafa. A. Losev in \[Lo1\] suggested to call them “the Commutativity Equations”.
b) If $`_0\phi _0=0,`$ then
$$(_0+\lambda _0)(e^\lambda \phi _0)=0.$$
3.2.2. Pencils of flat connections related to Frobenius manifolds. Any solution to the Associativity Equations produces a pencil of flat connections.
To explain this we will use the geometric language due to B. Dubrovin (and the notation of \[Ma\], I.1.5). Consider a Frobenius manifold $`(M,g,)`$ where
$$:𝒯_M_{𝒪_M}𝒯_M𝒯_M$$
is a (super)commutative associative multiplication on the tangent sheaf satisfying the potentiality condition, and $`g`$ is an invariant flat metric (no positivity condition is assumed, only symmetry and non–degeneracy). Denote by $`_0`$ the Levi–Civita connection of $`g`$. Finally, denote by $`𝒜`$ the operator obtained from the Frobenius multiplication in $`𝒯_M`$ (\[Ma\], I.1.4). In other words, consider the pencil of connections on $`=𝒯_M`$ whose covariant derivatives are
$$(_0+\lambda 𝒜)_X(Y):=_{0,X}(Y)+\lambda XY.$$
This pencil is flat (see \[Ma\], Theorem I.1.5, p. 20). In fact, $``$ written in a basis of $`_0`$–flat coordinates and the respective flat vector fields is simply the matrix of the second derivatives of a local potential $`\mathrm{\Phi }`$ (with one subscript raised). This is the first structure connection of $`M`$.
This pencil admits an infinite dimensional deformation: one should take the canonical extension of the potential to the large phase space and consider the coordinates with gravitational descendants as parameters of the deformation.
Another family of flat connections, this time on the cotangent sheaf of a Frobenius manifold $`M`$ admitting an Euler vector field $`E`$ (see \[Ma\], pp. 23–24), is defined as follows. Denote the scalar product on vector fields $`\stackrel{ˇ}{g}_\lambda (X,Y):=g((E\lambda )^1X,Y).`$ The inverse form induces a pencil of flat metrics on the cotangent sheaf, whose Levi–Civita connections however do not form a pencil of flat connections in our sense (see \[Du1\], Appendix D, and \[Du3\] for a general discussion of such setup). This is the second structure connection of $`M`$.
3.2.3. Flat coordinates and gravitational descendants. One can show that 1–forms on $`M`$ flat with respect to the dualized first structure connection are closed and therefore locally exact. Their integrals are called deformed flat coordinates. In \[Du2\], Example 2.3 and Theorem 2.2, B. Dubrovin gives explicit formal series in $`\lambda `$ ($`z`$ in his notation) for suitably normalized deformed flat coordinates. Coefficients of these series involve some correlators with gravitational descendants, namely those for which the non–trivial operators $`\tau _p`$ are applied only at one point. In \[KM2\] and \[Ma\], VI.7.2, p.278, it was shown that two–point correlators of this kind determine a linear operator in the large phase space which transforms the modified correlators with descendants into non–modified ones (in any genus). This is important because apriori only modified correlators are defined for an arbitrary Cohomological Field Theory in the sense of \[KM1\], which is not necessarily quantum cohomology of a manifold.
3.2.4. Pencils of flat connections in a global setting. Pencils of flat connections appear also in the context of Simpson’s non–abelian Hodge theory. Briefly, consider a smooth projective manifold $`M`$ over $``$. One can define two moduli spaces, $`Mod_1`$ and $`Mod_2`$. The first one classifies flat connections (on variable vector bundles $``$ with vanishing rational Chern classes) with semisimple Zariski closure of the monodromy group. The second one classifies semistable Higgs pairs $`(,𝒜)`$ where $`𝒜`$ is an operator as in 3.1, satisfying only the condition $`𝒜𝒜=0.`$ (In fact, one should only consider smooth points of the respective moduli spaces). N. Hitchin, C. Simpson, Fujiki et al. established that $`Mod_1`$ and $`Mod_2`$ are canonically isomorphic as $`C^{\mathrm{}}`$–manifolds, but their complex structures $`I`$, $`J`$ are different, and together with $`K=IJ`$ produce a hypercomplex manifold.
P. Deligne has shown that the respective twistor space is precisely the moduli space of the pencils of flat connections on $`M`$ (where the Higgs complex structure corresponds to the point $`\lambda =\mathrm{}`$ in our notation).
For details, see \[Si\].
3.3. Formal solutions to the Commutativity Equations and the homology of $`\overline{L}_n.`$ In \[KM1\] and \[KMK\] it was shown that formal solutions to the Associativity Equations are cyclic algebras over the cyclic genus zero homology modular operad $`(H_{}(\overline{M}_{0,n+1}))`$ (see also \[Ma\], III.4). The main goal of this section is to show the similar role of the homology of the spaces $`\overline{L}_n`$ in the theory of the Commutativity Equations. This was discovered and discussed on a physical level in \[Lo1\], \[Lo2\]. Here we supply precise mathematical statements with proofs.
Unlike the case of the Associativity Equations, we will have to deal here with modules over an algebra (depending explicitly on the base space) rather than with algebras over an operad. The main ingredient of the construction is the direct sum of the homology spaces of all $`\overline{L}_n`$ endowed with the multiplication coming from the boundary morphisms. We work with the combinatorial models of these spaces defined in 2.9.1.
We start with some preparations. Let $`V=_{n=1}^{\mathrm{}}V_n`$ be a graded associative $`k`$–algebra (without identity) in the category of vector $`k`$–superspaces over a field $`k`$. We will call it an $`𝕊`$–algebra, if for each $`n`$, an action of the symmetric group $`𝕊_n`$ on $`V_n`$ is given such that the multiplication map $`V_mV_nV_{m+n}`$ is compatible with the action of $`𝕊_m\times 𝕊_n`$ embedded in an obvious way into $`𝕊_{m+n}.`$
If $`V`$ is an $`𝕊`$–algebra, then the sum of subspaces $`J_n`$ spanned by $`(1s)v,s𝕊_n,vV_n,`$ is a double–sided ideal in $`V.`$ Hence the sum of the coinvariant spaces $`V_{𝕊_n}=V_n/J_n`$ is a graded ring which we denote $`V_𝕊`$.
If $`V`$, $`W`$ are two $`𝕊`$–algebras, then the diagonal part of their tensor product $`_{n=1}^{\mathrm{}}V_nW_n`$ is an $`𝕊`$–algebra as well.
Let $`T`$ be a vector superspace (below always assumed finite–dimensional). Its tensor algebra (without the rank zero part) is an $`𝕊`$–algebra.
As a less trivial example, consider $`H_{}:=_{n=1}^{\mathrm{}}H_n`$ where we write $`H_n`$ for $`H_{\{1,\mathrm{},n\}}`$. The multiplication law is given by what becomes the boundary morphisms in the geometric setting: if $`\tau ^{(1)}`$ (resp. $`\tau ^{(2)}`$) is a partition of $`\{1,\mathrm{},m\}`$ (resp. of $`\{1,\mathrm{},n\}`$), then
$$\mu (\tau ^{(1)})\mu (\tau ^{(2)})=\mu (\tau ^{(1)}\tau ^{(2)})$$
$`(3.5)`$
where the concatenated partition of $`\{1,\mathrm{},m,m+1,\mathrm{},m+n\}`$ is defined in an obvious way, shifting all the components of $`\tau ^{(2)}`$ by $`m`$.
Our main protagonist is the algebra of coinvariants of the diagonal tensor product of these examples:
$$H_{}T:=\left(_{n=1}^{\mathrm{}}H_nT^n\right)_𝕊.$$
$`(3.6)`$
We now fix $`T`$ and another vector superspace $`F`$ and assume that the ground field $`k`$ has characteristic zero.
###### \quad3.3.1. Theorem
There is a natural bijection between the set of representations of $`H_{}T`$ in $`F`$ and the set of pencils of flat connections on the trivial bundle with fiber $`F`$ on the formal completion of $`T`$ at the origin.
This bijection will be precisely defined and discussed below: see Proposition 3.6.1. Before passing to this definition and the proof of the Theorem, we will give a down–to–earth coordinate–dependent description of the representations of $`H_{}T`$.
3.4. Matrix correlators. Fix $`T`$ and choose its parity homogeneous basis $`(\mathrm{\Delta }_a|aI)`$ where $`I`$ is a finite set of indices.
For any $`n1`$, the space $`H_nT^n`$ is spanned by the elements
$$\mu (\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}$$
$`(3.7)`$
where $`\tau ^{(n)}`$ runs over all partitions of $`\{1,\mathrm{},n\}`$ whereas $`(a_1,\mathrm{},a_n)`$ runs over all maps $`\{1,\mathrm{},n\}I:ia_i.`$
In view of the Theorem 2.9, all linear relations between these elements are spanned by the following ones: choose $`(a_1,\mathrm{},a_n)`$ and $`(\tau ^{(n)},\tau _r^{(n)},ij\tau _r^{(n)})`$, then
$$\underset{\alpha :i\alpha j}{}\mu (\tau ^{(n)}(\alpha ))\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}\underset{\alpha :j\alpha i}{}\mu (\tau ^{(n)}(\alpha ))\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}=0$$
$`(3.8)`$
where the summation is taken over all 2–partitions $`\alpha `$ of $`\tau _r^{(n)}`$ separating $`i`$ and $`j`$.
The action of a permutation $`is(i)`$ on (3.7) is defined by
$$s\left(\mu (\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}\right)=\epsilon (s,(a_i))\mu (s(\tau ^{(n)}))\mathrm{\Delta }_{a_{s(1)}}\mathrm{}\mathrm{\Delta }_{a_{s(n)}}.$$
$`(3.9)`$
Here $`\epsilon (s,(a_i))=\pm 1`$ is the sign of the permutation induced by $`s`$ on the subfamily of odd $`\mathrm{\Delta }_{a_i}`$’s, and $`s(\tau ^{(n)})`$ is defined as follows:
$$s(i)s(\tau ^{(n)})_r\mathrm{iff}i\tau _r^{(n)}.$$
$`(3.10)`$
Finally, the multiplication rule between the generators in the diagonal tensor product is given by:
$$\mu (\tau ^{(m)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_m}\mu (\tau ^{(n)})\mathrm{\Delta }_{b_1}\mathrm{}\mathrm{\Delta }_{b_n}$$
$$=\mu (\tau ^{(m)}\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_m}\mathrm{\Delta }_{b_1}\mathrm{}\mathrm{\Delta }_{b_n}.$$
$`(3.11)`$
Any linear representation $`K:H_{}T\mathrm{End}F`$ can be described as a linear representation of the diagonal tensor product satisfying additional symmetry restrictions. To spell it out explicitly, we define the matrix correlators of $`K`$ as the following family of endomorphisms of $`F`$:
$$\tau ^{(n)}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}:=K(\mu (\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}).$$
$`(3.12)`$
###### \quad3.4.1. Claim
Matrix correlators of any representation satisfy the following relations:
(i) $`𝕊_n`$–symmetry:
$$s^1(\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}=\epsilon (s,(a_i))\tau ^{(n)}\mathrm{\Delta }_{a_{s(1)}}\mathrm{}\mathrm{\Delta }_{a_{s(n)}}.$$
$`(3.13)`$
(ii) Factorization:
$$(\tau ^{(m)}\tau ^{(n)})\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_m}\mathrm{\Delta }_{b_1}\mathrm{}\mathrm{\Delta }_{b_n}=\tau ^{(m)}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_m}\tau ^{(n)}\mathrm{\Delta }_{b_1}\mathrm{}\mathrm{\Delta }_{b_n}.$$
$`(3.14)`$
(iii) Linear relations:
$$\underset{\alpha :i\alpha j}{}\tau ^{(n)}(\alpha )\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}\underset{\alpha :j\alpha i}{}\tau ^{(n)}(\alpha )\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}=0$$
$`(3.15)`$
Conversely, any family of elements of $`\mathrm{End}F`$ defined for all $`n,(a_1,\mathrm{},a_n),\tau ^{(n)}`$ and satisfying (3.13)–(3.15) consists of matrix correlators of a well defined representation $`K:H_{}T\mathrm{End}F`$.
In fact, we obtain (3.13) by applying $`K`$ to (3.9) written for $`s^1(\tau ^{(n)})`$ in place of $`\tau ^{(n)}`$, because $`K`$, coming from $`H_{}T`$, vanishes on the image of $`1s`$. Moreover, (3.14) means the compatibility with the multiplication of the generators. Finally, (3.15) is a necessary and sufficient condition for the extendability of the system of matrix correlators to a linear map $`K`$.
Notice that we can replace here $`\mathrm{End}F`$ by an arbitrary associative superalgebra over $`k`$.
3.5. Top matrix correlators. Define top matrix correlators of $`K`$ as the subfamily of correlators corresponding to the identical partitions $`\epsilon ^{(n)}`$ of $`\{1,\mathrm{},n\}`$:
$$\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}:=\epsilon ^{(n)}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}.$$
###### \quad3.5.1. Proposition
Top matrix correlators satisfy the following relations:
$$\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}=\epsilon (s,(a_i))\mathrm{\Delta }_{a_{s(1)}}\mathrm{}\mathrm{\Delta }_{a_{s(n)}}$$
$`(3.16)`$
and
$$\underset{\sigma :i\sigma j}{}\epsilon (\sigma ,(a_k))\underset{k\sigma _1}{}\mathrm{\Delta }_{a_k}\underset{k\sigma _2}{}\mathrm{\Delta }_{a_k}\underset{\sigma :j\sigma i}{}\epsilon (\sigma ,(a_k))\underset{k\sigma _1}{}\mathrm{\Delta }_{a_k}\underset{k\sigma _2}{}\mathrm{\Delta }_{a_k}=0.$$
$`(3.17)`$
Here $`\sigma `$ runs over 2–partitions of $`\{1,\mathrm{},n\}`$. We choose additionally an arbitrary ordering of both parts $`\sigma _1,\sigma _2`$ determining the ordering of $`\mathrm{\Delta }`$’s in the angular brackets, and compensate this choice by the $`\pm 1`$–factor $`\epsilon (\sigma ,(a_k))`$.
Conversely, any family of elements $`\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}\mathrm{End}F`$ defined for all $`n`$ and $`(a_1,\mathrm{},a_n)`$ and satisfying (3.16), (3.17) is the family of top matrix correlators of a well defined representation $`K:H_{}T\mathrm{End}F`$.
Proof. Clearly, (3.16) is a particular case of (3.13). To get (3.17), we apply (3.15) to the identical partition $`\tau ^{(n)}=\epsilon ^{(n)}`$ and then replace each term by the double product of top correlators using (3.14).
Conversely, assume that we are given $`\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}`$ satisfying (3.16) and (3.17). There is a unique way to extend this system to a family of elements $`\tau ^{(n)}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}`$ defined for all $`N`$–partitions $`\tau ^{(n)}`$ and satisfying the factorization property (3.14) and at least a part of the symmetry relations (3.13):
$$\tau ^{(n)}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}:=\epsilon (\tau ^{(n)},(a_k))\underset{r=1}{\overset{r=N}{}}\underset{k\tau _r^{(n)}}{}\mathrm{\Delta }_{a_k}.$$
$`(3.18)`$
Here, as in (3.17), we choose arbitrary orderings of each $`\tau _r^{(n)}`$ and compensate this by the appropriate sign so that the result does not depend on the choices made. All the relations (3.13) become automatically satisfied with this definition. In fact, the left hand side of (3.13) puts into $`s^1(\tau ^{(n)})_r`$ those $`i`$ for which $`s(i)\tau _r^{(n)}`$ (see (3.10)) so that the expression of both sides of (3.13) through the top correlators consists of the same groups taken in the same order. The coincidence of the signs is left to the reader.
It remains to check that (3.18) satisfy the linear relations (3.15). Recall now that to write a concrete relation (3.15) down we choose $`\tau ^{(n)},r`$, $`i,j\tau _r^{(n)}`$ and $`(a_1,\mathrm{},a_n)`$ and then sum over 2–partitions $`\alpha `$ of $`\tau _r^{(n)}`$. Hence replacing each term of the left hand side of (3.15) by the prescriptions (3.18) we get
$$\underset{p=1}{\overset{r1}{}}\underset{k\tau _p^{(n)}}{}\mathrm{\Delta }_{a_k}\left(\underset{\alpha :i\alpha j}{}\pm \underset{k\alpha _1}{}\mathrm{\Delta }_{a_k}\underset{k\alpha _2}{}\mathrm{\Delta }_{a_k}\underset{\alpha :j\alpha i}{}\pm \underset{k\alpha _1}{}\mathrm{\Delta }_{a_k}\underset{k\alpha _2}{}\mathrm{\Delta }_{a_k}\right)$$
$$\underset{q=r+1}{\overset{N}{}}\underset{k\tau _q^{(n)}}{}\mathrm{\Delta }_{a_k}.$$
This expression vanishes because its middle term is an instance of (3.17).
3.6. Precise statement and proof of the Theorem 3.3.1. Assume that we are given a representation $`K:H_{}T\mathrm{End}F.`$ We will produce from it a formal solution of the Commutativity Equations using only its top correlators. Let $`(x^a)`$ be the basis of formal coordinates on $`T`$ dual to $`(\mathrm{\Delta }_a)`$. Put
$$=\underset{n=1}{\overset{\mathrm{}}{}}\underset{(a_1,\mathrm{},a_n)}{}\frac{x^{a_n}\mathrm{}x^{a_1}}{n!}\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}k[[x]]\mathrm{End}F.$$
$`(3.19)`$
###### \quad3.6.1. Proposition
a) We have
$$dd=0.$$
$`(3.20)`$
b) Conversely, let $`\mathrm{\Delta }(a_1,\mathrm{},a_n)\mathrm{End}F`$ be a family of linear operators defined for all $`n1`$ and all maps $`\{1,\mathrm{},n\}I:ia_i`$. Assume that the parity of $`\mathrm{\Delta }(a_1,\mathrm{},a_n)`$ coincides with the sum of the parities of $`\mathrm{\Delta }_{a_i}`$ and that for any $`s𝕊_n`$
$$\mathrm{\Delta }(a_{s(1)},\mathrm{},a_{s(n)})=\epsilon (s,(a_i))\mathrm{\Delta }(a_1,\mathrm{},a_n).$$
Finally, assume that the formal series
$$=\underset{n=1}{\overset{\mathrm{}}{}}\underset{(a_1,\mathrm{},a_n)}{}\frac{x^{a_n}\mathrm{}x^{a_1}}{n!}\mathrm{\Delta }(a_1,\mathrm{},a_n)k[[x]]\mathrm{End}F$$
$`(3.21)`$
satisfies the equations (3.20). Then there exists a well defined representation $`K:H_{}T\mathrm{End}F`$ such that $`\mathrm{\Delta }(a_1,\mathrm{},a_n)`$ are the top correlators $`\mathrm{\Delta }_{a_1}\mathrm{}\mathrm{\Delta }_{a_n}`$ of this representation.
Notice that any even element of $`k[[x]]\mathrm{End}F`$ without constant term can be uniquely written in the form (3.21).
Proof. Clearly, the equations $`dd=0`$ written for the series (3.21) are equivalent to a family of bilinear relations between the symmetric matrix–valued tensors $`\mathrm{\Delta }(a_1\mathrm{}a_n)`$. In view of the Proposition 3.5.1, it remains to check only that this family of relations is equivalent to the family (3.17). This is a straightforward exercise.
§4. Stacks $`\overline{L}_{g;A,B}`$ and the extended modular operad
4.1. Introduction. The basic topological operad $`(\overline{M}_{0;n+1},n2)`$ of Quantum Cohomology lacks the $`n=1`$ term which is usually formally defined as a point. We argued elsewhere (cf. \[MaZ\], sec. 7 and \[Ma\], VI.7.6) that it would be very desirable to find a non–trivial DM–stack which could play the role of $`\overline{M}_{0;2}`$. There are several tests that such an object should pass:
a) It must be a semigroup (because for any operad $`𝒫`$, the operadic multiplication makes a semigroup of $`𝒫(1)`$).
b) It must be a part of an extended genus zero operad, say, $`(\stackrel{~}{L}_{0;n+1},n1)`$ geometrically related to $`(\overline{M}_{0;n+1},n2)`$ in such a way that the theory of Gromov–Witten invariants with gravitational descendants could be formulated in this new context. In particular, it must geometrically explain two–point correlators with gravitational descendants.
c) In turn, the extended genus zero operad must be a part of an extended modular operad containing moduli spaces of arbitrary genus, in such a way that algebras over classical modular operads produce extended algebras.
In this section we will try to show that the space
$$\stackrel{~}{L}_{0;2}:=\underset{n1}{}\overline{L}_n$$
$`(4.1)`$
passes at least a part of these tests. (Another candidate which might be interesting is $`\mathrm{lim}\mathrm{proj}\overline{L}_n`$ with respect to the forgetful morphisms).
4.2. Semigroup structure. It is defined as the union of boundary (clutching) morphisms
$$b:=(b_{n_1,n_2}):\stackrel{~}{L}_{0;2}\times \stackrel{~}{L}_{0;2}\stackrel{~}{L}_{0;2}$$
$`(4.2)`$
where
$$b_{n_1,n_2}:\overline{L}_{n_1}\times \overline{L}_{n_2}\overline{L}_{n_1+n_2}$$
glues $`x_{\mathrm{}}`$ of the first curve to $`x_0`$ of the second curve and renumbers the black points of the second curve keeping their order (cf. \[MaZ\], section 7). This is the structure that induced our multiplication on $`H_{}`$ in 3.3 above.
4.3. Extended operads. In (4.2), only white points $`\{x_0,x_{\mathrm{}}\}`$ are used to define the operadic composition whereas the black ones serve only to stabilize the strings of $`^1`$’s which otherwise would be unstable. This is a key observation for our attempt to define an extended operad.
A natural idea would be to proceed as follows. Denote by $`\overline{M}_{g;A,B}`$ the stack of stable $`(A,B)`$–pointed curves of genus $`g`$ (see Definition 1.1). Check that it is a DM–stack. Put $`\stackrel{~}{M}_{g;m+1}:=_{n0}\overline{M}_{g;m+1,n}`$ and define the operadic compositions via boundary maps, using only white points as above. (We sometimes write here and below $`n`$ instead of $`\{1,\mathrm{},n\}`$).
It seems however that this object is too big for our purposes and that it must be replaced by a smaller stack which we will define inductively, by using the Construction 1.3 which we will call here simply the adjoining of a generic black point. The components of this stack will be defined inductively.
If $`g2`$, $`m0`$, we start with $`\overline{M}_{g;m}=\overline{M}_{g;m,\mathrm{}}`$ and add $`n`$ generic black points, one in turn. Denote the resulting stack by $`\overline{L}_{g;m,n}.`$
For $`g=1`$, one should add one more sequence of stacks, corresponding to $`m=0`$. Since we want to restrict ourselves to Deligne–Mumford stacks, we start at $`\overline{M}_{1;0,1}`$ identified with $`\overline{M}_{1;1}`$ (see 1.2 a)), and add black points to get the sequence $`\overline{L}_{1;0,n}`$, $`n1.`$ These spaces are needed to serve as targets for the clutching morphisms gluing $`x_0`$ to $`x_{\mathrm{}}`$ on the same curve of genus zero: cf. below.
Finally, for $`g=0`$ we obtain our series of spaces $`\overline{L}_n=\overline{L}_{0;2,n}`$, $`n1`$ and moreover $`\overline{L}_{0;m,n}`$, for all $`m3,n0`$.
4.3.1. Combinatorial types of fibers. Let us remind that combinatorial types of classical (semi)stable curves with (only white) points labeled by a finite set $`A`$ are isomorphism classes of graphs, whose vertices are labeled by “genera” $`g0`$ and tails are bijectively labeled by elements of $`A`$. Stability means that vertices of genus 0 bound $`3`$ flags, and vertices of genus 1 bound $`1`$ flags. Graphs can have edges with only one vertex, that is, simple loops. See \[Ma\], III.2 for more details.
Starting with such a graph $`\mathrm{\Gamma }`$, or rather with its geometric realization, we can obtain an infinite series of graphs, which will turn out to be exactly combinatorial types of (semi)stable $`(A,B)`$–pointed curves that are fibers of the families described above. Namely, subdivide edges and tails of $`\mathrm{\Gamma }`$ by a finite set of new vertices of genus zero (on each edge or tail, this set may be empty). If a tail was subdivided, move the respective label (from $`A`$) to the newly emerged tail. Distribute the black tails labeled by elements of $`B`$ arbitrarily among the old and the new vertices. Call the resulting graph a stringy stable combinatorial type if it becomes stable after repainting black tails into white ones. Clearly, new vertices depict strings of $`^1`$’s stabilized by black points and eventually two special points on the end components.
###### \quad4.3.2. Theorem
a) $`\overline{L}_{g;m,n}`$ is the Deligne–Mumford stack classifying $`(m,n)`$–pointed curves of genus $`g`$ of stringy stable combinatorial types. It is proper and smooth.
b) Therefore, one can define boundary morphisms gluing two white points of two different curves
$$\overline{L}_{g_1;m_1+1,n_1}\times \overline{L}_{g_2;m_2+1,n_2}\overline{L}_{g_1+g_2;m_1+m_2+1,n_1+n_2}$$
and gluing two white points of the same curve:
$$\overline{L}_{g;m+1,n}\overline{L}_{g+1,m1;n}$$
such that the locally finite DM–stacks
$$\stackrel{~}{L}_{g,m+1}:=\underset{n0}{}\overline{L}_{g;m+1,n}$$
will form components of a modular operad.
The statement a) can be proved in the same way as the respective statement 2.2 a).
It remains to see whether one can develop an extension of the Gromov–Witten invariants, preferably with descendants, to this context. The Remark 3.2.3 seems promising in this respect.
Appendix. Proof of the Technical Lemma
We break the proof into several steps whose content is indicated in the title of the corresponding subsection. An advice for the reader who might care to check the details: the most daunting task is to convince oneself that none of the alternatives has been inadvertently omitted.
A.1. The right hand side of (2.24) does not depend on the choice of $`i,j`$.
We must check that a different choice leads to the same answer modulo relations (2.23). We can pass from one choice to another by consecutively replacing only one element of the pair. Consider, say, the passage from $`(i,j)`$ to $`(i^{},j)`$. Form the difference of the right hand sides of (2.24) written for $`(i^{},j)`$ and for $`(i,j)`$.
In this difference, the terms corresponding to the partitions $`\beta `$ will cancel. The remaining terms will correspond to the partitions $`\alpha `$ of $`\tau _a`$ which separate $`i`$ and $`i^{}`$. Their difference will vanish in $`H_B`$ because of (2.23).
A.2. Multiplications by $`l_\sigma `$ are compatible with linear relations (2.23) between $`\mu (\tau )`$.
Choose and fix one linear relation (2.23), that is, a quadruple $`(\tau ,\tau _a,i,j\tau _a)`$, $`ij`$. Choose also a 2–partition $`\sigma `$. We want to check that after multiplying the left hand side of (2.23) by $`l_\sigma `$ according to the prescriptions (2.23)–(2.26) we will get zero modulo all relations of the type (2.23). There are several basic cases to consider.
(i) $`\sigma `$ breaks $`\tau `$ at $`\tau _b`$, $`ba`$. Then put $`\tau ^{}=\sigma \tau `$. After multiplication we will get again (2.23) written for $`\tau ^{}`$ and one of its components $`\tau _a`$.
(ii) $`\sigma `$ breaks $`\tau `$ at $`\tau _a`$. Let $`(\tau _{a1},\tau _{a2})`$ be the induced partition; it is now fixed. We must analyze $`l_\sigma \mu (\tau (\alpha ))`$ for variable 2–partitions $`\alpha `$ of $`\tau _a`$ with $`i\alpha j`$ or $`j\alpha i`$.
Those $`\alpha `$ which do not break $`(\tau _{a1},\tau _{a2})`$ will contribute zero because of (2.26).
Those $`\alpha `$ which break $`(\tau _{a1},\tau _{a2})`$ will produce a 3–partition of $`\tau _a`$, say $`(\tau _{a11},\tau _{a12},\tau _{a2})`$ or else $`(\tau _{a1},\tau _{a21},\tau _{a22})`$. Finally, there will be one $`\alpha `$ which is induced by $`\sigma `$ that is, coincides with $`(\tau _{a1},\tau _{a2})`$. We must show that the sum total of the respective terms vanishes. However, the pattern of cancellation will depend on the positions of $`i`$ and $`j`$. In order to present the argument more concisely, we will first introduce the numerotation of all possible positions with respect to a variable $`\alpha `$ as follows. Partitions which break $`(\tau _{a1},\tau _{a2})`$ with $`i\alpha j`$:
$$(\mathrm{I}):i\tau _{a11},j\tau _{a12}(\mathrm{II}):i\tau _{a11},j\tau _{a2}$$
$$(\mathrm{III}):i\tau _{a1},j\tau _{a22}(\mathrm{IV}):i\tau _{a21},j\tau _{a22}$$
Partitions which break $`(\tau _{a1},\tau _{a2})`$ and satisfy $`j\alpha i`$ will be denoted similarly, but with prime. Say, (III) means (III) with positions of $`i`$ and $`j`$ reversed.
Now we will explain the patterns of cancellation depending on the positions of $`i,j`$ with respect to $`\sigma `$. Recall that this latter data is fixed and determined by the choices we made at the beginning of this subsection.
If $`i,j\tau _{a1}`$, the only non–vanishing terms are of the types (I) and (I). Their sum over all $`\alpha `$ will vanish because of (2.23). Similarly, if $`i,j\tau _{a2}`$, (IV) and (IV) will cancel, and everything else will vanish.
Finally, assume that $`i\tau _{a1},j\tau _{a2}`$. that is, $`\sigma `$ separates $`i,j`$. Then we may have non–vanishing terms of the types (II) and (III) and in addition the terms coming from (the partition of $`\tau _a`$ induced by) $`\sigma `$ which must be treated using the formula (2.24), applied however to $`(\tau _1,\mathrm{},\tau _{a1},\tau _{a1},\tau _{a2},\tau _{a+1},\mathrm{})`$ in place of $`\tau `$. Half of these latter terms (with $`i\tau _{a11}`$) will cancel (II), whereas the other half (with $`j\tau _{a22}`$) will cancel (III).
The case $`j\tau _{a1},i\tau _{a2}`$ is treated similarly.
(iii) $`\sigma `$ breaks $`\tau `$ between $`\tau _b`$ and $`\tau _{b+1}`$. In this case $`\sigma `$ breaks any $`\tau (\alpha )`$ in (2.23) between two neighbors as well. A contemplation will convince the reader that only the cases $`b=a1`$ and $`b=a`$ may present non–obvious cancellations. Let us treat the first one; the second one is simpler.
For $`\alpha =(\tau _{a1},\tau _{a2})`$ we will calculate each term $`l_\sigma \mu (\tau (\alpha ))`$ using a formula of the type (2.24), first choosing some $`k\tau _{a1},l\tau _{a1}`$ (in place of $`i,j`$ of (2.24): these letters are already bound). The choice of $`k`$ does not matter, but we will choose $`l=i`$ if $`i\alpha j`$, and $`l=j`$ if $`j\alpha i`$. We get then for $`i\alpha j`$:
$$l_\sigma \mu (\tau (\alpha ))=l_\sigma \mu (\mathrm{}\tau _{a1}\tau _{a1}\tau _{a2}\mathrm{})$$
$$=\underset{\beta :k\tau _{a1,1}}{}\mu (\mathrm{}\tau _{a1,1}\tau _{a1,2}\tau _{a1}\tau _{a2}\mathrm{})\underset{\gamma :i\tau _{a12}}{}\mu (\mathrm{}\tau _{a1}\tau _{a11}\tau _{a12}\tau _{a2}\mathrm{})$$
$`(A.1)`$
where $`\beta `$ runs over 2–partitions of $`\tau _{a1}`$ and $`\gamma `$ runs over 2–partitions of $`\tau _{a1}.`$ Write now a similar expression for $`j\alpha i`$ (with the choice $`l=j`$). The second sum in this expression will term–by–term cancel the second sum in (A.1), because our choices force $`i\tau _{a12},j\tau _{a2}`$ in both cases.
If we sum first over $`\alpha `$, we will see that the first two sums cancel modulo relations (2.23) because our choices imply $`i\tau _{a1},j\tau _{a2}`$ in the first sum of (A.1) and the reverse relation in the first sum written for $`j\alpha i.`$
(iv) $`\sigma `$ does not break $`\tau `$. In this case we choose a bad pair $`(\tau _b,\tau _{b+1})`$ for $`\sigma `$ and $`\tau `$ (see Lemma 2.4.2(iii)). One easily sees that if $`ab,b+1`$, then it remains a bad pair for $`\sigma `$ and $`\tau (\alpha )`$ for any $`\alpha `$ in (2.23). Therefore, $`ł_\sigma `$ annihilates all terms of (2.23) in view of (2.26).
We will show that in the exceptional cases we still can find a bad pair for $`\sigma `$ and $`\tau (\alpha )`$, but it will depend on $`\alpha =(\tau _{a1},\tau _{a2})`$, which does not change the remaining argument.
Assume that $`b=a`$ that is, $`\tau _a\sigma _1\mathrm{},\tau _{a+1}\sigma _1\mathrm{}`$ (see (2.12)). Then $`(\tau _{a2},\tau _{a+1})`$ is a bad pair for $`\sigma `$ and $`\tau (\alpha )`$, unless $`\tau _{a2}\sigma _1`$, in which case $`\sigma _1`$ cannot contain $`\tau _{a1}`$ so that $`(\tau _{a1},\tau _{a2})`$ form a bad pair.
Similarly, if $`b=a1`$, then $`\tau _{a1},\tau _{a1}`$ will be a bad pair unless $`\tau _{a1}\sigma _1=\mathrm{}`$, in which case $`(\tau _{a1},\tau _{a2})`$ will be a bad pair.
By this time we have checked that multiplications by $`l_\sigma `$ are well defined linear operators on the space $`H_B`$. We will now prove that they pairwise commute and therefore define an action of $`_B`$ upon $`H_B`$.
A.3. Multiplications by $`l_\sigma `$ pairwise commute.
We start with fixing $`\tau `$, $`\sigma ^{(1)}`$ and $`\sigma ^{(2)}`$. We want to check that
$$l_{\sigma ^{(1)}}(l_{\sigma ^{(2)}}\mu (\tau ))=l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau )).$$
We may and will assume that $`\sigma ^{(1)}\sigma ^{(2)}`$. The following alternatives can occur for $`\sigma ^{(1)}`$ and $`\sigma ^{(2)}`$ separately:
(i) $`\sigma ^{(1)}`$ breaks $`\tau `$ at $`\tau _a`$.
(ii) $`\sigma ^{(1)}`$ breaks $`\tau `$ between $`\tau _a`$ and $`\tau _{a+1}`$.
(iii) $`\sigma ^{(1)}`$ does not break $`\tau `$.
(i) $`\sigma ^{(2)}`$ breaks $`\tau `$ at $`\tau _b`$.
(ii) $`\sigma ^{(2)}`$ breaks $`\tau `$ between $`\tau _b`$ and $`\tau _{b+1}`$.
(iii) $`\sigma ^{(2)}`$ does not break $`\tau `$.
We will have to consider the combined alternatives (i)(i), (i)(ii), $`\mathrm{}`$ , (iii)(iii) in turn. The symmetry of $`\sigma ^{(1)}`$ and $`\sigma ^{(2)}`$ allows us to discard a few of them.
Subcase (i)(i)
We will first assume that $`ab`$, say $`a<b`$. Denote by $`\alpha `$ (resp. $`\beta `$) the partition induced by $`\sigma ^{(1)}`$ (resp. $`\sigma ^{(2)}`$) on $`\tau _a`$ (resp. $`\tau _b`$). Then
$$l_{\sigma ^{(1)}}(l_{\sigma ^{(2)}}\mu (\tau ))=l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau ))=\mu (\tau (\alpha )(\beta ))=\mu (\tau (\beta )(\alpha )).$$
Now assume that $`a=b`$. If $`\alpha `$ breaks $`\beta `$, we will again have the desired equality, because $`\alpha \beta ==\beta \alpha `$. If $`\alpha `$ does not break $`\beta `$, the both sides will vanish.
After having treated this subcase, we add one more remark which will be used below, in the subsection A.5. Namely, $`\alpha `$ does not break $`\beta `$ exactly when $`\sigma ^{(1)}`$ does not break $`\sigma ^{(2)}`$. Therefore, if $`l_{\sigma ^{(1)}}l_{\sigma ^{(2)}}`$ is one of the quadratic generators of $`I_B`$, then consecutive multiplication by the respective elements annihilates $`\mu (\tau )`$.
Subcase (i)(ii)
If $`ab,b+1`$, the modifications induced in $`\tau `$ by the two multiplications are made in mutually disjoint places and therefore commute as above. Consider now the case $`a=b`$, the case $`a=b+1`$ being similar.
Denote by $`(\tau _{a1},\tau _{a2})`$ the partition induced by $`\sigma `$ on $`\tau _a.`$ Then we have
$$l_{\sigma ^{(1)}}\mu (\tau )=\mu (\mathrm{}\tau _{a1}\tau _{a1}\tau _{a2}\tau _{a+1}\mathrm{})=\mu (\tau ^{}).$$
Clearly, $`\sigma ^{(2)}`$ breaks $`\tau ^{}`$ between $`\tau _{a2}`$ and $`\tau _{a+1}`$ so that, after choosing $`i\tau _{a2},j\tau _{a+1}`$ we have
$$l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau ))=\underset{\alpha :i\alpha }{}\mu (\tau ^{}(\alpha ))\underset{\beta :\beta j}{}\mu (\tau ^{}(\beta ))$$
$`(A.2)`$
where $`\alpha `$ runs over 2–partitions of $`\tau _{a2}`$ and $`\beta `$ runs over 2–partitions of $`\tau _{a+1}`$.
On the other hand, with the same choice of $`i,j`$ we have
$$l_{\sigma ^{(2)}}\mu (\tau )=\underset{\gamma :i\gamma }{}\mu (\tau (\gamma ))\underset{\beta :\beta j}{}\mu (\tau (\beta ))$$
$`(A.3)`$
where $`\gamma `$ runs over 2–partitions of $`\tau _a`$ and $`\beta `$ runs over 2–partitions of $`\tau _{a+1}`$. After multiplication of (A.3) by $`l_{\sigma ^{(1)}}`$, the second sum in (A.3) will become the second sum of (A.2). In the first sum, only partitions $`\delta `$ breaking $`(\tau _{a1},\tau _{a2})`$ will survive, and they will produce exactly the first sum in (A.2).
Subcase (i)(iii)
Here $`\sigma ^{(1)}`$ breaks $`\tau `$ at $`\tau _a`$, and there exists a bad pair $`(\tau _b,\tau _{b+1})`$ for $`\sigma ^{(2)}`$ and $`\tau `$. Since $`l_{\sigma ^{(2)}}\mu (\tau )=0`$, it remains to check that $`l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau ))=0`$. But $`l_{\sigma ^{(1)}}\mu (\tau )=\mu (\tau ^{})`$ as in the previous subcase. So it remains to find a bad pair for $`\sigma ^{(2)}`$ and $`\tau ^{}`$.
If $`ab,b+1`$, then $`(\tau _b,\tau _{b+1})`$ will will be such a bad pair
If $`a=b`$, denote by $`(\tau _{a1},\tau _{a2})`$ the partition of $`\tau _a`$ induced by $`\sigma ^{(1)}`$. If $`\tau _{a2}`$ is not contained in $`\sigma ^{(2)}`$, $`(\tau _{a2},\tau _{a+1})`$ will form a bad pair. Otherwise this role will pass to $`(\tau _{a1},\tau _{a2})`$.
Finally, consider the case when $`a=b+1`$. In the previous notation, if $`\sigma _1^{(2)}\tau _{a1}\mathrm{}`$, then $`(\tau _{a1},\tau _{a1})`$ is the bad pair we are looking for, otherwise we should take $`(\tau _{a1},\tau _{a2})`$.
Subcase (ii)(ii)
Here $`\sigma ^{(1)}`$ (resp. $`\sigma ^{(2)}`$) breaks $`\tau `$ between $`a`$ and $`a+1`$ (resp. between $`b`$ and $`b+1`$), and $`ab`$.
If $`ab1,b+1`$, the modifications indiced in $`\tau `$ by $`\sigma ^{(1)}`$ and $`\sigma ^{(2)}`$ do not interact and the respective multiplications commute.
By symmetry, it remains to consider the case $`a=b1`$. Choose $`i\tau _a,j\tau _{a+1}.`$ Summing first over partitions $`\alpha =(\tau _{a1},\tau _{a2})`$ and $`\beta =(\tau _{a+1,1},\tau _{a+1,2})`$ we have
$$l_{\sigma ^{(1)}}\mu (\tau )=\underset{\alpha :i\alpha }{}\mu (\mathrm{}\tau _{a1}\tau _{a2}\mathrm{})\underset{\beta :\beta j}{}\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2}\mathrm{}).$$
Now, $`\sigma ^{(2)}`$ will break the terms of the first (resp. second) sum between $`\tau _{a+1}`$ and $`\tau _{a+2}`$ (resp. between $`\tau _{a+1,2}`$ and $`\tau _{a+2}`$). In order to multiply by $`l_{\sigma ^{(2)}}`$ each term of these sums we choose the same $`j\tau _{a+1}`$ and some $`l\tau _{a+2}`$. Below we sum additionally over 2–partitions $`\beta =(\tau _{a+1,1},\tau _{a+1,2})`$ and $`\gamma =(\tau _{a+2,1},\tau _{a+2,2})`$ in the first two sums. In the second two sums the respective notation is $`\beta ^{}=(\tau _{a+1,2,1},\tau _{a+1,2,2})`$:
$$l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau ))=$$
$$\frac{=}{\frac{\alpha :i\alpha \beta :j\beta \mu (\mathrm{}\tau _{a1}\tau _{a2}\tau _{a+1,1}\tau _{a+1,2}\mathrm{})+}{\alpha :i\alpha \gamma :\gamma l\mu (\mathrm{}\tau _{a1}\tau _{a2}\tau _{a+1}\tau _{a+2,1}\tau _{a+2,2}\mathrm{})}}$$
$$\frac{+}{\frac{\beta :\beta j\beta ^{}:j\beta ^{}\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2,1}\tau _{a+1,2,2}\tau _{a+2}\mathrm{})+}{\beta :\beta j\gamma :\gamma l\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2}\tau _{a+2,1}\tau _{a+2,2}\mathrm{})}}$$
$`\left(A.4\right)`$
On the other hand, with the same notation we have:
$$l_{\sigma ^{(2)}}\mu (\tau )=\underset{\beta :j\beta }{}\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2}\mathrm{})\underset{\gamma :\gamma l}{}\mu (\mathrm{}\tau _{a+2,1}\tau _{a+2,2}\mathrm{})$$
and
$$l_{\sigma ^{(1)}}(l_{\sigma ^{(2)}}\mu (\tau ))=$$
$$\frac{=}{\frac{\alpha :i\alpha \beta :j\beta \mu (\mathrm{}\tau _{a1}\tau _{a2}\tau _{a+1,1}\tau _{a+1,2}\mathrm{})+}{\beta :j\beta \beta ^{\prime \prime }:\beta ^{\prime \prime }j\mu (\mathrm{}\tau _{a+1,1,1}\tau _{a+1,1,2}\tau _{a+1,2}\mathrm{})}}$$
$$\frac{+}{\frac{\alpha :i\alpha \gamma :\gamma l\mu (\mathrm{}\tau _{a1}\tau _{a2}\tau _{a+1}\tau _{a+2,1}\tau _{a+2,2}\mathrm{})+}{\beta :\beta j\gamma :\gamma l\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2}\tau _{a+2,1}\tau _{a+2,2}\mathrm{})}}$$
$`\left(A.5\right)`$
where $`\beta ^{\prime \prime }=(\tau _{a+1,1,1},\tau _{a+1,1,2})`$. Three of the four sums in (A.4) and (A.5) obviously coincide. The third sum in (A.4) coincides with the second sum in (A.5) because both are taken over 3–partitions of $`\tau _{a+1}`$ with $`j`$ in the middle part.
Subcase (ii)(iii)
Here $`\sigma ^{(1)}`$ breaks $`\tau `$ between $`a`$ and $`a+1`$, $`\sigma ^{(2)}`$ does not break $`\tau `$. We must check that $`l_{\sigma ^{(2)}}(l_{\sigma ^{(1)}}\mu (\tau ))=0`$, by finding a bad pair for $`\sigma ^{(2)}`$ and each term in the right hand side of
$$l_{\sigma ^{(1)}}\mu (\tau )=\underset{\alpha :i\alpha }{}\mu (\mathrm{}\tau _{a1}\tau _{a2}\mathrm{})\underset{\beta :\beta j}{}\mu (\mathrm{}\tau _{a+1,1}\tau _{a+1,2}\mathrm{}).$$
Denote by $`(\tau _b,\tau _{b+1})`$ a bad pair for $`\sigma ^{(2)}`$ and $`\tau `$. As in the subcase (i)(iii), it will remain the bad pair unless $`b\{a1,a,a+1\},`$ and will change somewhat in the exceptional cases.
More preciasely, if $`b=a1`$, then for the terms of the second sum $`(\tau _{a1},\tau _a)`$ will be bad. For the first sum, if $`\sigma _1^{(2)}\tau _{a1}\mathrm{}`$, the bad pair will be $`(\tau _{a1},\tau _{a1})`$. Otherwise it will be $`(\tau _{a1},\tau _{a2})`$.
If $`b=a`$, then for the terms of the first sum $`(\tau _{a2},\tau _{a+1})`$ will be bad. For the second sum, if $`\sigma _1^{(2)}\tau _{a+1,1}\mathrm{}`$, the bad pair will be $`(\tau _a,\tau _{a+1,1})`$. Otherwise it will be $`(\tau _{a+1,1},\tau _{a+1,2})`$.
Finally, if $`b=a+1`$, then for the terms of the first sum $`(\tau _{a+1},\tau _{a+2})`$ will be bad. For the second sum, it will be $`(\tau _{a+1,2},\tau _{a+2})`$.
In the last remaining case (iii)(iii) both multiplications produce zero.
To complete the proof of the Technical Lemma, it now remains to check that the elements (2.14), (2.15) generating $`I_B`$ annihilate $`H_B`$.
A.4. Elements $`r_{ij}^{(1)}`$ annihilate $`H_B`$.
Fix $`i,j`$ and a partition $`\tau `$. If $`\tau `$ does not separate $`i`$ and $`j`$, we have $`i,j\tau _a`$ for some $`a`$, and then
$$r_{ij}^{(1)}\mu (\tau )=\left(\underset{\sigma :i\sigma j}{}l_\sigma \underset{\sigma :j\sigma i}{}l_\sigma \right)\mu (\tau )$$
$$=\underset{\alpha :i\alpha j}{}\mu (\tau (\alpha ))\underset{\alpha :j\alpha i}{}\mu (\tau (\alpha ))$$
$`(A.6)`$
where $`\alpha `$ runs over partitions of $`\tau _a`$. This expression vanishes because of (2.23).
Assume now that $`\tau `$ separates $`i`$ and $`j`$, say, $`i\tau _a,j\tau _b,a<b.`$ In this case $`l_\sigma \mu (\tau )=0`$ for all $`\sigma `$ with $`j\sigma i.`$ The remaining terms of (A.6) vanish unless $`\sigma `$ breaks $`\tau `$ at some $`\tau _c,acb`$, or else between $`\tau _c`$ and $`\tau _{c+1}`$ for $`acb1.`$ In the latter cases each term corresponding to one $`\sigma `$ can be replaced by a sum of terms corresponding to the 2–partitions $`\alpha _c`$ of $`\tau _c`$ with the help of (2.24) and (2.25).
Let us choose $`k_c\tau _c`$ for all $`acb`$ so that $`k_a=i,k_b=j`$ and spell out the resulting expression:
$$\underset{\sigma :i\sigma j}{}l_\sigma \mu (\tau )=\underset{c=a}{\overset{c=b}{}}{}_{}{}^{}(\underset{\alpha _c:k_c\alpha _c}{}+\underset{\alpha _c:\alpha _ck_c}{})\mu (\tau (\alpha _c))$$
$$\underset{c=a}{\overset{c=b1}{}}\left(\underset{\alpha _c:k_c\alpha _c}{}+\underset{\alpha _{c+1}:\alpha _{c+1}k_{c+1}}{}\right)\mu (\tau (\alpha _{c+1})).$$
Here prime at the first sum indicates that the terms with $`\alpha _ai`$ and $`j\alpha _b`$ should be skipped.
All the terms of this expression cancel.
A.5. Elements $`r^{(2)}(\sigma ^{(1)},\sigma ^{(2)})`$ annihilate $`H_B`$.
These elements correspond to the pairs ($`\sigma ^{(1)},\sigma ^{(2)}`$) that do not break each other. If at least one of them, say $`\sigma ^{(1)}`$, does not break $`\tau `$ either, then $`l_{\sigma ^{(1)}}\mu (\tau )=0`$ so that $`r^{(2)}(\sigma ^{(1)},\sigma ^{(2)})\mu (\tau )=0.`$ If both $`\sigma ^{(1)},\sigma ^{(2)}`$ break $`\tau `$, a contemplation will convince the reader that they must break $`\tau `$ at one and the same component $`\tau _a.`$ This is the subcase (i)(i) of A.3, and we made the relevant comment there.
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warning/0001/math-ph0001021.html | ar5iv | text | # Transformations of ordinary differential equations via Darboux transformation technique
## 1. Introduction
The nonlinear ordinary differential equations (ODE’s) admitting a representation as the compatibility condition of the overdetermined system of linear equations (Lax pair) are intensively studied in modern mathematical physics. The most famous equations belonging to this class are the six Painlev$`\stackrel{´}{\text{e}}`$ equations $`\mathrm{P}_\mathrm{I}`$$`\mathrm{P}_{\mathrm{VI}}`$ . It was found that these equations are closely connected with nonlinear partial differential equations (PDE’s) integrable in the frameworks of the inverse scattering transformation (IST) method . Various approaches developed in the theory of nonlinear integrable PDE’s were applied to study ODE’s of such a class (e.g., see reviews and references therein).
It is well known that $`\mathrm{P}_{\mathrm{II}}`$$`\mathrm{P}_{\mathrm{VI}}`$ have the transformations that map the solutions of a given Painlev$`\stackrel{´}{\text{e}}`$ equation into solutions of the same equation but with different values of the parameters of the equation. Such transformations for $`\mathrm{P}_{\mathrm{II}}`$ were found in as a generalization of corresponding formulas of the rational solutions transformations . The transformations of solutions of $`\mathrm{P}_{\mathrm{II}}`$ were also derived by means of the Bäcklund transformation technique . A representation of the Painlev$`\stackrel{´}{\text{e}}`$ equations as the systems of the first order equations was used to obtain the transformations of $`\mathrm{P}_{\mathrm{III}}`$ in the cases of two different choices of values of the parameters , $`\mathrm{P}_{\mathrm{IV}}`$ and of $`\mathrm{P}_\mathrm{V}`$ with special parameter values . In similar manner the transformations for $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_{\mathrm{IV}}`$ were obtained in . Moreover, an equivalence of nonlinear ODE to a system of equations was used to construct the transformations of solutions of the higher order nonlinear ODE’s appearing as the self–similar reduction in the KdV hierarchy . An unified approach to investigate the transformation properties of Painlev$`\stackrel{´}{\text{e}}`$ equations, which utilizes the singularity structure of their solutions, were developed in . The known transformations and new ones for $`\mathrm{P}_\mathrm{V}`$, that differs from the equation considered in by a region of the parameters values, and for $`\mathrm{P}_{\mathrm{VI}}`$ were derived. The truncation method was modified to find the transformations of nonlinear ODE’s . Its extension, which is based on considering the transformations that preserve the locations of subsets of the singularities of solutions, was suggested . The transformations of solutions of $`\mathrm{P}_{\mathrm{II}}`$ and $`\mathrm{P}_{\mathrm{IV}}`$ were, in particular, obtained using this approach.
The transformations of nonlinear ODE representable as the compatibility condition of the Lax pair cause evidently the transformations of the Lax pair solutions. Studying the monodromy preserving deformations of matrix linear ODE’s with regular singularities, Schlesinger constructed, in pure algebraic way, the transformations that keep the monodromy data of the solutions except for the exponents of formal monodromy. These transformations, which are now referred to as the Schlesinger transformations, were generalized for linear equations with irregular singular points by Jimbo and Miwa . The monodromy preserving condition is equivalent to having a set of additional linear ODE’s on solution of the isomonodromic equation , which establish a hierarchy of the Lax pairs. However, the authors mentioned seem to apply no the Schlesinger transformations for the derivation of transformations for nonlinear ODE’s (such as Painlev$`\stackrel{´}{\text{e}}`$ equations; see , p. 437) that admit the compatibility condition representation. The transformations of $`\mathrm{P}_{\mathrm{II}}`$$`\mathrm{P}_{\mathrm{VI}}`$, including new ones, were obtained by means of the Schlesinger transformations of associated Lax pairs in .
In the present paper we show that the Darboux transformation (DT) technique can be used to construct the transformations of solutions of nonlinear ODE’s and associated Lax pairs. This technique is suitable for generating the infinite hierarchies of solutions satisfying the reduction constraints imposed on the coefficients of the Lax pairs. To extend the DT technique for obtaining the transformations of ODE’s we suggest that the points, in which the DT possesses a kernel, coincide with the singular points of the Lax pairs. The cases of $`\mathrm{P}_{\mathrm{II}}`$ and certain nonlinear ODE of the second order are examined by way of illustration.
The paper is organized as follows. The Lax pairs of the ODE’s under consideration are written in Sec.2. We present the explicit formulas of the DT’s in Sec.3 and discuss the properties of them. The asymptotic expansions at the singular point neighborhood of the Lax pairs solutions are given in Sec.4. These expansions are used there to build the transformations of the ODE’s on the basis of the DT formulas.
## 2. Lax Pairs and Nonlinear ODE’s
Let us consider direct Lax pair
$$\{\begin{array}{c}\mathrm{\Psi }_x=P\mathrm{\Psi }\hfill \\ \mathrm{\Psi }_\lambda =Q\mathrm{\Psi }\hfill \end{array}$$
(1)
for matrix function $`\mathrm{\Psi }\mathrm{\Psi }(x,\lambda )`$ and dual Lax pair
$$\{\begin{array}{c}\mathrm{\Xi }_x=\mathrm{\Xi }P\hfill \\ \mathrm{\Xi }_\lambda =\mathrm{\Xi }Q\hfill \end{array}$$
(2)
for matrix $`\mathrm{\Xi }\mathrm{\Xi }(x,\lambda )`$. The compatibility condition of the direct pair $`\mathrm{\Psi }_{x\lambda }=\mathrm{\Psi }_{\lambda x}`$ (or, for dual one, $`\mathrm{\Xi }_{x\lambda }=\mathrm{\Xi }_{\lambda x}`$) leads to the condition on the matrix coefficients $`P`$ and $`Q`$:
$$P_\lambda Q_x+[P,Q]=0.$$
(3)
The equation $`\mathrm{P}_{\mathrm{II}}`$
$$u_{xx}=2u^3+xu\alpha $$
(4)
arises from the compatibility condition of Lax pairs (1,2) with coefficients
$$P=\lambda \sigma _3+u\sigma _1,$$
(5)
$$Q=4\lambda ^2\sigma _34\lambda u\sigma _1+\sigma _3\left(2u_x\sigma _12u^2ExE\right)+\lambda ^1\alpha \sigma _1,$$
(6)
where $`\sigma _i`$ ($`i=1,2,3`$) are the Pauli spin matrices, $`E`$ is $`22`$ unit matrix.
In the sequel the DT technique will be applied for obtaining the transformations of the solutions of $`\mathrm{P}_{\mathrm{II}}`$ and following ODE
$$x^2uu_{xx}=(xu_x)^2xuu_x+\alpha ^2u^42\gamma u^32(2\beta +1)xu4x^2$$
(7)
($`\alpha `$, $`\beta `$ and $`\gamma `$ are constants), for which the transformations were not written explicitly in the literature. This equation is connected with special cases of $`\mathrm{P}_{\mathrm{III}}`$ and $`\mathrm{P}_\mathrm{V}`$ .
Let us consider direct and dual Lax pairs (1,2) with coefficients
$$P=\lambda \sigma _3+\frac{1}{x}\left(\begin{array}{cc}0& v_1\\ v_2& 0\end{array}\right),$$
(8)
$$Q=x\sigma _3+\lambda ^1\left(\begin{array}{cc}\beta & v_1\\ v_2& \beta \end{array}\right)+\lambda ^2\left(\begin{array}{cc}q& w_1\\ w_2& q\end{array}\right),$$
(9)
where $`v_1=uzw,v_2={\displaystyle \frac{v(\alpha z)}{w}},w_1=zw,w_2={\displaystyle \frac{\alpha z}{w}},q=z{\displaystyle \frac{\alpha }{2}}.`$ The compatibility condition (3) in this case yields
$$\{\begin{array}{c}u_x=2\frac{2\beta }{x}u+\frac{2z\alpha }{x}u^2,\hfill \\ v_x=2+\frac{2\beta }{x}v\frac{2z\alpha }{x}v^2,\hfill \\ z_x=\frac{vu}{x}z(z\alpha ),\hfill \\ \frac{w_x}{w}=\frac{uz+v\left(z\alpha \right)}{x}.\hfill \end{array}$$
The first three equations are reduced to Eq.(7).
## 3. Darboux Transformation Technique
The first equations of Lax pairs (1,2) with coefficient $`P`$ given by Eq.(5) or Eq.(8) are well known in the theory of IST as direct and dual Zakharov–Shabat (ZS) spectral problems
$$\mathrm{\Psi }_x=(\lambda P^{(1)}+P^{(0)})\mathrm{\Psi },$$
(10)
$$\mathrm{\Xi }_x=\mathrm{\Xi }(\lambda P^{(1)}+P^{(0)}).$$
(11)
Potential $`P^{(0)}`$ of these problems has following general form:
$$P^{(0)}=\left(\begin{array}{cc}0& u_1\\ u_2& 0\end{array}\right).$$
The DT technique allows one to produce new solutions of the ZS problems and corresponding matrix potential, starting from the initial ones. Let $`\phi =(\phi _1,\phi _2)^T`$ be a vector solution of direct ZS problem (10) with $`\lambda =\mu `$. The first elementary DT (EDT) of direct problem $`\{\mathrm{\Psi },\mathrm{\Xi },P^{(0)}\}\{\stackrel{~}{\mathrm{\Psi }},\stackrel{~}{\mathrm{\Xi }},\stackrel{~}{P}^{(0)}\}`$ is defined as given :
$$\stackrel{~}{\mathrm{\Psi }}=\sigma \left(\begin{array}{cc}\lambda \mu +\frac{u_1\phi _2}{2\phi _1}& \frac{u_1}{2}\\ \frac{\phi _2}{\phi _1}& 1\end{array}\right)\mathrm{\Psi },\stackrel{~}{\mathrm{\Xi }}=\frac{\mathrm{\Xi }}{\sigma }\left(\begin{array}{cc}1& \frac{u_1}{2}\\ \frac{\phi _2}{\phi _1}& \lambda \mu +\frac{u_1\phi _2}{2\phi _1}\end{array}\right),$$
(12)
$$\stackrel{~}{P}^{(0)}=\left(\begin{array}{cc}0& \frac{u_{1,x}}{2}+\frac{u_1^2\phi _2}{2\phi _1}\mu u_1\\ 2\frac{\phi _2}{\phi _1}& 0\end{array}\right)$$
(13)
(it is assumed hereafter that $`\sigma `$ is a scalar function of $`\lambda `$). The direct and dual ZS problems are covariant with respect to this EDT: matrix functions $`\stackrel{~}{\mathrm{\Psi }}`$ and $`\stackrel{~}{\mathrm{\Xi }}`$ are the solutions of Eqs.(10,11) with potential $`\stackrel{~}{P}^{(0)}`$. The second EDT of direct problem has a form:
$$\stackrel{~}{\mathrm{\Psi }}=\sigma \left(\begin{array}{cc}1& \frac{\phi _1}{\phi _2}\\ \frac{u_2}{2}& \lambda \mu \frac{u_2\phi _1}{2\phi _2}\end{array}\right)\mathrm{\Psi },\stackrel{~}{\mathrm{\Xi }}=\frac{\mathrm{\Xi }}{\sigma }\left(\begin{array}{cc}\lambda \mu \frac{u_2\phi _1}{2\phi _2}& \frac{\phi _1}{\phi _2}\\ \frac{u_2}{2}& 1\end{array}\right),$$
(14)
$$\stackrel{~}{P}^{(0)}=\left(\begin{array}{cc}0& 2\frac{\phi _1}{\phi _2}\\ \frac{u_{2,x}}{2}\frac{u_2^2\phi _1}{2\phi _2}\mu u_2& 0\end{array}\right).$$
(15)
In similar manner one can define the EDT’s of dual problem, which depend explicitly on a solution of the dual ZS problem (11). Successive carrying out of an EDT of direct problem and proper one of dual problem leads to so–called binary DT (BDT) :
$$\stackrel{~}{\mathrm{\Psi }}=\sigma \left(E\frac{\mu \nu }{\lambda \nu }R\right)\mathrm{\Psi },\stackrel{~}{\mathrm{\Xi }}=\frac{\mathrm{\Xi }}{\sigma }\left(E\frac{\mu \nu }{\mu \lambda }R\right),$$
(16)
$$\stackrel{~}{P}^{(0)}=P^{(0)}+(\mu \nu )[P^{(1)},R],$$
(17)
where matrix R is the projector: $`R_{ij}=\phi _i\chi _j/(\chi ,\phi )`$, $`\phi =(\phi _1,\phi _2)^T`$ is the vector solution of Eq.(10) with $`\lambda =\mu `$, $`\chi =(\chi _1,\chi _2)`$ is the vector solution of Eq.(11) with $`\lambda =\nu `$.
It is seen that the potential in the Lax pairs for $`\mathrm{P}_{\mathrm{II}}`$ is the symmetric matrix. The DT technique is convenient for keeping the reduction constraints imposed on the coefficients of the spectral problems . We can obtain the transformed potentials to satisfy this reduction, performing the iterations of BDT .
The DT’s presented in this section allow us to build infinite hierarchy of solutions of nonlinear PDE’s integrable in the frameworks of the IST method. It is supposed in so doing that vectors $`\phi `$ and $`\chi `$ are the solutions of Lax pairs associated with given PDE. Unfortunately, this approach is unfit for obtaining the transformations, which generate infinite hierarchy of solutions of nonlinear ODE’s representable as the compatibility condition of Lax pairs. Vectors $`\phi `$ and $`\chi `$ cannot be regarded as the solutions of systems (1,2), since any transformation of the second equations of these Lax pairs ought to have no the kernel. The basic idea of the generalization of the DT technique for ODE’s is to put the points, in which the DT of the first equations of corresponding Lax pairs has the kernel, into the singular points of the second equations of Lax pairs. For this aim we will assume in the next section that $`\phi `$ and $`\chi `$ are the solutions of systems (1,2) in the points $`\lambda =\mu `$ and $`\lambda =\nu `$. Then we will consider the limits in the formulas of DT’s, tending the points $`\mu `$ and $`\nu `$, in which the kernel exists, to singular point $`\lambda =0`$.
## 4. Transformations of ODE’s
To fulfill the procedure suggested in the previous section we need the explicit expressions of the asymptotic expansions of solutions of Lax pairs (1,2) with coefficients given by Eqs.(5,6) or Eqs.(8,9) at singular point neighborhood.
If $`\alpha `$ is unequal to half–integer, the asymptotic expansion of the Lax pairs solutions of $`\mathrm{P}_{\mathrm{II}}`$ at neighborhood of point $`\lambda =0`$ have the form
$$\mathrm{\Psi }=\mathrm{\Psi }_0(E+\lambda \rho +o(\lambda ))\mathrm{\Lambda },$$
(18)
$$\mathrm{\Xi }=\mathrm{\Lambda }^1(E\lambda \rho +o(\lambda ))\mathrm{\Psi }_0^1.$$
(19)
Here we use notations
$$\mathrm{\Psi }_0=\left(\begin{array}{cc}d& d^1\\ d& d^1\end{array}\right),\mathrm{\Lambda }=\left(\begin{array}{cc}\lambda ^\alpha & 0\\ 0& \lambda ^\alpha \end{array}\right),\rho =\left(\begin{array}{cc}0& \mathrm{\Delta }_1\\ \mathrm{\Delta }_2& 0\end{array}\right),$$
(20)
$$\mathrm{\Delta }_1=\frac{(2u_x+2u^2+x)d^2}{2\alpha 1},\mathrm{\Delta }_2=\frac{(2u_x2u^2x)d^{\mathrm{\hspace{0.17em}2}}}{2\alpha +1},d=\text{exp}\left(u𝑑x\right).$$
The expansions in series of solutions of Lax pairs contain the logarithmic terms in the case of half–integer $`\alpha `$.
The asymptotic expansions of the solutions of Lax pairs of Eq.(7) at neighborhood of point $`\lambda =0`$ are given by Eqs.(18,19), in which the coefficients are defined as follows:
$$\mathrm{\Psi }_0=\left(\begin{array}{cc}c_1& c_2\\ \frac{\alpha z}{wz}c_1& \frac{c_2}{w}\end{array}\right),\mathrm{\Lambda }=\left(\begin{array}{cc}\text{e}^{\left({\displaystyle \frac{\alpha }{2\lambda }}+{\displaystyle \frac{\gamma }{\alpha }}\mathrm{ln}\lambda \right)}& 0\\ 0& \text{e}^{\left({\displaystyle \frac{\alpha }{2\lambda }}{\displaystyle \frac{\gamma }{\alpha }}\mathrm{ln}\lambda \right)}\end{array}\right),$$
(21)
where functions $`c_1`$ and $`c_2`$ solve equations: $`c_{1,x}=\frac{\alpha z}{x}uc_1,c_{2,x}=\frac{z}{x}uc_2`$.
These asymptotic expansions allow us to apply formulas of the DT’s for deriving the transformations of ODE’s. Having the reduction constraint on potential $`P^{(0)}`$, we start from BDT (16,17) in the case of $`\mathrm{P}_{\mathrm{II}}`$. To satisfy condition $`\text{Sp}Q=0`$ under performing the transformation we put $`\sigma =1`$. Substituting
$$\phi =\mathrm{\Psi }|_{\lambda =\mu }\left(\genfrac{}{}{0pt}{}{a_1}{a_2}\right),$$
(22)
$$\chi =(b_1,b_2)\mathrm{\Xi }|_{\lambda =\nu }$$
(23)
($`a_1`$, $`a_2`$, $`b_1`$, $`b_2`$ are constants) and considering the limits $`\mu 0`$ and $`\nu 0`$ in the BDT formulas, taking into account Eqs.(18,19,20), we obtain two transformations of Lax pairs solutions and corresponding well–known transformations of $`\mathrm{P}_{\mathrm{II}}`$ (e.g., see ):
$$\{\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}=\mathrm{\Psi }+\mathrm{ln}_x(\mathrm{\Delta }_1)(\sigma _3+i\sigma _2)\mathrm{\Psi }/(2\lambda )\hfill \\ \stackrel{~}{\mathrm{\Xi }}=\mathrm{\Xi }\mathrm{ln}_x(\mathrm{\Delta }_1)(\sigma _3+i\sigma _2)\mathrm{\Xi }/(2\lambda )\hfill \\ \stackrel{~}{u}=u+\mathrm{ln}_x\mathrm{\Delta }_1\hfill \\ \stackrel{~}{\alpha }=1\alpha \hfill \end{array},\{\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}=\mathrm{\Psi }+\mathrm{ln}_x(\mathrm{\Delta }_2)(\sigma _3i\sigma _2)\mathrm{\Psi }/(2\lambda )\hfill \\ \stackrel{~}{\mathrm{\Xi }}=\mathrm{\Xi }\mathrm{ln}_x(\mathrm{\Delta }_2)(\sigma _3i\sigma _2)\mathrm{\Xi }/(2\lambda )\hfill \\ \stackrel{~}{u}=u\mathrm{ln}_x\mathrm{\Delta }_2\hfill \\ \stackrel{~}{\alpha }=1\alpha \hfill \end{array}.$$
The formulas of transformations coincide with presented above or follow from them as limits in the case of half–integer $`\alpha `$.
The transformations of Eq.(7) are constructed similarly to that of $`\mathrm{P}_{\mathrm{II}}`$. Substituting the asymptotic expansions (18,21) in formulas of EDT’s (12,14) accordingly to Eq.(22) and considering limit $`\mu 0`$, we obtain four transformations of solutions and coefficients of Lax pairs:
$$\{\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}=\left(\begin{array}{cc}\lambda +\frac{u_1}{2}_{}\epsilon _1^\pm & \frac{u_1}{2}\\ \epsilon _1^\pm & 1\end{array}\right)_{}\frac{\mathrm{\Psi }}{\sqrt{\lambda }},\stackrel{~}{\mathrm{\Xi }}=\frac{\mathrm{\Xi }}{\sqrt{\lambda }}\left(\begin{array}{cc}1& \frac{u_1}{2}_{}\\ \epsilon _1^\pm & \lambda +\frac{u_1}{2}\epsilon _1^\pm \end{array}\right),\hfill \\ \stackrel{~}{\beta }=\beta +\frac{1}{2},\stackrel{~}{v}_1=A_1^\pm +\frac{v_1}{2x}_{},\stackrel{~}{v}_2=2x\epsilon _1^\pm ,\stackrel{~}{w}_1=\frac{v_1}{2x}(2q\frac{v_1v_2}{2x}+\hfill \\ +w_1\epsilon _1^\pm +A_1^\pm \epsilon _1^\pm )_{},\stackrel{~}{w}_2=v_22\beta \epsilon _1^\pm v_1\epsilon _1^{\pm \mathrm{\hspace{0.17em}2}},\stackrel{~}{q}=q\frac{v_1v_2}{2x}+A_1^\pm \epsilon _1^\pm ,\hfill \end{array}$$
$$\{\begin{array}{c}\stackrel{~}{\mathrm{\Psi }}=\left(\begin{array}{cc}1& \epsilon _2^\pm \\ \frac{u_2}{2}_{}& \lambda \frac{u_2}{2}\epsilon _2^\pm \end{array}\right)_{}\frac{\mathrm{\Psi }}{\sqrt{\lambda }},\stackrel{~}{\mathrm{\Xi }}=\frac{\mathrm{\Xi }}{\sqrt{\lambda }}\left(\begin{array}{cc}\lambda \frac{u_2}{2}_{}\epsilon _2^\pm & \epsilon _2^\pm \\ \frac{u_2}{2}& 1\end{array}\right),\hfill \\ \stackrel{~}{\beta }=\beta \frac{1}{2},\stackrel{~}{v}_1=2x\epsilon _2^\pm ,\stackrel{~}{v}_2=A_2^\pm \frac{v_2}{2x}_{},\stackrel{~}{w}_1=v_1+2\beta \epsilon _2^\pm v_2\epsilon _2^{\pm \mathrm{\hspace{0.17em}2}},\hfill \\ \stackrel{~}{w}_2=\frac{v_2}{2x}\left(2q+\frac{v_1v_2}{2x}+w_2\epsilon _2^\pm +A_2^\pm \epsilon _2^\pm \right)_{},\stackrel{~}{q}=q\frac{v_1v_2}{2x}A_2^\pm \epsilon _2^\pm ,\hfill \end{array}$$
where $`\epsilon _1^\pm ={\displaystyle \frac{1}{\epsilon _2^\pm }}={\displaystyle \frac{2q\pm \alpha }{2w_1}},A_1^\pm =w_1+{\displaystyle \frac{\beta }{x}}v_1+{\displaystyle \frac{v_1^2}{2x}}\epsilon _1^\pm ,A_2^\pm =w_2+{\displaystyle \frac{\beta }{x}}v_2{\displaystyle \frac{v_2^2}{2x}}\epsilon _2^\pm .`$ (To keep condition $`\text{Sp}Q=0`$ we assume $`\sigma =\lambda ^{1/2}`$ in the formulas of EDT’s.)
These transformations lead to the following ones of solutions and parameters of Eq.(7) $`\{u,\alpha ,\beta ,\gamma \}\stackrel{T_i}{}\{\stackrel{~}{u},\stackrel{~}{\alpha },\stackrel{~}{\beta },\stackrel{~}{\gamma }\}`$ $`(i=1,\mathrm{},4)`$:
$$T_1:\stackrel{~}{u}=\frac{2x(xu_x+\alpha u^22x2(\beta +1)u)}{u(\alpha (xu_x+\alpha u^22x)2\gamma u)},\stackrel{~}{\alpha }=\alpha ,\stackrel{~}{\beta }=\beta +\frac{1}{2},\stackrel{~}{\gamma }=\gamma \frac{\alpha }{2};$$
$$T_2:\stackrel{~}{u}=\frac{2x(xu_x\alpha u^22x2(\beta +1)u)}{u(\alpha (xu_x\alpha u^22x)+2\gamma u)},\stackrel{~}{\alpha }=\alpha ,\stackrel{~}{\beta }=\beta +\frac{1}{2},\stackrel{~}{\gamma }=\gamma +\frac{\alpha }{2};$$
$$T_3:\stackrel{~}{u}=\frac{2x(xu_x+\alpha u^2+2x+2\beta u)}{u(\alpha (xu_x+\alpha u^2+2x)2\gamma u)},\stackrel{~}{\alpha }=\alpha ,\stackrel{~}{\beta }=\beta \frac{1}{2},\stackrel{~}{\gamma }=\gamma \frac{\alpha }{2};$$
$$T_4:\stackrel{~}{u}=\frac{2x(xu_x\alpha u^2+2x+2\beta u)}{u(\alpha (xu_x\alpha u^2+2x)+2\gamma u)},\stackrel{~}{\alpha }=\alpha ,\stackrel{~}{\beta }=\beta \frac{1}{2},\stackrel{~}{\gamma }=\gamma +\frac{\alpha }{2}.$$
The EDT’s of dual pair give the same transformations. The transformations obtained commute and transformations $`T_4`$ and $`T_3`$ are inverse to $`T_1`$ and $`T_2`$ respectively. One can note that Eq.(7) is invariant under changing $`\alpha \alpha `$. This symmetry connects transformations $`T_1`$ and $`T_2`$, $`T_3`$ and $`T_4`$. It should be stressed that a form of these transformations differs from one considered in .
If $`\gamma =\pm \alpha \beta `$ or $`\gamma =\pm \alpha (\beta +1)`$, particular solutions of Eq.(7) satisfy the Riccati equations. In these cases we obtain, performing the transformations of Eq.(7), which keep this condition, two transformations
$$\{\begin{array}{c}\stackrel{~}{\psi }=\sqrt{x}\left(\psi _x\frac{\beta +1}{x}\psi \right)_{}\hfill \\ \stackrel{~}{\beta }=\beta +\frac{1}{2}\hfill \end{array},\{\begin{array}{c}\stackrel{~}{\psi }=\sqrt{x}\left(\psi _x+\frac{\beta }{x}\psi \right)_{}\hfill \\ \stackrel{~}{\beta }=\beta \frac{1}{2}\hfill \end{array}$$
of the linearization of mentioned Riccati equations: $`\psi _{xx}\left({\displaystyle \frac{2\alpha }{x}}+{\displaystyle \frac{\beta (\beta +1)}{x^2}}\right)\psi =0.`$
So, the modification of the DT technique presented here allows us to obtain the transformations of solutions of the nonlinear ODE’s admitting the compatibility condition representation. At the same time we derive the Schlesinger transformations of associated Lax pairs. This approach can be applied for finding the transformations of nonlinear systems of ODE’s and their Lax pairs that have different types of the singularities and the reduction constraints of the coefficients. |
warning/0001/math0001125.html | ar5iv | text | # Topological obstructions to nonnegative curvature
## 1. Introduction
According to the soul theorem of J. Cheeger and D. Gromoll a complete open manifold of nonnegative sectional curvature is diffeomorphic to the total space of the normal bundle of a compact totally geodesic submanifold which is called the soul. One of the harder questions in the subject of is what kind of normal bundles can occur.
Cheeger and Gromoll also proved that a finite cover of any closed nonnegatively curved manifold (throughout the paper by a nonnegatively curved manifold we mean a complete Riemannian manifold of nonnegative sectional curvature) is diffeomorphic to a product of a torus and a simply-connected closed nonnegatively curved manifold. It turns out that a similar statement holds for open complete nonnegatively curved manifolds (see \[Wil98\] and section 2 where a Ricci version of the statement is proved).
We use this fact to find new obstructions to nonnegative curvature and build many examples of vector bundles whose total spaces admit no complete metric of nonnegative curvature. Basic obstructions are provided by the following proposition (of which there is a more general version incorporating the Euler class).
###### Proposition 1.1.
Let $`N`$ be an open complete nonnegatively curved manifold such that $`Q(TN)0`$ for some polynomial $`Q`$ in rational Pontrjagin classes. Then $`Q(T\stackrel{~}{N})0`$ for the universal (and hence any) cover $`\pi :\stackrel{~}{N}N`$.
Note that 1.1 is true for finite covers without any curvature assumptions (because finite covers induce injective maps on rational cohomology). In general, the results of this paper are only interesting for manifold with infinite fundamental groups.
Previously, obstructions to the existence of nonnegatively curved metrics on vector bundles were only known for a flat soul \[ÖW94\]. No obstructions are known when the soul is simply-connected. Examples of nonnegatively curved metrics on vector bundles can be found in \[Che73, Rig78, Yan95, GZ99, GZ\].
###### Corollary 1.2.
Let $`\eta `$ be a vector bundle over a closed smooth manifold $`C`$ and let $`\xi `$ be a vector bundle over a closed flat manifold $`F`$ such that the total space of $`\eta \times \xi `$ admits a complete nonnegatively curved metric. Then $`\xi `$ becomes stably trivial after passing to a finite cover. Furthermore, if either $`rank(\eta )=0`$, or $`\eta `$ is orientable and has nonzero rational Euler class, then $`\xi `$ becomes trivial in a finite cover.
Note that a vector bundle over a flat manifold $`F`$ becomes trivial in a finite cover iff its rational Euler and Pontrjagin classes vanish. Similarly, a bundle over $`F`$ is stably trivial in a finite cover iff its rational Pontrjagin classes vanish (see 4.4).
In case $`C`$ is a point 1.2 says that any vector bundle $`F`$ with nonnegatively curved total space becomes trivial in a finite cover. Also since the Euler and Pontrjagin classes determine a vector bundle up to finite ambiguity (see e.g. \[Bel98\]), in every rank there are only finitely many vector bundles over $`F`$ with nonnegatively curved total spaces. Thus, 1.2 is a generalization of the main result of \[ÖW94\].
To see how 1.2 works, note that if $`T`$ is a torus of dimension $`4`$, then there are infinitely many vector bundles over $`T`$ of every rank $`2`$ with (pairwise) different first Pontrjagin classes. Also there are infinitely many rank $`2`$ vector bundles over $`T^2`$ and $`T^3`$ with different Euler classes. We now deduce the following.
###### Corollary 1.3.
Let $`B`$ be a closed nonnegatively curved manifold. If $`\pi _1(B)`$ contains a free abelian subgroup of rank four (two, respectively), then for each $`k2`$ (for $`k=2`$, respectively) there exists a finite cover of $`B`$ over which there exist infinitely many rank $`k`$ vector bundles whose total spaces admit no nonnegatively curved metrics.
By contrast, any vector bundle over $`S^2\times S^1`$ admits a nonnegatively curved metric as we observe in 7.3. Thus 1.3 cannot be generalized to the case when $`\pi _1(B)`$ is virtually-$``$.
Passing to finite covers in 1.2 and 1.3 seems necessary, in general, in order to obtain bundles without nonnegatively curved metrics. For example, one can easily construct flat $`SO(n)`$ vector bundles over a torus with nonzero Stiefel-Whitney classes, and obviously their total spaces are complete flat manifolds. Here is an example when we get a complete picture without passing to a finite cover.
###### Corollary 1.4.
Let $`\xi `$ be a vector bundle over $`S^3\times S^1`$ whose total space has a nonnegatively curved metric. Then either $`\xi `$ is the trivial bundle or $`\xi `$ is the product of a trivial bundle over $`S^3`$ and the Möbius band bundle over $`S^1`$.
We emphasize that our method does not apply when $`B`$ is simply-connected, or more generally if after passing to a finite cover $`C\times TB`$ the bundle $`\xi `$ becomes a pullback of a bundle over $`C`$ via the projection $`C\times TC`$. (Here, and until the end of the section $`C`$ is a simply-connected manifold and $`T`$ is a torus.) For instance, if $`B`$ is a closed flat manifold which is an odd-dimensional rational homology sphere \[Szc83\], then any vector bundle over $`B`$ becomes trivial in a finite cover and it is unclear whether there are bundles over $`B`$ which are not nonnegatively curved.
A reasonable goal is to find an example of a rank $`k`$ vector bundle over $`C\times T`$ with no nonnegatively curved metric, whenever there is a rank $`k`$ vector bundle over $`C\times T`$ that does not become the pullback of a bundle over $`C`$ in a finite cover. This is achieved in 1.3 when $`dim(T)4`$. Otherwise, the answer may depends on the topology of $`C\times T`$. For example, any bundle of rank $`3`$ over $`2`$-torus becomes trivial, and hence nonnegatively curved, in a finite cover.
While we do not quite settle the case $`dim(T)<4`$, we get various partial results. For instance, given a closed orientable $`2n`$-manifold $`B`$ and an integer $`d0`$, there always exists a map $`f:BS^{2n}`$ of degree $`d`$. Then, if $`\pi _1(B)`$ is infinite, we show that the total space of the pullback bundle $`f^\mathrm{\#}TS^{2n}`$ admits no complete metric of nonnegative curvature. To state further results we need to review some basic bundle theory.
By a simple obstruction-theoretic argument $`H^{even}(C\times T,C)=0`$ implies that any vector bundle over $`C\times T`$ becomes the pullback of a bundle over $`C`$ after passing to a finite cover. This is the case, for example, for bundles over $`CP^n\times S^1`$. However, once $`H^{even}(C\times T,C)0`$ we immediately get a bundle with no nonnegatively curved metric.
###### Corollary 1.5.
If $`H^{2i}(C\times T,C)0`$ for some $`i>0`$, then there exist infinitely many rank $`2i`$ vector bundles over $`C\times T`$ with different Euler classes whose total spaces are not nonnegatively curved.
The Euler class is unstable and, in fact, the bundles constructed in the proof of 1.5 become pullbacks of bundles over $`C`$ after taking Whitney sum with a trivial line bundle and passing to a finite cover.
To get examples that survive stabilization one has to deal with Pontrjagin classes which live in $`H^4(C\times T)`$. Generally, if $`H^4(C\times T,C)=0`$, then after adding a trivial line bundle, any vector bundle over $`C\times T`$ becomes the pullback of a bundle over $`C`$ in a finite cover. If $`H^4(C\times T,C)0`$, one hopes to find a vector bundle without nonnegatively curved metric that survives stabilization and passing to finite covers. We do this in several cases, the simplest being when the rank of the bundle is $`dim(C)`$ (see section 5 for other results involving various assumptions on Pontrjagin classes of $`TC`$).
###### Corollary 1.6.
If $`H^{4i}(C\times T,C)0`$ for some $`i>0`$, then for each $`kdim(C)`$ there exist infinitely many rank $`k`$ vector bundles over $`C\times T`$ with different Pontrjagin classes whose total spaces admit no metric of nonnegative curvature.
The main geometric ingredient of this paper is that a finite cover of any complete nonnegatively curved manifold $`N`$ is diffeomorphic to a product of a torus $`T`$ and a simply connected manifold $`M`$ and this diffeomorphism can be chosen to take a soul $`S`$ to the product of $`T`$ and a simply-connected submanifold of $`M`$. There is also a Ricci version of this statement described in section 2. For example, the above conclusion holds if $`N`$ has nonnegative Ricci curvature, $`S`$ is an isometrically embedded compact submanifold of $`N`$ such that the inclusion $`SN`$ induces an isomorphism of fundamental groups, and either $`S`$ is totally convex, or there exists a distance nonincreasing retraction $`NS`$. In particular, all the theorems stated above hold in these cases.
Our methods also yield obstructions to existence of metrics of nonnegative Ricci curvature on closed manifolds (after all, 1.1 can be applied to closed manifolds). Here is an example. It was shown in \[GW00\] that the total space of the sphere bundle associated with the normal bundle to the soul has a nonnegatively curved metric. Thus, potentially, sphere bundles provide a good source of closed nonnegatively curved manifolds. Among other things, we prove the following.
###### Corollary 1.7.
Let $`\xi `$ be a bundle over a flat manifold $`F`$ with associated sphere bundle $`S(\xi )`$ and let $`C`$ be a closed smooth simply-connected manifold. If $`C\times S(\xi )`$ admits a metric of nonnegative Ricci curvature , then $`\xi `$ becomes trivial in a finite cover.
Finally note that obstructions to to the existence of nonnegatively curved metrics on total spaces of vector bundles give rise to obstructions to the existence of $`G`$-invariant nonnegatively curved metrics on the associated $`G`$-principal bundles. Indeed, any vector bundle $`\xi `$ with a structure group $`G`$ can be written as $`(P\times ^k)/G`$ where $`P`$ is a principal $`G`$-bundle and $`G`$ acts on $`^k`$ via a representation $`GSO(k)`$. By the O’Neill curvature submersion formula, if $`P`$ has a $`G`$-invariant nonnegatively curved metric, then so does the total space of $`\xi `$.
The structure of the paper is as follows. Section 2 contains the above mentioned splitting theorem for nonnegatively curved manifolds. Section 3 summarizes the obstructions to nonnegative curvature coming from the splitting theorem. In section 4 we develop general existence and uniqueness results for bundles over $`C\times T`$. Section 5 contains concrete examples of vector bundles with no nonnegatively curved metrics. Various obstructions to the existence of metrics of nonnegative Ricci curvature on sphere bundles are described in section 6. Theorem 1.4 is proved in the section 7.
We are grateful to William Goldman, Burkhard Wilking, and Wolfgang Ziller for many illuminating conversations. The first author is thankful to the Geometry-Topology group of the McMaster University for support and excellent working conditions.
## 2. Splitting in a finite cover
Cheeger and Gromoll proved in \[CG72\] that a finite cover of a closed nonnegatively curved manifold is diffeomorphic to a product of a torus and a simply connected manifold. The main geometric tool we employ in this paper is the following generalization of this result to open manifolds.
###### Lemma 2.1.
Let $`(N,g)`$ be a complete nonnegatively curved manifold. Then there exists a finite cover $`N^{}`$ of $`N`$ diffeomorphic to a product $`M\times T^k`$ where $`M`$ is a complete open simply connected nonnegatively curved manifold. Moreover, if $`S^{}`$ is a soul of $`N^{}`$, then this diffeomorphism can be chosen in such a way that it takes $`S^{}`$ onto $`C\times T^k`$ where $`C`$ is a soul of $`M`$.
After obtaining this result we have learned that it follows from a more general theorem which was proved earlier by B. Wilking \[Wil98\]. We then realized that our proof of 2.1 in fact gives the following stronger statement.
###### Proposition 2.2.
Let $`(N,g)`$ be a complete manifold of nonnegative Ricci curvature. Let $`q:\stackrel{~}{N}N`$ be the universal cover of $`N`$ and let $`\rho :\pi Iso(\stackrel{~}{N})`$ be the deck transformation representation of $`\pi =\pi _1(N)`$.
Suppose that there exists a closed manifold $`SN`$ isometrically embedded into $`N`$ such that the inclusion $`SN`$ induces an isomorphism of the fundamental groups, and any line in $`\stackrel{~}{S}=q^1(S)`$ with respect to the induced metric from $`\stackrel{~}{N}`$ is also a line in $`\stackrel{~}{N}`$.
Then $`\pi `$ is virtually abelian and, if $`\pi `$ has no torsion, then there exists a smooth path $`\rho (t):[0,1]\mathrm{Hom}(\pi ,Iso(\stackrel{~}{N}))`$ such that
1. $`\rho (0)=\rho `$;
2. for each $`t`$ the action of $`\pi `$ on $`\stackrel{~}{N}`$ is free and properly discontinuous;
3. A finite cover of $`N_1=\stackrel{~}{N}/\rho (1)(\pi )`$ splits isometrically as $`M\times T^k`$ where $`k=rank(\pi )`$;
4. There is a family of closed submanifolds $`S_tN_t=\stackrel{~}{N}/\rho (t)(\pi )`$ such that
1. $`S_0=S`$
2. Under the splitting from (iii) the cover of $`S_1`$ corresponds to the Riemannian product $`C\times T^kM\times T^k`$ where $`C`$ is a closed isometrically embedded submanifold of $`M`$.
3. for each $`t`$ there exists a diffeomorphism $`\varphi _t:(N_t,S_t)(N_0,S_0)`$
The assumption that any line in $`\stackrel{~}{S}=q^1(S)`$ is also a line in $`\stackrel{~}{N}`$ is satisfied if $`S`$ is totally convex in $`N`$ or if there is a distance nonincreasing retraction $`NS`$. Both of these conditions are true if $`N`$ is an open manifold of nonnegative sectional curvature and $`SN`$ is its soul. In this case one can also describe the souls of the deformed manifolds $`N_t`$. Namely we have the following
###### Proposition 2.3.
Let $`(N,g)`$ be a complete nonnegatively curved manifold with a free abelian fundamental group $`\pi `$. Let $`\rho :\pi Iso(\stackrel{~}{N})`$ be the deck transformation representation of $`\pi `$. Then there exists a smooth path $`\rho (t):[0,1]\mathrm{Hom}(\pi ,Iso(\stackrel{~}{N}))`$ such that in addition to (i)-(iv) of 2.2 the following holds.
(v) If $`S`$ is a soul of $`N`$, then there exists an isometric splitting $`\stackrel{~}{N}=M\times ^k`$ where $`k=rank(\pi )`$ and a soul $`C`$ of $`M`$ such that, for every $`t[0,1]`$, the projection $`S_t`$ of $`C\times ^k`$ to $`N_t=\stackrel{~}{N}/\rho (t)(\pi )`$ is a soul of $`N_t`$. Also for each $`t`$, there exists a diffeomorphism $`\varphi _t:(N_t,S_t)(N_0,S_0)`$.
###### Remark 2.4.
The above mentioned result of Wilking \[Wil98\] implies the existence of the deformation $`\rho (t)`$ as in 2.2 for an arbitrary virtually abelian group. He also gives an upper bound on the order of the covering in question in terms of $`\pi `$ and the number of connected components of $`Iso(M)`$. Nevertheless, we will present our proof of 2.2 for it is considerably easier than the one in \[Wil98\]. (In fact, our proof is very similar to the original argument of Cheeger and Gromoll in the closed manifold case.) Besides, the statements of 2.2 and 2.3 are tailored to our applications, for example the parts $`(iv)(v)`$ are not discussed in \[Wil98\].
###### Proof of 2.2.
Let $`q:\stackrel{~}{N}N`$ be the universal cover of $`N`$ and let $`\stackrel{~}{S}=q^1(S)`$. Then since inclusion $`SN`$ induces an isomorphism of the fundamental groups $`q|_{\stackrel{~}{S}}:\stackrel{~}{S}S`$ is the universal cover of $`S`$.
Let $`\stackrel{~}{S}=C\times ^k`$ be the de Rham decomposition of $`\stackrel{~}{S}`$ so that $`C`$ does not split off a Euclidean factor. We claim that $`C`$ is compact. (Indeed, suppose $`C`$ is not compact. Then $`C`$ contains a ray $`\gamma `$. Since $`S`$ is compact, there exists a point $`pS`$ such that $`q(\gamma (i))p`$ as $`i\mathrm{}`$. Let $`\stackrel{~}{p}`$ be a point in $`q^1(p)`$. By above there exists a sequence $`g_i\pi `$ such that $`g_i(\gamma (i))\stackrel{~}{p}`$. Passing to a subsequence we can assume that $`g_i(\gamma ^{}(i))vT_{\stackrel{~}{p}}C`$. Just as in \[CG72\] we readily conclude that $`\sigma (t)=\mathrm{exp}(tv):\stackrel{~}{S}`$ is a line in $`C\stackrel{~}{S}`$. By assumption $`\sigma (t)`$ is also a line in $`\stackrel{~}{N}`$. Therefore, the splitting theorem \[CG72\] implies $`\stackrel{~}{N}`$ splits off $`\sigma (t)`$ isometrically. So $`v`$ is invariant under the $`Hol(\stackrel{~}{N})`$ and hence under $`Hol(\stackrel{~}{S})`$ which contradicts the fact that $`C`$ does not split off a Euclidean factor.)
Since any isometry of $`\stackrel{~}{S}`$ takes lines to lines, the isometry group $`Iso(\stackrel{~}{S})`$ splits as a direct product $`Iso(\stackrel{~}{S})Iso(C)\times Iso(^k)`$. Therefore, the natural deck transformation action of $`\pi `$ on $`\stackrel{~}{S}`$ gives a monomorphism $`\rho =(\rho _1,\rho _2):\pi Iso(C)\times Iso(^k)`$.
Since $`Iso(C)`$ is compact and $`\pi `$ is discrete, the group $`\rho _2(\pi )`$ is a discrete subgroup of $`Iso(^k)`$. Also $`\mathrm{ker}(\rho _2)`$ is compact and hence it is finite. Thus, $`\pi `$ is an extension of a finite group by a crystallographic one. It is well-known (see e.g. the proof of \[Wil98, Thm 2.1\] ) that any such a group is virtually abelian.
Now suppose that $`\pi `$ is free abelian. Then $`\rho _2(\pi )`$ is a discrete torsion-free subgroup of $`Iso(^k)`$, in particular, it acts on $`^k`$ by translations and $`^k/\rho _2(\pi )`$ is isometric to a flat torus $`T^k`$.
By above the splitting $`\stackrel{~}{S}=C\times ^k`$ is just a part of a bigger isometric splitting $`\stackrel{~}{N}=\stackrel{~}{M}\times ^k`$ where $`\stackrel{~}{M}`$ is a complete open simply connected manifold containing $`C`$ as an isometrically embedded submanifold.
Since the action of $`\pi `$ on $`\stackrel{~}{N}`$ leaves $`\stackrel{~}{S}`$ invariant, it sends lines parallel to $`^k`$ into lines parallel to $`^k`$. Hence the map $`\rho `$ is a restriction of a natural monomorphism $`\pi Iso(\stackrel{~}{M})\times Iso(^k)`$ which with a slight abuse of notations we will still denote by $`\rho =(\rho _1,\rho _2)`$. In fact, the image of $`\rho _1`$ lies in the subgroup $`GIso(\stackrel{~}{M})`$ of isometries leaving $`C`$ invariant. Since $`C`$ is compact it follows that $`G`$ is a compact subgroup of $`Iso(\stackrel{~}{M})`$.
Next consider the homomorphism $`\rho _1:\pi G`$. Let $`H`$ be the closure of $`\rho _1(\pi )`$ in $`G`$. Then $`H`$ is a compact abelian subgroup of $`G`$. Let $`H_0`$ be the identity component of $`H`$. Consider the short exact sequence $`1H_0H\mathrm{\Gamma }1`$ where $`\mathrm{\Gamma }=H/H_0`$ is a finite abelian group. We claim that this sequence splits and hence $`HH_0\times \mathrm{\Gamma }`$.
Indeed, the group $`\mathrm{\Gamma }`$ is a product of finite cyclic subgroups and, since $`H`$ is abelian, it is enough to define the splitting on generators of these subgroups. Let $`g\mathrm{\Gamma }`$ be a generator of order $`m`$ and let $`\overline{g}H`$ be a preimage of $`g`$. The endomorphism of $`H`$ sending $`x`$ to $`x^m`$, takes $`\overline{g}`$ to $`H_0`$, and maps $`H_0`$ onto itself. Hence, there is $`hH_0`$ such that $`h^m=\overline{g}^m`$, and we can define a splitting by mapping $`g`$ to $`\overline{g}h^1`$.
Thus, $`\rho _1:\pi HH_0\times \mathrm{\Gamma }`$ can be written as a product of two representations $`\rho ^{}:\pi H_0`$ and $`\rho ^{\prime \prime }:\pi \mathrm{\Gamma }`$. Since $`H_0T^l`$, the representation variety $`\mathrm{Hom}(\pi ,H_0)`$ is diffeomorphic to a torus $`T^{kl}`$. Hence, we can find a smooth deformation $`\rho _1^{}(t)\mathrm{Hom}(\pi ,H_0)`$ such that $`\rho _1^{}(0)=\rho ^{}`$ and $`\rho _1^{}(1)=1`$, the trivial representation.
Crossing $`\rho _1^{}(t)`$ with $`\rho ^{\prime \prime }`$ and $`\rho _2`$, we obtain a smooth path $`\rho (t)\mathrm{Hom}(\pi ,Iso(\stackrel{~}{N}))`$ such that $`\rho (0)=\rho `$ and $`\rho (1)=1\times \rho ^{\prime \prime }\times \rho _2`$. For every $`t`$ the action of $`\pi `$ on $`\stackrel{~}{N}`$ via $`\rho _t`$ is free and properly discontinuous because so is the action of $`\pi `$ on $`^k`$. Therefore, we get a smooth family of manifolds of nonnegative Ricci curvature $`N_t=\stackrel{~}{N}/\rho (t)(\pi )`$ with $`N_0=N`$. We also get the family $`S_t=\stackrel{~}{S}/\rho (t)(\pi )N_t`$ of closed isometrically embedded submanifolds with $`S_0=S`$.
The finite cover of $`N_1`$ corresponding to the kernel of $`\rho ^{\prime \prime }`$ splits isometrically as $`M\times T^k`$. Under the splitting the cover of $`S_1`$ corresponds to the Riemannian product $`C\times T^kM\times T^k`$.
By the (relative) covering homotopy theorem the family $`(N_t,S_t)`$, considered as a bundle over $`[0,1]`$ is smoothly isomorphic to the trivial bundle $`[0,1]\times (N,S)`$. In particular, all $`S_t`$’s are mutually diffeomorphic and, moreover, have isomorphic normal bundles. ∎
###### Proof of 2.3.
Let $`S`$ be a soul of $`N`$ and let $`p:\stackrel{~}{N}N`$ be the universal cover of $`N`$. By the Cheeger-Gromoll soul theorem \[CG72\] $`S`$ is totally convex and the inclusion $`SN`$ is a homotopy equivalence. Thus, 2.2 applies and it only remains to deduce (v).
Let $`h:N`$ be the Cheeger-Gromoll exhaustion function generating $`S`$ and let $`\stackrel{~}{h}=hq:\stackrel{~}{N}`$ be its lift to the universal cover $`\stackrel{~}{N}`$. Clearly, $`\stackrel{~}{h}`$ is convex. Moreover, since every line in $`\stackrel{~}{N}`$ parallel to $`^k`$ projects to an infinite geodesic lying in a compact set, $`\stackrel{~}{h}`$ is constant along any such line. Hence $`\stackrel{~}{h}`$ is given by the formula $`\stackrel{~}{h}(m,t)=\overline{h}(m)`$ for some convex function $`\overline{h}:M`$. It is easy to see that $`\overline{h}`$ is an exhaustion function. Let $`CM`$ is the soul generated by $`\overline{h}`$.
By construction $`\overline{h}`$ is invariant under the action of $`\rho _1(\pi )`$ and, hence, under the action of $`H`$. In particular, $`\overline{h}`$ is invariant under the action of $`\rho _1(t)(\pi )`$ for any $`t`$. Therefore, $`\overline{h}`$ descends to a well defined convex exhaustion function $`h_t:N_t`$ generating the soul $`S_t=(C\times ^k)/\rho (t)(\pi )`$. ∎
###### Remark 2.5.
Actually, it follows from the proof of 2.2 that some versions of 2.2 and 2.1 hold without any curvature assumptions. For example, let $`N`$ be a complete Riemannian manifold whose universal cover is isometric to $`M\times ^n`$ where $`Iso(M)`$ is compact. Then a finite cover of $`N`$ is diffeomorphic to the product $`M\times T^k\times ^{nk}`$. See \[Wil98\] for a stronger result.
## 3. Basic obstructions
In this section we obtain simple topological obstructions to nonnegative curvature coming from the results of the section 2.
In this section and throughout the rest of this paper we use the notation $`e`$ for the Euler class, $`p_i`$ for the $`i`$th Pontrjagin class, and $`p=_{i0}p_i`$ for total Pontrjagin class. Unless stated otherwise, all the characteristic classes live in cohomology with rational coefficients. (However, it is useful to keep in mind that $`e`$ and $`p_i`$ are in fact integral classes, that is they lie in the image of $`H^{}(B,)H^{}(B,)`$.)
Let $`S`$ be a closed manifold smoothly embedded into an open manifold $`N`$ such that the inclusion $`SN`$ induces an isomorphism of fundamental groups. Let $`q:\stackrel{~}{N}N`$ be the universal cover of $`N`$; then $`q:\stackrel{~}{S}=q^1(S)S`$ is the universal cover of $`S`$. Assume that after passing to a finite cover $`N`$ becomes diffeomorphic to $`M\times T`$ where $`\pi _1(M)=1`$ and $`T`$ is a torus of positive dimension, Further, suppose that this diffeomorphism takes (a finite cover of) $`S`$ onto $`C\times T`$ where $`C`$ is a submanifold of $`M`$. Denote the normal bundles of $`S`$ in $`N`$ by $`\nu _S`$.
###### Lemma 3.1.
Suppose there is a polynomial $`Q`$ with rational coefficients such that $`Q(e(\nu _S),p_1(TN|_S),p_2(TN|_S),\mathrm{})0`$ where $`\nu _S`$ is assumed to be oriented if $`Q`$ depends on $`e`$. Then $`Q(e(q^\mathrm{\#}\nu _S),p_1(T\stackrel{~}{N}|_{\stackrel{~}{S}}),p_2(T\stackrel{~}{N}|_{\stackrel{~}{S}}),\mathrm{})0`$.
###### Proof.
Note that $`Q_S=Q(e(\nu _S),p_1(TN|_S),p_2(TN|_S),\mathrm{})0`$ remains true after passing to any finite cover because finite covers induce injective maps on rational cohomology. Thus, we can assume without loss of generality that $`N`$ is diffeomorphic to $`M\times T`$ as above and this diffeomorphism identifies $`S`$ with $`C\times T`$.
Then the normal bundle $`\nu _C^M`$ of $`C`$ in $`M`$ is the pullback of $`\nu _S`$ via the inclusion $`i_C:CS`$ and, also $`\nu _S`$ is the pullback of $`\nu _C^M`$ via the projection $`\pi _C:SC`$. Similarly, since $`T`$ is parallelizable, $`TN|_C`$ is stably isomorphic to $`i_C^\mathrm{\#}TN|_S`$ and $`TN|_S`$ is stably isomorphic to $`\pi _C^\mathrm{\#}TN|_C`$. In particular,
$$Q_C=Q(e(\nu _C^M),p_1(TN|_C),p_2(TN|_C),\mathrm{})=i_C^{}Q_S$$
and $`\pi _C^{}Q_C=Q_S`$. The latter implies that $`Q_C0`$.
Since $`C`$ is simply-connected, $`i_C`$ factors through the universal covering $`q:\stackrel{~}{S}S`$. In particular, $`q^{}Q_S0`$ as desired. ∎
###### Remark 3.2.
Clearly, 3.1 remains true for any (not necessarily universal) cover $`q^{}`$ in place of $`q`$ because $`q^{}Q_S=0`$ implies $`q^{}Q_S=0`$.
###### Remark 3.3.
A particular case of 3.1 remains true even without mentioning $`S`$. Namely, assume only that $`N`$ is an open manifold whose finite cover is diffeomorphic to $`M\times T`$. Let $`Q`$ be a polynomial in rational Pontrjagin classes such that $`Q(TN)0`$. Then the same proof implies $`Q(T\stackrel{~}{N})0`$. This applies to the geometric situation discussed in 2.5.
### Base versus soul.
Now we specialize to the case when $`N`$ is the total space of a smooth vector bundle over a closed manifold $`B`$. We identify $`B`$ with the zero section. Then the universal cover of $`B`$ is $`q:\stackrel{~}{B}=q^1(B)B`$. Assume also that the inclusion $`SN`$ is a homotopy equivalence.
First, we need to see how the characteristic classes of the normal bundles to the $`S`$ and $`B`$ are related. The homotopy equivalence $`h:BS`$ (defined as the composition of the inclusion $`BN`$ and a homotopy inverse of $`SN`$) clearly has the property that $`h^{}p_i(TN|_S)=p_i(TN|_B)`$ for any $`i`$.
Furthermore, if the manifolds $`N`$, $`B`$, $`S`$ are oriented, then for the rational Euler class we have $`h^{}e(\nu _S)=\mathrm{deg}(h)e(\nu _B)`$. (Indeed, suppressing the inclusions we have $`h^{}e(\nu _B),\alpha =e(\nu _B),h_{}\alpha =[B],h_{}\alpha =\mathrm{deg}(h)[S],\alpha =\mathrm{deg}(h)e(\nu _S),\alpha `$.)
By possibly changing orientation on $`S`$ we can arrange that $`\mathrm{deg}(h)=1`$ so that $`h^{}e(\nu _B)=e(\nu _S)`$. Thus, since $`h^{}`$ is an algebra homomorphism, we get
$$h^{}Q(e(\nu _B),p_1(TN|_B),p_2(TN|_B),\mathrm{})=Q(e(\nu _S),p_1(TN|_S),p_2(TN|_S),\mathrm{}).$$
###### Proposition 3.4.
Let $`\xi `$ be a vector bundle over a closed smooth manifold $`B`$ whose total space $`N`$ admits a complete Riemannian metric of nonnegative sectional curvature. Suppose there is a polynomial $`Q`$ with rational coefficients such that $`Q(e(\xi ),p_1(TN|_B),p_2(TN|_B),\mathrm{})0`$ where $`\xi `$ is assumed to be oriented if $`Q`$ depends on $`e`$. Then $`Q(e(q^\mathrm{\#}\xi ),p_1(T\stackrel{~}{N}|_{\stackrel{~}{B}}),p_2(T\stackrel{~}{N}|_{\stackrel{~}{B}}),\mathrm{})0`$.
###### Proof.
Let $`S`$ be a soul of $`N`$. By above the homotopy equivalence $`h`$ takes $`Q(e(\xi ),p_1(TN|_B),p_2(TN|_B),\mathrm{})`$ to $`Q(e(\nu _S),p_1(TN|_S),p_2(TN|_S),\mathrm{})`$ hence the latter is nonzero.
By 3.1 we have $`Q(e(q^\mathrm{\#}\nu _S),p_1(T\stackrel{~}{N}|_{\stackrel{~}{S}}),p_2(T\stackrel{~}{N}|_{\stackrel{~}{S}}),\mathrm{})0`$. Let $`\stackrel{~}{h}`$ be the lift of $`h`$ to the universal covers; note that $`\stackrel{~}{h}`$ is a homotopy equivalence. Then by commutativity $`\stackrel{~}{h}^{}e(q^\mathrm{\#}\nu _B)=e(q^\mathrm{\#}\nu _S)`$ and $`\stackrel{~}{h}^{}p_i(T\stackrel{~}{N}|_{\stackrel{~}{B}})=p_i(T\stackrel{~}{N}|_{\stackrel{~}{S}})`$. So the homotopy inverse of $`\stackrel{~}{h}`$ takes $`Q(e(q^\mathrm{\#}\nu _S),p_1(T\stackrel{~}{N}|_{\stackrel{~}{S}}),p_2(T\stackrel{~}{N}|_{\stackrel{~}{S}}),\mathrm{})`$ to the corresponding polynomial for $`\stackrel{~}{B}`$ which is therefore nonzero as claimed. ∎
###### Remark 3.5.
The statement of 3.4 becomes especially simple if $`p(TB)=1`$. Indeed, it implies that $`p(TN|_B)=p(\xi TB)=p(\xi )p(TB)=p(\xi )`$. We also get $`p(T\stackrel{~}{B})=1`$ which implies $`p(T\stackrel{~}{N}|_{\stackrel{~}{B}})=p(q^\mathrm{\#}\xi )`$.
We shall often use the following variation of 3.4.
###### Proposition 3.6.
Let $`\xi `$ be an vector bundle over $`B=C\times T`$ where $`C`$ is a closed connected smooth manifold and $`T`$ is a torus. Assume the total space $`N`$ of $`\xi `$ admits a complete Riemannian metric of nonnegative sectional curvature. Suppose there is a polynomial $`Q`$ with rational coefficients such that $`Q(e(\xi ),p_1(TN|_B),p_2(TN|_B),\mathrm{})0`$ where $`\xi `$ is assumed to be oriented if $`Q`$ depends on $`e`$. Then $`Q(e(i_C^\mathrm{\#}\xi ),p_1(TN|_C),p_2(TN|_C),\mathrm{})0`$ where $`i_C:CC\times T`$ is the inclusion.
###### Proof.
The universal cover $`q:\stackrel{~}{B}=\stackrel{~}{C}\times ^kC\times T=B`$ clearly factors through the inclusion $`i_C:CC\times T`$. By 3.4, $`Q(e(q^\mathrm{\#}\xi ),p_1(T\stackrel{~}{N}|_{\stackrel{~}{B}}),p_2(T\stackrel{~}{N}|_{\stackrel{~}{B}}),\mathrm{})0`$, therefore, $`Q(e(i_C^\mathrm{\#}\xi ),p_1(TN|_C),p_2(TN|_C),\mathrm{})`$ must be nonzero. ∎
## 4. Producing vector bundles
In this section we discuss some methods of building vector bundles. We start from several general methods and then concentrate on the case when the base is $`C\times T`$ where $`C`$ is a finite connected CW-complex and $`T`$ is a torus of positive dimension.
###### Example 4.1.
Let $`B`$ be a closed orientable $`2n`$-manifold and let $`\xi `$ be a bundle over $`S^{2n}`$. Since there always exists a degree one map $`f:BS^{2n}`$, we get a pullback bundle $`f^\mathrm{\#}\xi `$. Now if $`\xi `$ has a nonzero rational characteristic class (that necessarily lives in $`H^{2n}(S^{2n},)`$), so does $`f^\mathrm{\#}\xi `$ because $`f`$ induces an isomorphism on the $`2n`$-dimensional cohomology. In particular, every even integer $`2d`$ can be realized as the Euler number of a rank $`2n`$ bundle over $`B`$ (by taking $`\xi `$ to be the pullback of $`TS^{2n}`$ via a self-map of $`S^{2n}`$ of degree $`d`$).
###### Example 4.2.
Any element of $`H^2(B,)`$ can be realized as the Euler class of an oriented rank two bundle over $`B`$ (where $`B`$ is any paracompact space) \[Hir66, I.4.3.1\].
###### Example 4.3.
If $`B`$ is a finite CW-complex of dimension $`d`$, then it is well-known that a multiple of any element of $`_{i>0}H^{4i}(B,)`$ can be realized as the Pontrjagin character of a vector bundle over $`B`$ of rank $`d`$ (and hence of any rank $`d`$).
In particular, a multiple of any element $`xH^{4k}(B,)`$ can be realized as the $`k`$th Pontrjagin class of a bundle of rank $`d`$. (Indeed, let $`X`$ be the image of $`x`$ under the inclusion $`H^{4k}(B,)_{i>0}H^{4i}(B,)`$. Realize a multiple of $`X`$ as the Pontrjagin character of a bundle. Then this bundle has zero Pontrjagin classes $`p_i`$ for $`0<i<k`$ and the $`k`$th Pontrjagin class is a multiple of $`x`$.)
We now prove a uniqueness and existence theorem for vector bundles over $`C\times T`$.
###### Theorem 4.4.
Let $`C`$ be a finite connected CW-complex and let $`\xi `$ and $`\eta `$ be oriented rank $`n`$ vector bundles over $`C\times T`$ such that
(1)
$$\xi |_{C\times }\eta |_{C\times }$$
and
(2)
$$\xi \text{ and }\eta \text{ have the same rational characteristic classes.}$$
Then there exists a finite cover $`\pi :C\times TC\times T`$ such that $`\pi ^\mathrm{\#}\xi \pi ^\mathrm{\#}\eta `$.
###### Proof.
First, note that after passing to a finite cover we can assume that $`\xi `$ and $`\eta `$ have the same integral cohomology classes. (Indeed, look for example at the integral $`i`$-th Pontrjagin class $`p_i`$. By assumption $`p_i(\xi )p_i(\eta )`$ is a torsion element of $`H^{4i}(C\times T,)`$. By the Künneth formula we can write $`p_i(\xi )p_i(\eta )=_{j=0}^{4i}c^{4ij}t^j`$ where $`c^sH^s(C,)`$ and $`t^jH^j(T^k,)`$. Condition (1) implies $`t^0c^{4i}=p_i(\xi |_{C\times })p_i(\eta |_{C\times })=0`$. Since any torsion element of the form $`_{j=1}^{4i}c^{4ij}t^j`$ becomes zero when mapped to an appropriate finite cover $`C\times TC\times T`$ along $`T`$, the integral Pontrjagin classes $`p_i(\xi )`$, $`p_i(\eta )`$ become equal in such a finite cover.)
Let $`f,g:C\times TBSO(n)`$ be the classifying maps for $`\xi `$ and $`\eta `$ respectively. Let $`\gamma ^n`$ be the universal bundle over $`BSO(n)`$. For each $`i[n/2]`$, we view the classes $`p_i(\gamma ^n)H^{4i}(BSO(n))`$ as maps $`p_i:BSO(n)K(,4i)`$, and similarly if $`n`$ is even, $`e(\gamma ^n)H^n(BSO(n))`$ is thought of as a map $`e:BSO(n)K(,n)`$.
Consider the combined map $`c`$ from $`BSO(n)`$ to the product of Eilenberg-MacLane spaces given by the formula
$$c=(p_1,p_2\mathrm{}p_{(n1)/2}):BSO(n)X=\times _{s=1}^{(n1)/2}K(,4s)\text{ if }n\text{ is odd}$$
and,
$$c=(e,p_1,p_2\mathrm{},p_{n/21}):BSO(n)X=K(,n)\times (\times _{s=1}^{n/21}K(,4s)).$$
if $`n`$ is even.
It is well known that $`c`$ is a rational homotopy equivalence (i.e. the homotopy fiber $`F`$ of $`c`$ has finite homotopy groups) and the spaces $`F,BSO(n),X`$ are simply-connected (see e.g. \[Bel98\]).
By the condition (2) $`pf`$ is homotopic to $`pg`$. We shall now try to lift this homotopy to the homotopy of $`f`$ and $`g`$. We view the pair $`(C\times T,C)`$ (where we identify $`C\times `$ with $`C`$) as a relative CW-complex and we try to construct the homotopy between $`f`$ and $`g`$ inductively on the dimension of the skeletons. By (1) we can assume that the homotopy is already constructed on the zero skeleton $`(C\times T,C)_0`$.
Suppose that we have already constructed the homotopy on $`(C\times T,C)_{i1}`$ for $`i>0`$. We want to show that after possibly passing to a finite cover we can extend it over $`(C\times T,C)_i`$. The relative obstruction $`O_i`$ to the extension over the $`i`$-th skeleton lives in the cohomology $`H^i((C\times T,C),\pi _i(F))`$. Let $`m=|\pi _i(F)|`$ and $`k=dim(T)`$; by assumption $`k>0`$. Consider the $`m^k`$ cover $`\mathrm{\Pi }:C\times TC\times T`$ given by the formula $`(c,z_1,\mathrm{},z_k)(c,z_1^m,\mathrm{},z_k^m)`$. Notice that by the Künneth formula for the pair $`(C\times T,C)=(C,\mathrm{})\times (T,)`$ we have
$$H^i((C\times T,C),\pi _i(F))=\underset{j=0}{\overset{i}{}}H^{ij}((C,\mathrm{}),\pi _i(F))H^j((T,),\pi _i(F))=$$
$$=\underset{j=1}{\overset{i}{}}H^{ij}((C,\mathrm{}),\pi _i(F))H^j((T,),\pi _i(F))$$
where the last equality is due to the fact that $`H^0((T,),\pi _i(F))=0`$. Therefore, for any $`\delta H^i((C\times T,C),\pi _i(F))`$, its pullback $`\mathrm{\Pi }^{}(\delta )`$ is an $`m`$-th multiple of some class, and thus is equal to zero. In particular, $`\mathrm{\Pi }^{}(O_i)=0`$. On the other hand, by the naturality of obstructions $`\mathrm{\Pi }^{}(O_i)`$ is the obstruction to extending the homotopy between $`f\mathrm{\Pi }`$ and $`g\mathrm{\Pi }`$ over the relative $`i`$-skeleton $`(C\times T,C)_i`$. ∎
###### Theorem 4.5.
Let $`\xi `$ be an oriented rank $`n`$ vector bundle over a finite connected CW-complex $`C`$ and let $`i:CC\times T`$ be the canonical inclusion onto $`C\times `$. Let $`e^{}H^n(C\times T),p_1^{}H^4(C\times T),\mathrm{},p_{[n/2]}^{}H^{4[n/2]}(C\times T)`$ be a collection of integral cohomology classes such that their restrictions onto $`C\times `$ give corresponding integral characteristic classes of $`\xi `$ and, furthermore, $`e^{}=0`$ if $`n`$ is odd and $`p_{[n/2]}^{}=e^{}e^{}`$ if $`n`$ even. Then there exists a finite cover $`\mathrm{\Pi }:C\times TC\times T`$ and a rank $`n`$ vector bundle $`\eta `$ over $`C\times T`$ such that the integral characteristic classes of $`\eta `$ satisfy $`e(\eta )=\mathrm{\Pi }^{}(e^{})`$ and $`p_i(\eta )=\mathrm{\Pi }^{}(p_i^{})`$ for $`i=1,\mathrm{},[n/2]`$.
###### Proof.
Again, consider the universal fibration $`FBSO(n)\stackrel{c}{}X`$ where $`X`$ is the product of appropriate Eilenberg-MacLane spaces.
The collection of characteristic classes $`(e^{},p_1^{},p_2^{},\mathrm{},p_{[n/2]}^{})`$ defines a natural map $`c^{}:C\times TX`$ (where we exclude $`e^{}`$ for odd $`n`$ and $`p_{[n/2]}^{}`$ for even $`n`$). It suffices to show that after passing to a finite cover there exists a lift of this map to the map $`f:C\times TBSO(n)`$.
By assumptions, we can construct the lift $`f`$ over $`C\times `$ by letting $`f|_{C\times }`$ to be equal to a classifying map of $`\xi `$.
Now we are faced with the relative lifting problem of extending the lift from $`C\times `$ to $`C\times T`$. As in the proof of 4.4 we will proceed by induction on the dimension of the relative skeleton $`(C\times T,C\times )_i`$.
Suppose $`f`$ is already defined on $`(C\times T,C\times )_{i1}`$ for some $`i>0`$. As before the primary obstruction $`O_i`$ to extending the lift over $`(C\times T,C\times )_i`$ lives in the group $`H^i((C\times T,C),\pi _{i1}(F))`$. Arguing exactly as in the proof of 4.4, we see that the cover $`\mathrm{\Pi }:C\times TC\times T`$ given by the formula $`(c,z_1,\mathrm{},z_k)(c,z_1^m,\mathrm{},z_k^m)`$ where $`m=|\pi _{i1}(F)|`$ has the property that $`\mathrm{\Pi }^{}(O_i)=0`$. Therefore, the lift of $`c^{}\mathrm{\Pi }`$ given by $`f\mathrm{\Pi }`$ can be extended over the relative $`i`$-th skeleton of $`(C\times T,C\times )`$. This completes the proof of the induction step and hence the proof of the theorem. ∎
###### Remark 4.6.
Note that by construction the cover $`\mathrm{\Pi }`$ depends only on $`n`$ and $`dim(C\times T)`$.
###### Remark 4.7.
The proof of 4.5 shows how to compute the characteristic classes of a bundle with the classifying map $`f`$. For example, represent $`e^{}`$ as $`_{j=0}^ke_j^{}`$ where $`e_j^{}H^{nj}(C,)H^j((T,),)`$. Then the Euler class of $`f`$ is given by $`_{j=0}^km^je_j^{}`$ where $`dim(T)=k`$. The same result is of course true for any Pontrjagin class of $`f`$.
In particular, if $`e^{}=e_j^{}`$ for some $`j>0`$, then an integer multiple of $`e^{}`$ is realized as the Euler class of some bundle over $`C\times T`$.
###### Example 4.8.
We shall often use 4.5 in the following situation. Assume $`H^{4i}(C\times T,C,)`$ is infinite and let $`p_i^{}H^{4i}((C\times T,C),)`$ be a nontorsion class and $`j`$ be any nonzero integer. Let $`\xi `$ be the trivial bundle of some rank $`>2i`$. Then by 4.5 there exists a bundle $`\eta _j`$ over $`C\times T`$ and a finite cover $`\mathrm{\Pi }:C\times TC\times T`$ such that the restriction of $`\eta _j`$ to $`C\times `$ is isomorphic to $`\xi `$ and $`\mathrm{\Pi }^{}(jp_i^{})=p_i(\eta _j)`$. Clearly, the bundles $`\eta _j`$ are pairwise nonisomorphic.
## 5. Vector bundles with no nonnegatively curved metrics
In this section we obtain concrete examples of bundles without nonnegatively curved metrics. Throughout this section $`T`$ is a torus of positive dimension and $`C`$ is a closed connected smooth manifold. (Note that $`C`$ is not assumed to be simply-connected so the results of this section are slightly more general than the ones stated in the introduction.) We shall often use that the tangent bundle to $`C\times T`$ is stably isomorphic to the pullback of $`TC`$ via the projection $`C\times TC`$. All (co)homology groups and characteristic classes in this section have rational coefficients.
###### Corollary 5.1.
Let $`\eta `$ be a vector bundle over $`C`$ and let $`\xi `$ be a vector bundle over $`T`$ such that the total space of $`\eta \times \xi `$ admits a complete nonnegatively curved metric. Then $`\xi `$ becomes stably trivial in a finite cover. Furthermore, if either $`rank(\eta )=0`$, or $`\eta `$ is orientable with $`e(\eta )0`$, then $`\xi `$ becomes trivial in a finite cover.
###### Proof.
Denote $`\eta TC`$ by $`\eta ^{}`$ so that the tangent bundle to the total space of $`\eta \times \xi `$ restricted to the zero section is stably isomorphic to $`\eta ^{}\times \xi `$. Let $`i`$ be the largest nonnegative integer such that $`p_i(\eta ^{})0`$. Arguing by contradiction assume that $`p_k(\xi )0`$ for some $`k>0`$. Using the product formula $`p(\eta ^{}\times \xi )=p(\eta ^{})\times p(\xi )`$, we conclude that the component of $`p_{i+k}(\eta ^{}\times \xi )`$ in the group $`H^{4i}(C)H^{4k}(T)`$ is equal to $`p_i(\eta ^{})\times p_k(\xi )`$. Since the cross product of nonzero classes is nonzero, $`p_{i+k}(\eta ^{}\times \xi )`$ is nonzero. On the other hand, the component of $`p_{i+k}(\eta ^{}\times \xi )`$ in $`H^{4i+4k}(C)H^0(T)`$ is $`p_{i+k}(\eta ^{})\times 1=0\times 1=0`$. We now apply 3.6 for $`Q=p_{i+k}`$ to get a contradiction. By 4.4, $`\xi `$ becomes stably trivial in a finite cover.
Now assume that either $`rank(\eta )=0`$, or $`e(\eta )0`$. By 4.4 it suffices to show that $`e(\xi )=0`$. Of course, we can assume that $`rank(\xi )>0`$. The pullback of $`\eta \times \xi `$ to $`C`$ has zero Euler class because the pullback is the Whitney sum of $`\eta `$ and a trivial bundle of the same rank as $`\xi `$. Hence according to 3.6, $`e(\eta \times \xi )=0`$. Thus if $`rank(\eta )=0`$, we get $`e(\xi )=0`$. Otherwise, note that $`e(\eta \times \xi )=e(\eta )\times e(\xi )`$ and since $`e(\eta )0`$ it implies $`e(\xi )=0`$ as wanted. ∎
###### Remark 5.2.
The assumption that $`e(\eta )0`$ is certainly necessary, in general. For example, let $`\eta `$ be the trivial line bundle over $`C`$ and let $`\xi `$ be the bundle over $`T^{2n}`$ which is the pullback of $`TS^{2n}`$ via a degree one map $`T^{2n}S^{2n}`$. Then the bundle $`\eta \times \xi `$ is trivial so its total space is nonnegatively curved whenever $`sec(C)0`$.
###### Theorem 5.3.
Let $`H^{4i}(C\times T,C)0`$ for some $`i>0`$ and let $`Q^{}`$ be a polynomial in rational Pontrjagin classes $`p_j`$ where $`0<j<i`$ such that the projection of $`p_i(TC)+Q^{}(TC)`$ to $`H^{4i}(C)`$ is zero. Then, in each rank $`>2i`$, there exist infinitely many vector bundles over $`B=C\times T`$ whose total spaces are not nonnegatively curved.
###### Proof.
Set $`Q=p_i+Q^{}`$. Since $`H^{4i}(C\times T,C)0`$, we can use 4.5 to find a vector bundle $`\xi `$ over $`C\times T`$ of rank $`2i+1`$ such that $`p_j(\xi )=0`$ for $`0<j<i`$ and $`p_i(\xi )`$ is a nonzero class whose projection to $`C`$ is zero. Using the Whitney sum formula for Pontrjagin classes we get $`p_i(TB\xi )=p_i(TB)+p_i(\xi )`$ and $`p_j(TB\xi )=p_j(TB)`$ for $`0<j<i`$. Thus, looking at the projection to $`H^{4i}(C\times T)`$, we get
$$Q(TB\xi )=p_i(TB\xi )+Q^{}(TB\xi )=p_i(TB)+p_i(\xi )+Q^{}(TB)=p_i(\xi )$$
where the last equality is true because
$$p_i(TB)+Q^{}(TB)=(p_i(TC)+Q^{}(TC))\times 1=0\times 1=0.$$
Thus, we can apply 3.6. ∎
###### Corollary 5.4.
Let $`H^{4i}(C\times T,C)0`$ for some $`i>0`$ and let $`ph_i(TC)=0`$, where $`ph_i`$ is the component of the Pontrjagin character that lives in the $`4i`$th cohomology. Then, in each rank $`>2i`$, there exist infinitely many vector bundles over $`B=C\times T`$ whose total spaces are not nonnegatively curved.
###### Proof.
Take $`Q^{}=ph_ip_i`$. ∎
###### Corollary 5.5.
Let $`H^{4i}(C\times T,C)0`$ and $`p_i(TC)=0`$ for some $`i>0`$, then there exists a vector bundle $`\xi `$ over $`C\times T`$ of any rank $`>2i`$ such that $`E(\xi )`$ has no metric of nonnegative curvature.
###### Proof.
Take $`Q=0`$. ∎
###### Remark 5.6.
In particular, 5.5 shows that if $`H^4(C\times T,C)0`$ and $`p_1(TC)=0`$, then there exist infinitely many bundles of every rank $`>2`$ whose total spaces are not nonnegatively curved.
###### Corollary 5.7.
If $`H^{2i}(C\times T,C)0`$ for some $`i>0`$, then there exist infinitely many rank $`2i`$ vector bundles over $`C\times T`$ with different Euler classes whose total spaces are not nonnegatively curved.
###### Proof.
It follows from 4.5 that there exists a bundle $`\xi `$ (and, in fact, infinitely many such bundles) of rank $`2i`$ over $`C\times T`$ such that $`e(\xi )`$ is a nonzero class whose $`H^{2i}(C)`$ component is zero. By 3.6 $`E(\xi )`$ is not nonnegatively curved. ∎
###### Corollary 5.8.
Let $`dim(C)=4m+2`$ and $`p_m(TC)0`$, and let $`dim(T)2`$. Then there exist infinitely many bundles of each rank $`2`$ over $`C\times T`$ whose total spaces are not nonnegatively curved.
###### Proof.
By Poincaré duality find $`yH^2(C)`$ such that $`p_m(TC)y0H^{4m+2}(C)`$. Take any $`tH^2(T)`$ and realize $`y1+1t`$ as the Euler class of an oriented rank $`2`$ bundle $`\xi `$ over $`C\times T=B`$. Then $`p_1(\xi )=(y1+1t)^2=y^21+2yt+1t^2`$. Note that $`p_m(TB)p_1(\xi )0`$. (Indeed, it suffices to show that the projection of $`p_m(TB)p_1(\xi )`$ to $`H^{4m+2}(C)H^2(T)`$ is nonzero. which is true because the projection is equal to $`2p_m(TC)yt`$.) Also the Whitney sum formula implies that $`p_{m+1}(TB\xi )=p_m(TB)p_1(\xi )0`$ and the projection of $`p_{m+1}(TB\xi )`$ to $`C`$ vanishes because $`H^{4m+4}(C)=0`$. Hence, we are done by 3.6. By adding trivial bundles to $`\xi `$ one can make its rank arbitrary large. ∎
###### Corollary 5.9.
If $`H^{4i}(C\times T,C)0`$ for some $`i>0`$, then there exist infinitely many vector bundles over $`B=C\times T`$ of any rank $`dim(C)`$ whose total spaces are not nonnegatively curved.
###### Proof.
Let $`\nu (C)`$ be a rank $`dim(C)`$ bundle over $`C`$ which is stably isomorphic to stable normal bundle of $`C`$. By 4.5 we can find a bundle $`\xi `$ (and, in fact, infinitely many such bundles) of rank $`dim(C)`$ over $`C\times T`$ such that the pullback of $`\xi `$ to $`C`$ is isomorphic to $`\nu (C)`$ and $`p_i(\xi )`$ has a nonzero projection to $`H^{4i}(C\times T,C)`$.
Look at the bundle $`TE(\xi |_B)\xi TB`$. Note that the pullback of $`\xi TB`$ to $`C`$ is isomorphic to $`\nu (C)TC`$ which is stably trivial, hence $`p_i(TE(\xi |_C))=0`$. On the other hand the projection of $`p(\xi TB)=p(\xi )p(TB)`$ to $`H^{4i}(C\times T,C)`$ is equal to $`p_i(\xi )`$, in particular the projection of $`p_i(\xi TB)`$ to $`H^{4i}(C\times T,C)`$ is nontrivial. By 3.6 $`E(\xi )`$ is not nonnegatively curved. ∎
###### Remark 5.10.
The method of 5.9 can be used with some other bundles in place of $`\nu (C)`$. To illustrate the idea we discuss the case when $`C`$ is the total space $`S(\eta )`$ of the sphere bundle associated with a vector bundle $`\eta `$ over $`S^4`$.
First, let us handle the easier case when $`e(\eta )0`$. Then $`rank(\eta )`$ is necessarily $`4`$ and it follows from the Gysin sequence that $`S(\eta )`$ is a rational homology $`7`$-sphere. In particular, all the Pontrjagin classes of $`S(\eta )`$ vanish. Now 4.5 implies that there are infinitely many rank $`4`$ bundles over $`S(\eta )\times T`$ with nonzero $`p_2`$. By 5.5 their total spaces admit no nonnegatively curved metrics and as usual we can add trivial bundles to make the rank $`4`$.
Now assume $`e(\eta )=0`$. Recall that $`p_1(\eta )[S^4]`$ is necessarily even and furthermore, any even integer can be realized as $`p_1(\xi )[S^4]`$ where $`\xi `$ is a $`4`$-bundle \[Mil56\]. Thus we can find a $`4`$-bundle $`\xi ^{}`$ over $`S^4`$ with $`p_1(\xi ^{})[S^4]=p_1(\eta )[S^4]`$ so that $`p_1(\eta \xi ^{})=0`$. Note that $`TS(\eta )`$ is stably isomorphic to the pullback of $`\eta `$ via the bundle projection $`\pi :S(\eta )S^4`$. Setting $`\xi =\pi ^\mathrm{\#}\xi ^{}`$, we get that $`p_1(TS(\eta )\xi )=0`$. Since $`e(\eta )=0`$, the Gysin sequence implies that $`\pi `$ induces an isomorphism on $`H^4`$, and by the Poincaré duality $`H^3(S(\eta ))0`$. Hence by 4.5 there are infinitely many rank $`4`$ bundles over $`S(\eta )\times T`$ with nonzero $`p_1`$ such that their pullback to $`S(\eta )`$ is $`\xi `$. So the proof of 5.9 applies and these bundles admit no nonnegatively curved metrics.
It is interesting to see whether $`rank(\xi )`$ can be lowered to $`3`$. Recall that an integer $`k`$ can be realized as $`p_1(\xi )[S^4]`$ for a $`3`$-bundle over $`S^4`$ iff $`k`$ is a multiple of $`4`$ \[Mil56\]. Thus, if $`p_1(\eta )[S^4]`$ is divisible by $`4`$, the argument of the previous paragraph applies and we get infinitely many $`3`$-bundles over $`S(\eta )\times T`$ with no nonnegatively curved metrics.
Now assume that $`p_1(\eta )[S^4]2\mathrm{m}\mathrm{o}\mathrm{d}(4)`$. We are looking for a $`3`$-bundle $`\xi `$ over $`S(\eta )`$ such that $`p_1(TS(\eta )\xi )=0`$. Since $`e(\eta )=0`$, the bundle $`S(\eta )S^4`$ has a section $`s`$. Setting $`\xi ^{}=s^\mathrm{\#}\xi `$, we would get a $`3`$-bundle $`\xi ^{}`$ over $`S^4`$ with $`p_1(\eta \xi ^{})=0`$. In particular, $`p_1(\xi ^{})[S^4]2\mathrm{m}\mathrm{o}\mathrm{d}(4)`$ which is impossible for a $`3`$-bundle.
Thus, the methods of this paper fail here. For instance, we do not have examples of $`3`$-bundles over $`S(\eta )\times S^1`$ that admit no nonnegatively curved metrics whenever $`p_1(\eta )[S^4]2\mathrm{m}\mathrm{o}\mathrm{d}(4)`$. Note that 4.5 produces many $`3`$-bundles over $`S(\eta )\times S^1`$ which do not become pullbacks of bundles over $`S(\eta )`$.
### Metastable range, sphere bundles, and surgery.
We now describe yet another variation of 5.9. When the method works, it gives a result similar to 5.9 with sometimes lower rank. We showed above that, under certain assumptions on the Pontrjagin classes of $`B=C\times T`$, there are vector bundles over $`B`$ whose total spaces admit no nonnegatively curved metric. Now the idea is to replace $`C`$ by a homotopy equivalent closed manifold $`C^{}`$ with “nicer” (e.g. trivial) Pontrjagin classes. Then theorems of this section can be used to produce a vector bundle over $`C^{}\times T`$ whose total space admits no nonnegatively curved metric, and can often use it to get a similar bundle over $`C\times T`$.
Indeed, let $`f:BB^{}`$ be a homotopy equivalence of closed smooth manifolds and let $`\xi `$ be a vector bundle over $`B^{}`$ with total space $`E(\xi )`$. Assume now that $`2\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{k}(\xi )dim(B)+35`$, that is, we are in the metastable range. By \[Hae61\], the homotopy equivalence $`f:BE(\xi )`$ is homotopic to a smooth embedding $`e:BE(\xi )`$. The above inequality implies that $`\mathrm{rank}(\xi )3`$, hence by \[Sie69, Thm 2.2\] $`E(\xi )`$ is diffeomorphic to the total space of the normal bundle to $`E(\nu _e)`$. Clearly, $`E(\xi )`$ is nonnegatively curved iff so is $`E(\nu _e)`$.
###### Theorem 5.11.
Let $`T`$ be a torus and let $`C`$ be a closed smooth manifold homotopy equivalent to a closed manifold $`C^{}`$ such that $`T`$ and $`C^{}`$ satisfy the assumptions of 5.3 or 5.8. Then, in each rank $`>1+dim(C\times T)/2`$, there exist infinitely many vector bundles over $`B=C\times T`$ whose total spaces are not nonnegatively curved.
###### Proof.
By 5.3 or 5.8 we can find a bundle $`\xi `$ over $`C^{}\times T`$ whose total space $`E(\xi )`$ is not nonnegatively curved in any rank $`>2i`$. Assume now that the rank is $`>1+dim(C\times T)/2`$ (note that $`1+dim(C\times T)/2>2i`$ because $`H^{4i}(C\times T,C)0`$. This puts us in the in metastable range so the homotopy equivalence $`f\times \mathrm{id}:C\times TC^{}\times TE(\xi )`$ is homotopic to an embedding whose normal bundle has total space diffeomorphic to $`E(\xi )`$. Of course, the total space of this normal bundle is not nonnegatively curved. By varying $`\xi `$ (or, rather, the Pontrjagin class of $`\xi `$), we get infinitely many such examples. ∎
One way to replace $`C`$ by a manifold $`C^{}`$ with “nicer” Pontrjagin classes is by surgery. Namely, assume $`\pi _1(C)=1`$ and let $`\tau `$ be a vector bundle over $`C`$ so that $`\tau `$ and $`TC`$ are stably fiber homotopy equivalent. Then, if the surgery obstruction vanishes (which always happens if $`dim(C)`$ is odd \[Bro72, II.3.1\]), then there is a closed smooth manifold $`C^{}`$ and a homotopy equivalence $`f:C^{}C`$ such that $`TC^{}`$ is stably isomorphic to $`f^\mathrm{\#}\tau `$.
In general, it is not easy to decide when a given vector bundle, such as $`TC`$, is stably fiber homotopy equivalent to a bundle with “nicer” Pontrjagin classes. However, each bundle with “nice” Pontrjagin classes is usually stably fiber homotopy equivalent to infinitely many different bundles.
Indeed, recall that two vector bundles are stably fiber homotopy equivalent if the corresponding spherical fibrations are stably equivalent. The stable equivalence spherical fibrations over $`C`$ (or any finite simply connected cell complex) are in one-to-one correspondence with $`[C,BSG]`$. The Whitney sum gives $`BSG`$ and $`BSO`$ an $`H`$-group structure, and the natural map $`BSOBSG`$ that assigns to a vector bundle the corresponding spherical fibration induces a group homomorphism $`[C,BSO][C,BSG]`$. After tensoring with rational the group $`[C,BSO]`$ becomes $`_{i>0}H^{4i}(C,)`$ while $`[C,BSG]`$ becomes the trivial group. In particular, if $`_{i>0}H^{4i}(C,)0`$, each stable fiber homotopy equivalence class contains infinitely many vector bundles in any rank $`dim(C)`$.
For example, let $`C`$ be the total space of a sphere bundle over closed simply-connected manifold $`V`$ associated with a vector bundle $`\eta `$. (Due to \[GW00\] such manifolds could be a good source of nonnegatively curved manifolds.) The bundle $`TC`$ is the pullback of $`TV\eta `$ via the bundle projection $`CV`$, hence $`TC`$ is stably fiber homotopy trivial whenever so are $`TV`$ and $`\eta `$. This construction gives many manifolds with stably fiber homotopy trivial tangent bundles.
## 6. Sphere bundles with no metric of nonnegative curvature
It was shown in \[GW00\] that the total space of the sphere bundle associated with the normal bundle to the soul has a nonnegatively curved metric. Thus, potentially, sphere bundles provide a good source of closed nonnegatively curved manifolds.
###### Theorem 6.1.
For $`k>0`$, let $`EF`$ be a $`k`$-sphere Serre fibration over a flat manifold $`F`$ with nonzero rational Euler class. Let $`P`$ be closed smooth manifold such that there is a map $`PE`$ that induces an isomorphism of fundamental groups. Then $`P`$ admits no metric of nonnegative Ricci curvature.
###### Proof.
Arguing by contradiction, assume that $`P`$ admits a metric of nonnegative Ricci curvature. Pass to a finite cover $`\stackrel{~}{P}P`$ so that $`\stackrel{~}{P}`$ is diffeomorphic to $`C\times T`$ where $`C`$ is simply connected and $`T`$ is a torus.
Look at the corresponding covers $`\stackrel{~}{E}E`$ and $`\stackrel{~}{F}F`$. Note that $`\pi _1(\stackrel{~}{F})`$ is free abelian because $`\pi _1(\stackrel{~}{F})`$ is a torsion free group which is the image of a finitely generated abelian group $`\pi _1(\stackrel{~}{E})\pi _1(\stackrel{~}{P})`$. The $`k`$-sphere fibration $`\stackrel{~}{E}\stackrel{~}{F}`$ still has nonzero rational Euler class since since it is a pullback of $`EF`$ and since finite covers induce injective maps on rational cohomology.
First consider the case $`k=1`$. The circle fibration $`\stackrel{~}{E}\stackrel{~}{F}`$ induces an epimorphism $`\varphi :\pi _1(\stackrel{~}{E})\pi _1(\stackrel{~}{F})`$ of finitely generated free abelian groups. Therefore, $`\varphi `$ has a section. Since $`\stackrel{~}{E}`$ is aspherical, this section is induced by a continuous map $`\stackrel{~}{F}\stackrel{~}{E}`$ which defines a homotopy section of the circle fibration $`\stackrel{~}{E}\stackrel{~}{F}`$. Thus, the Euler class must be zero which is a contradiction.
Now assume that $`k>1`$ so that $`\stackrel{~}{E}\stackrel{~}{F}`$ induces a $`\pi _1`$-isomorphism. Then the inclusion $`T(C\times T)=\stackrel{~}{P}`$ followed by the map $`\stackrel{~}{P}\stackrel{~}{E}\stackrel{~}{F}`$ induces a $`\pi _1`$-isomorphism hence is a homotopy equivalence. Let $`s`$ be its homotopy inverse. Then $`s`$ followed by the inclusion $`T\stackrel{~}{P}`$ and the map $`\stackrel{~}{P}\stackrel{~}{E}`$ is a homotopy section of the fibration $`\stackrel{~}{E}\stackrel{~}{F}`$. The Euler class then must be zero which gives a contradiction. ∎
###### Remark 6.2.
The above argument is a special case of the following phenomenon. Suppose we have a Serre fibration $`CPF`$ where $`F`$ is a flat manifold and $`C`$ is connected and simply-connected. Look at the spectral sequence of this fibration with rational coefficients. Then if there exists a nonzero differential, then $`P`$ does not admit a nonnegatively curved metric.
Indeed, if $`P`$ is nonnegatively curved, then a finite cover $`\stackrel{~}{P}`$ of $`P`$ splits topologically as $`M\times T`$ where $`M`$ is simply connected and $`T`$ is a torus. By naturality we can see that spectral sequence of the pullback fibration $`C\stackrel{~}{P}T`$ also has a nonzero differential. Since the universal cover of $`\stackrel{~}{P}`$ is homotopy equivalent to both $`M`$ and $`C`$, they are homotopy equivalent to each other. In particular $`dimH^{}(M)=dimH^{}(C)`$ and hence $`dimH^{}(\stackrel{~}{P})=dimH^{}(M)H^{}(T)=dim(H^{}(C)H^{}(T))`$. On the other hand if there is a nonzero differential we should have that $`dimH^{}(\stackrel{~}{P})<dim(H^{}(C)H^{}(T))`$ which is a contradiction.
###### Theorem 6.3.
Let $`E(\xi )`$ be the total space of a vector bundle $`\xi `$ over a closed smooth manifold $`B`$ and let $`S(\xi )B`$ be the associated sphere bundle. Assume that $`\xi `$ has zero rational Euler class and there exists a polynomial $`Q`$ in rational Pontrjagin classes such that $`Q(TE(\xi ))0`$ and $`Q(T\stackrel{~}{E}(\xi ))=0`$ for the universal cover $`\pi :\stackrel{~}{E}(\xi )E(\xi )`$. Then $`S(\xi )`$ admits no metric of nonnegative Ricci curvature.
###### Proof.
First, we introduce several notations. Let $`q:\stackrel{~}{B}B`$ be the universal covering and $`j:S(\xi )E(\xi )`$ be the inclusion. Also denote by $`\pi `$ and $`i`$ the bundle projection and the zero section of $`\xi `$, respectively.
Since $`S(\xi )E(\xi )`$ is an oriented codimension one hypersurface, $`TS(\xi )`$ is stably isomorphic to $`j^\mathrm{\#}TE(\xi )`$, hence $`Q(TS(\xi ))=j^{}Q(TE(\xi ))`$. Also $`TE(\xi )`$ is isomorphic to $`(i\pi )^\mathrm{\#}TE(\xi )`$ since $`i\pi `$ is homotopic to the identity of $`TE(\xi )`$. We get $`Q(TS(\xi ))=j^{}\pi ^{}Q(i^\mathrm{\#}TE(\xi ))`$. By assumption $`Q(TE(\xi ))`$, and hence $`Q(i^\mathrm{\#}TE(\xi ))`$ is nonzero. Also $`j^{}\pi ^{}=(\pi j)^{}`$ where $`\pi j:S(\xi )B`$ is the bundle projection. It follows from the Gysin sequence that $`\pi j`$ is injective in cohomology because the kernel of $`(\pi j)^{}`$ consists of the cup-multiples of the Euler class which is zero by assumption. Thus, $`Q(TS(\xi ))0`$.
On the other hand, the inclusion $`S(q^\mathrm{\#}\xi )E(q^\mathrm{\#}\xi )`$ takes $`Q(TE(q^\mathrm{\#}\xi ))=0`$ to $`Q(TS(q^\mathrm{\#}\xi ))`$. Hence, $`Q(TS(q^\mathrm{\#}\xi ))=0`$ and we are in position to apply the theorem 3.4. ∎
###### Corollary 6.4.
Let $`\xi `$ be a bundle over a flat manifold $`F`$ with associated sphere bundle $`S(\xi )`$ and let $`C`$ be a closed smooth simply-connected manifold. If $`C\times S(\xi )`$ admits a metric of nonnegative Ricci curvature, then $`\xi `$ becomes trivial in a finite cover.
###### Proof.
By 4.4 it suffices to show that $`e(\xi )=0`$ and $`p(\xi )=1`$. Vanishing of $`e(\xi )`$ follows from 6.1. Vanishing of all Pontrjagin classes follows exactly as in the proof of 5.1 where instead of referring to 3.4 we use 6.3. ∎
## 7. The classification of nonnegatively curved vector bundles over $`S^1\times S^3`$
In this section we prove the theorem 1.4. Note that the converse of 1.4 is trivially true, i.e. both the trivial bundle and the product of the trivial bundle over $`S^3`$ and Möbius band line bundle over $`S^1`$ are nonnegatively curved.
###### Proof of 1.4.
Any vector bundle over a $`4`$-complex is the Whitney sum of a trivial bundle and a bundle of rank $`4`$, hence it suffices to consider the bundles of rank $`k`$ at most $`4`$.
First, assume that $`\xi `$ is orientable. Let $`q:\times S^3S^1\times S^3`$ be the universal cover of $`S^1\times S^3`$. Then since any vector bundle over $`S^3`$ is trivial we have that $`q^\mathrm{\#}(\xi )`$ is trivial. In particular, $`p_1(q^\mathrm{\#}(\xi ))=e(q^\mathrm{\#}(\xi ))=0`$. Therefore, according to 3.4. the classes $`p_1(\xi )`$ and $`e(\xi )`$ vanish. Thus, it suffices to prove the following.
###### Lemma 7.1.
Let $`\eta `$ be an orientable vector bundle over $`S^1\times S^3`$ such that $`p_1(\eta )=e(\eta )=0`$. Then $`\eta `$ is trivial.
###### Proof.
Since $`H^1(S^1\times S^3,/2)=0=H^2(S^1\times S^3,)`$, any rank one or rank two orientable bundle over $`S^1\times S^3`$ is trivial.
Assume that $`\eta `$ is an orientable bundle of rank $`4`$. Let $`f:S^1\times S^3BSO(4)`$ denote a classifying map for $`\eta `$, i.e $`\eta f^{}\gamma ^4`$ where $`\gamma ^4`$ is the universal $`4`$-bundle over $`BSO(4)`$. The first four homotopy groups of $`BSO(4)`$ are as follows: $`\pi _0(BSO(4))=0,\pi _1(BSO(4))=0,\pi _2(BSO(4))=/2,\pi _3(BSO(4))=0`$ and $`\pi _4(BSO(4))=`$. Consider the standard product cell decomposition of $`S^1\times S^3`$ coming from canonical cell decompositions $`S^1=e^0e^1`$ and $`S^3=e^0e^3`$. Then the $`3`$-skeleton of $`S^1\times S^3`$ is the wedge $`S^1S^3`$. Since any orientable vector bundle over $`S^1`$ or $`S^3`$ is trivial, $`f|_{S^1{\scriptscriptstyle S^3}}`$ is homotopic to a point and therefore by the homotopy extension property we can assume that $`f`$ send $`S^1S^3`$ to a point to begin with. In other words, $`f`$ can be written as a composition $`f=\overline{f}\pi `$ where $`\pi `$ is the factorization map $`\pi :S^1\times S^3S^1\times S^3/(S^1S^3)S^4`$. Since $`\pi `$ induces an isomorphism on $`H^4`$, the bundle $`\overline{f}^{}(\gamma ^4)`$ has zero Euler and Pontrjagin classes. It is a well known that a bundle over $`S^4`$ with zero Euler and Pontrjagin classes is trivial. (Indeed, the map $`(e,p_1):\pi _4(BSO(4))`$ which associates to a $`4`$-bundle over $`S^4`$ its Euler and Pontrjagin classes is a rational homotopy equivalence. Then the induced map on $`\pi _4`$ has finite, and hence trivial, kernel because $`\pi _4(BSO(4))\times `$.) Thus $`\overline{f}`$, and hence $`f`$, is nullhomotopic.
A very similar argument shows that any orientable $`3`$-bundle over $`S^1\times S^3`$ with zero first Pontrjagin class is trivial. Again, everything can be reduced to 3-bundles over $`S^4`$ with zero $`p_1`$. The map $`p_1:BSO(3)K(,4)`$ is a rational homotopy equivalence, in particular, the induced map on $`\pi _4`$ has finite, and hence trivial, kernel because $`\pi _4(BSO(3))`$. Hence, only the trivial $`3`$-bundle over $`S^4`$ has zero $`p_1`$. ∎
Now suppose that $`\xi `$ is not orientable and its total space admits a metric of nonnegative curvature. Since $``$ has a unique subgroup of index $`2`$ the orientation double cover for $`\xi `$ is given by the map $`\pi _S=(zz^2)\times \mathrm{id}:S^1\times S^3S^1\times S^3`$. Then the pullback $`\pi _S^\mathrm{\#}(\xi )`$ is orientable and also admits a metric of nonnegative curvature. By above, the pullback $`\pi _S^\mathrm{\#}(\xi )`$ is trivial. The following lemma completes the proof of 1.4 in the nonorientable case. ∎
###### Lemma 7.2.
Let $`\eta `$ be a nonorientable rank $`k`$ bundle over $`S^1\times S^3`$ whose orientation lift is a trivial bundle. Then $`\eta `$ is isomorphic to the product $`\mu ^1\times ϵ^{k1}`$ of the Möbius band line bundle $`\mu ^1`$ over $`S^1`$ and a trivial rank $`(k1)`$-bundle $`ϵ^{k1}`$ over $`S^3`$.
###### Proof.
Since $`H^1(S^1\times S^3,/2)=/2`$, there is only one nonorientable line bundle over $`S^1\times S^3`$, namely, $`\mu ^1\times ϵ^0`$.
Case of rank four. Let $`f:S^1\times S^3BO(4)`$ be the classifying map for $`\eta `$ and $`f_0:S^1\times S^3BO(4)`$ be the classifying map for $`\mu ^1\times ϵ^3`$. We want to show that these maps are homotopic. The same argument as in the proof of 7.1 shows that $`f`$ and $`f_0`$ are homotopic on the $`3`$-skeleton. Let us show that this homotopy can be extended over the $`4`$-cell.
Let $`\pi _B:BSO(4)BO(4)`$ be the canonical double cover. Then each of the maps $`f\pi _S`$ and $`f_0\pi _S`$ lifts to a map $`\stackrel{~}{f}:S^1\times S^3BSO(4)`$ which is the classifying map for $`\pi _S^{}\eta `$. In other words, we have the following commutative diagram
$$\begin{array}{ccc}S^1\times S^3& \stackrel{\stackrel{~}{f}}{}& BSO(4)\\ \pi _S& & \pi _B\\ S^1\times S^3& \stackrel{𝑓}{}& BO(4)\end{array}$$
By construction, the map $`\stackrel{~}{f}`$ is equivariant under the action of the group of deck transformations $`/2`$ where the nontrivial element $`i/2`$ acts on $`S^1\times S^3`$ by the formula $`(z,q)(z,q)`$ and acts on $`BSO(4)`$ by reversing orientations of $`4`$-planes.
Clearly the maps $`f`$ and $`f_0`$ are homotopic iff the maps $`\stackrel{~}{f}`$ and $`\stackrel{~}{f}_0`$ are equivariantly homotopic. By above we can assume that $`\stackrel{~}{f}`$ and $`\stackrel{~}{f}_0`$ are equivariantly homotopic on the 3-skeleton of $`S^1\times S^3`$. Next we compute the equivariant cohomology group $`H_{eq}^4(S^1\times S^3,\{\pi _4(BSO(4))\})`$ that contains the obstruction for extending the homotopy over the $`4`$-skeleton and show that in our situation the obstruction has to vanish.
In order to explicitly describe equivariant cochains we have to identify the action of $`/2`$ on $`\pi _4(BSO(4))`$. Recall that $`\pi _4(BSO(4))`$ classifies the isomorphism classes of orientable $`4`$-bundles over $`S^4`$. On the other hand, $`\pi _4(BSO(4))\pi _3(SO(4))`$ where the last isomorphism can be described explicitly as $`(m,n)(qq^nvq^m)`$ where we identify $`^4`$ with the quaternions $``$ and $`S^3`$ with the set of unit quaternions $`q`$. According to \[Mil56\], the classes $`p_1,eH^4(S^4,)`$ of the bundle $`(m,n)`$ are given by $`p_1(m,n)=2(mn)`$ and $`e(m,n)=m+n`$. The action of $`i`$ on $`BSO(4)`$ sends the canonical oriented $`4`$-bundle $`\gamma ^4`$ to $`\gamma ^4`$ (i.e the same bundle with its orientation reversed). Therefore, $`i^{}(p_1(\gamma ^4))=p_1(\gamma ^4)`$ and $`i^{}(e(\gamma ^4))=e(\gamma ^4)`$, and hence the action of $`i`$ on $`\pi _3(SO(4))`$ is given by $`i(m,n)=(n,m)`$.
Now once the action is identified, a straightforward computation shows that $`H_{eq}^4(S^1\times S^3,\{\pi _4(BSO(4))\})()/\mathrm{diagonal}`$. Let $`O_4H_{eq}^4(S^1\times S^3,\{\pi _4(BSO(4))\})`$ be the obstruction for the equivariant extension of the homotopy between $`\stackrel{~}{f}`$ and $`\stackrel{~}{f_0}`$ over the $`4`$-skeleton. It remains to show that $`O_4`$ vanishes. The double cover $`\pi _S`$ induces a homomorphism
$$\pi _S^{}:H_{eq}^4(S^1\times S^3,\{\pi _4(BSO(4))\})H_{eq}^4(S^1\times S^3,\pi _S^\mathrm{\#}\{\pi _4(BSO(4))\})$$
where the last group is equal to $`H^4(S^1\times S^3,\pi _4(BSO(4))`$ because the pullback bundle of coefficients $`\pi _S^\mathrm{\#}\{\pi _4(BSO(4))\}`$ is trivial. We claim that this map is injective. (Indeed, since both groups are isomorphic to $``$, it suffices to show that the map is nonzero. If $`\pi _S^{}`$ were zero, the orientation lift of any nonorientable $`4`$-bundle over $`S^1\times S^3`$ would be trivial which is certainly not the case since there exist nonorientable bundles over $`S^1\times S^3`$ with nonzero $`p_1`$. An example of such a bundle is the Whitney sum of a nontrivial line bundle and the pullback of a $`3`$-bundle over $`S^4`$ with nonzero $`p_1`$ via a degree one map $`S^1\times S^3S^4`$.) Since both $`f\pi _S`$ and $`f_0\pi _S`$ are null homotopic we know that $`\pi _S^{}(O_4)=0`$ and hence $`O_4`$ vanishes.
Case of rank three and two. Again, the classifying maps $`f`$ and $`f_0`$ are homotopic on the $`3`$-skeleton and one has to compute the obstruction $`O_4`$ to extending the homotopy over the $`4`$-skeleton.
If the rank is three, $`/2`$ action on the coefficient group is trivial and the obstruction group $`H_{eq}^4(S^1\times S^3,\{\pi _4(BSO(3))\})`$ reduces to $`H^4(S^1\times S^3,)`$. As before we have $`\pi _S^{}(O_4)=0`$ and since in this case $`\pi _S^{}`$ is clearly injective, we conclude that $`O_4`$ vanishes. In the rank two case the obstruction is always zero simply because $`\pi _4(BSO(2))=0`$. ∎
###### Proposition 7.3.
The total space of any vector bundle over $`S^1\times S^2`$ has a complete metric of nonnegative curvature such that the zero section is a soul.
###### Proof.
Since all vector bundles over $`S^1`$ and $`S^2`$ admit nonnegatively curved metric such that the zero sections are souls, it suffices to show that any vector bundle over $`S^1\times S^2`$ is isomorphic to the product of a bundle over $`S^1`$ and a bundle over $`S^2`$.
First of all observe that two rank $`k`$ vector bundles over $`S^1\times S^2`$ are isomorphic iff their restrictions to the two-skeleton $`S^1S^2`$ are isomorphic. Indeed, we only need to extend the homotopy of the classifying maps $`S^1\times S^2BO(k)`$ from $`S^1S^2`$ to the remaining $`3`$-cell. This is always possible since $`\pi _3(BO(k))\pi _2(O(n))=0`$.
Now let $`\xi `$ be a vector bundle of rank $`k`$ over $`S^1\times S^2`$ with the classifying map $`f:S^1\times S^2BO(k)`$.
The case $`k=1`$ is obvious because line bundles are classified by $`w_1`$ and the inclusion $`S^1S^1\times S^2`$ induces an isomorphism on $`H^1(,/2)`$ so that any line bundle over $`S^1\times S^2`$ is a pullback of a bundle over $`S^1`$. Similarly, if $`k=2`$ and $`\xi `$ is orientable, then $`\xi `$ is completely determined by its Euler class. Since the inclusion $`S^2S^1\times S^2`$ induces an isomorphism on $`H^2(,)`$, we conclude that $`\xi `$ is a pullback of a bundle over $`S^2`$.
Assume that $`k3`$. Since $`\pi _2(BSO(k))=/2`$, there are exactly two $`k`$-bundles over $`S^2`$, namely the trivial bundle and the Whitney sum of a trivial bundle and a $`2`$-bundle with nonzero $`w_2`$. Let $`\kappa `$ be a $`2`$-bundle over $`S^2`$ that has the same $`w_2`$ as the restriction of $`\xi `$ to $`S^2`$ and let $`\lambda `$ be a line bundle over $`S^1`$ with the same $`w_1`$ as the restriction of $`\xi `$ to $`S^1`$. Finally, let $`g`$ be the classifying map for the the Whitney sum of $`\kappa \times \lambda `$ and the trivial bundle of rank $`(k3)`$. By construction the restrictions of $`f`$ and $`g`$ to the two-skeleton $`S^1S^2`$ are homotopic as needed.
Finally, suppose that $`k=2`$ and $`\xi `$ is not orientable. Note that the orientable two-fold cover $`\stackrel{~}{\xi }`$ of $`\xi `$ has zero Euler class. (Indeed, since $`S^2`$ represents the generator of $`H_2(S^1\times S^2,)`$ it suffices to show that the intersection number of $`S^2`$ and the zero section of $`\stackrel{~}{\xi }`$ inside the total space $`E(\stackrel{~}{\xi })`$ is zero. To compute the intersection number put $`S^2`$ in the general position to the zero section of $`\xi `$ and then look at the the preimage of the manifolds inside $`E(\stackrel{~}{\xi })`$. The covering action of $`/2`$ on $`E(\stackrel{~}{\xi })`$ preserves the orientation on the base and changes the orientation of the total space. Thus, points of intersection come in pairs: one with plus sign and the other with minus sign. So the intersection number is zero.) This implies that the restriction of $`\xi `$ to $`S^2`$ has zero Euler class, and so $`\xi |_{S^2}`$ is a trivial bundle. By above $`\xi `$ is isomorphic to the product of $`\xi |_{S^1}`$ and the rank zero bundle over $`S^2`$. ∎ |
warning/0001/astro-ph0001220.html | ar5iv | text | # PROBING THE SITE FOR r-PROCESS NUCLEOSYNTHESIS WITH ABUNDANCES OF BARIUM AND MAGNESIUM IN EXTREMELY METAL-POOR STARS
## 1 INTRODUCTION
Truran (1981) suggested that the observed presence of elements heavier than iron (Fe) in extremely metal-poor stars is due to the nucleosynthesis products of rapid neutron capture reactions (the $`r`$-process). Europium (Eu), for which about 97% of its abundance in the solar system is of pure $`r`$-process origin (Käppler, Beer, & Wisshak 1989), has therefore been considered as one of the most useful elements for locating the astrophysical sites of $`r`$-process nucleosynthesis. Assuming that such sites would most likely be identified with Type II supernovae (SNe II), many authors have attempted to constrain the applicable mass range of SN II progenitors from an observed \[Eu/Fe\] vs. \[Fe/H\] relation for metal-poor stars. However, a satisfactory mass range for the progenitor stars has not yet been agreed upon. Mathews, Bazan, & Cowan (1992) suggested a mass range of $`M_{\mathrm{ms}}=78M_{}`$ for the astrophysical $`r`$-process site, while Travaglio et al. (1999) supported a somewhat higher mass range, $`M_{\mathrm{ms}}=810M_{}`$. Most recently, Ishimaru & Wanajo (1999) concluded that the \[Eu/Fe\] versus \[Fe/H\] relation alone is unable to distinguish an $`r`$-process site with progenitor masses $`M_{\mathrm{ms}}=810M_{}`$ from that resulting from a site with $`M_{\mathrm{ms}}>30M_{}`$.
The discovery of numerous Galactic halo stars with extremely low metal abundances from the ongoing HK survey of Beers and colleagues (Beers, Preston, & Shectman 1992; Beers 1999), in particular those in the abundance range $`4.0[\mathrm{Fe}/\mathrm{H}]2.5`$, have opened the door for much more detailed investigations of the elemental production from SNe II in the early Galaxy. Shigeyama & Tsujimoto (1998) recently argued that the elemental abundance patterns observed in the atmospheres of extremely metal-poor stars retain those produced by individual SNe II (see also Audouze & Silk 1995), and thus can be used to estimate the heavy element yields of the first generation of stars in the Galaxy. They also argued, on the basis of theoretical SN progenitor models, that the synthesized mass of magnesium, $`M_{\mathrm{Mg}}`$, as a function of $`M_{\mathrm{ms}}`$, is likely to be far more certain than the predicted mass of iron produced, $`M_{\mathrm{Fe}}`$, which is dependent on the mass cut chosen in the stellar core. These considerations indicate that the observed \[Eu/Mg\] vs. \[Mg/H\] relation in extremely metal-poor stars should give a much firmer constraint on the $`r`$-process site. According to Tsujimoto & Shigeyama (1998), stars with the lowest value of \[Mg/H\]$`=2.7`$ (CS 22892-052) that exhibit absorption lines of Eu in their spectra in the sample of McWilliam et al. (1995) set the lower limit to $`M_{\mathrm{ms}}20M_{}`$, above which SNe II produce the $`r`$-process elements.
For most stars of extremely low metallicity (\[Fe/H\] $`\stackrel{<}{}`$$``$ 2.0), obtaining an abundance measurement of Eu with 2.5–4m-class telescopes is a difficult observational task. For example, McWilliam et al. (1995) only obtained Eu abundances for 11 of the 33 metal-poor stars in their sample. This incompleteness implies that one might miss valuable insights to $`r`$-process production mechanisms by concentrating exclusively on observations of Eu. McWilliam (1998) showed that the \[Ba/Eu\] ratios for stars with $`3\stackrel{<}{}`$\[Fe/H\]$`\stackrel{<}{}2`$ are consistent with pure $`r`$-process ratios, and concluded that Ba in extremely metal-poor stars is also of $`r`$-process origin. Burris et al. (1999) arrived at a similar result.
In Figure 1, which uses the data of McWilliam et al. (1995) and McWilliam (1998), the \[Ba/Mg\] values for a sample of metal-poor stars are plotted against \[Mg/H\], together with four heavier $`r`$-process elements, La, Ce, Nd, and Eu. It is indeed remarkable that there appears a nearly vertical boundary at \[Mg/H\]$`2.5`$, with \[Ba/Mg\] spanning the range from \[Ba/Mg\]$`=2`$ to +0.6, whereas a horizontal line of \[Ba/Mg\]$`1.4`$ emerges from \[Mg/H\]$`2.5`$ down to $`3.7`$. The data of Ryan, Norris, & Beers (1996) for seven of their extremely metal-poor stars show a similar trend in the \[Ba/Mg\]$``$\[Mg/H\] plane.
In §2 we discuss the $`r`$-process sites which might be consistent with the vertical \[Ba/Mg\]$``$\[Mg/H\] boundary, based on the models proposed by Shigeyama & Tsujimoto (1998) and Tsujimoto & Shigeyama (1998). In §3 the entire range of observed \[Ba/Mg\]$``$\[Mg/H\] abundances in halo stars is used to explore the early stages of inhomogeneous chemical evolution in the Galaxy in the context of the SN-induced star-formation model described by Tsujimoto, Shigeyama, & Yoshii (1999, hereafter TSY).
## 2 PRODUCTION SITE FOR r-PROCESS ELEMENTS
Figure 2a illustrates the hypothesis we seek to defend, that the metal-poor stars of the Galaxy populate two separate branches in the \[Ba/Mg\]$``$\[Mg/H\] plane. The first branch extends rightward of the vertical boundary which begins at (\[Ba/Mg\], \[Mg/H\])=($`2.0,2.3`$) and ends at (\[Ba/Mg\], \[Mg/H\])=($`+0.6,2.7`$). The second branch is horizontal from \[Mg/H\]$`2.5`$ to $`3.7`$ at a constant value of \[Ba/Mg\]$`1.4`$.
Other $`r`$-process elements, such as La, Ce, Nd, and Eu, shown in Figure 1, populate the first branch identified from the behavior of \[Ba/Mg\] in Figure 2, but we presently lack measurements of their abundances in the most metal-poor stars, which are required in order to confirm the presence of the second branch. Tsujimoto & Shigeyama (1998) derived the Eu yield as a function of SN II progenitor mass, and showed that average predicted value, weighted by the Salpeter initial mass function (IMF), successfully reproduces a plateau value of \[Eu/Fe\] seen at $`2\stackrel{<}{}`$\[Fe/H\]$`\stackrel{<}{}1`$. Thus, hereafter, we refer to the first branch as the $`y`$-branch, where the letter “$`y`$” stands for its origination in individual SNe II yields.
As shown in Figure 2a, the $`y`$-branch is confined to stars with \[Mg/H\]$`>2.7`$; there is no star with \[Mg/H\]$`\stackrel{<}{}`$$`2.7`$ that belongs to the $`y`$-branch. On the other hand, Shigeyama & Tsujimoto (1998) showed that the metallicity \[Mg/H\] of stars born from an SN remnant (SNR) is well approximated by the average \[Mg/H\] inside the shell swept up by the SNR. This gives a relation between the metallicity \[Mg/H\] of stars and the mass $`M_{\mathrm{ms}}`$ of SN II progenitor as shown in Figure 2b. As a consequence, it is indicated that only SNe II with $`M_{\mathrm{ms}}`$$`\stackrel{>}{}`$$`20M_{}`$ produce Ba via the $`r`$-process. If we assume that the vertical boundary to the $`y`$-branch has a one-to-one correspondence to the Ba yield from individual SNe II, the progenitor mass range is confined to $`M_{\mathrm{ms}}=2025M_{}`$. It is straightforward to derive the Ba yield from the observed \[Ba/Mg\]-value along the vertical boundary and the synthesized Mg mass, leading to $`M_{\mathrm{Ba}}=8.5\times 10^6M_{}`$ for $`M_{\mathrm{ms}}=20M_{}`$, and $`M_{\mathrm{Ba}}=4.5\times 10^8M_{}`$ for $`M_{\mathrm{ms}}=25M_{}`$.
We now consider the origin of the second branch of \[Ba/Mg\], in the range of $`3.7\stackrel{<}{}`$\[Mg/H\]$`\stackrel{<}{}2.7`$. This \[Mg/H\] range corresponds to $`M_{\mathrm{ms}}=1220M_{}`$, for which SNe II, according to our present models, do not produce significant amounts of Ba. Therefore, our hypothesis is that the Ba abundances which are observed in the atmospheres of these stars come only from the interstellar matter (ISM) that was enriched in Ba by the preceding SNe II (with $`M_{\mathrm{ms}}=2025M_{}`$), and that was swept up in the shells by later SNe II with $`M_{\mathrm{ms}}=1220M_{}`$. Thus, hereafter we refer to this branch, as the $`i`$-branch, where the letter “$`i`$” stands for an ISM origin. We note that the \[Ba/Mg\] value in this branch should always be below its plateau value \[Ba/Mg\]$`0.6`$ at higher metallicities. We also predict from the above argument that stars born from the shell swept up by SNe II with $`M_{\mathrm{ms}}>25M_{}`$ would form another $`i`$-branch at \[Mg/H\]$`\stackrel{>}{}2.5`$, evidence for which is not seen in the sample of McWilliam et al. (1995).
## 3 CHEMICAL EVOLUTION OF $`R`$-PROCESS ELEMENTS IN THE GALACTIC HALO
In this section, we discuss the chemical evolution of Ba and Eu in the Galactic halo, based on the formulation presented in TSY. The essence of the TSY model is that the star-forming process is assumed to be confined to separate clouds which make up the entire halo, and that the chemical evolution in these clouds proceeds through a successive sequence of SN II explosion, shell formation, and resulting star formation. Details can be found in TSY.
An overview of our model applied to the Galactic halo is as follows: – (1) The metal-free Population III stars (Pop III) form by some (as yet unspecified) mechanism in primordial-composition gas clouds of the Galactic halo. (2) The most massive stars among them explode as Pop III SNe II, which trigger a series of star formation events. (3) Star formation terminates when SNRs become unable to the formation of dense shells. (4) Roughly 90 % of the cloud mass remains unused in star formation, and may fall onto the still-forming Galactic disk.
The free parameters in our model are the mass fraction $`x_{\mathrm{III}}`$ of metal-free Pop III stars initially formed in each cloud, and the mass fraction $`ϵ`$ of stars formed in a dense shell swept up by each SNR. A value of $`x_{\mathrm{III}}`$ sensitively determines the level of \[Ba/Mg\] for the $`i`$-branch, so we take $`x_{\mathrm{III}}=2.5\times 10^4`$ to be consistent with the observed value \[Ba/Mg\]=$`1.4`$ at \[Mg/H\]$`3.7`$ (McWilliam 1998). We take $`ϵ=4.3\times 10^3`$ to reproduce the observed \[Fe/H\] distribution function of halo field stars for \[Fe/H\]$`<1`$ (TSY).
Adopting the result of §2, that SNe II with $`M_{\mathrm{ms}}=2025M_{}`$ are the dominant site for $`r`$-process nucleosynthesis, and using the derived Ba yield dependent on $`M_{\mathrm{ms}}`$, we calculate the Ba evolution at early epochs of the Galactic halo. For later evolution, which pertains to the abundance range $`2`$$`\stackrel{<}{}`$\[Mg/H\]$`\stackrel{<}{}`$$`1`$, we add the contribution from $`s`$-process nucleosynthesis, assuming that the production site for early-epoch $`s`$-process elements is $`23`$ $`M_{}`$ AGB stars (Burris et al. 1999). The stellar IMF is assumed to be of the Salpeter form. The upper and lower mass limits for stars that explode as SNe II are taken to be $`50M_{}`$ and $`12M_{}`$, respectively.
Figure 3a is a color-coded predicted frequency distribution of stars in the \[Ba/Mg\]$``$\[Mg/H\] plane, normalized to unity when integrated over the entire area (see the color bar for the scale). In order to enable a direct comparison with the data, the frequency distribution has been convolved with a Gaussian with $`\sigma =0.2`$ dex for \[Ba/Mg\] and $`\sigma =0.1`$ dex for \[Mg/H\]. As already pointed by Burris et al. (1999), there is a systematic difference in the derived Ba abundance between Burris et al. (1999) and McWilliam (1998). Burris et al. (1999) obtain \[Ba/H\] ratios which are higher, in the mean, by 0.46 dex for the seven stars in common with the data of McWilliam (1998). We therefore decrease their \[Ba/H\] ratios by 0.46 dex to enable a consistent comparison with the data of McWilliam (1998) in Figure 3a. The \[Ba/Mg\] ratio “plateau” in the range of $`2`$$`\stackrel{<}{}`$\[Mg/H\]$`\stackrel{<}{}`$$`1`$ is obviously higher than that of $`i`$-branch (\[Mg/H\] $`\stackrel{<}{}`$$`2.5`$), which implies a significant contribution from $`s`$-process nucleosynthesis. For comparison, we show by the dashed line in this figure the conventional one-zone chemical evolution model assuming that the production sites for $`r`$-process and $`s`$-process are $`810`$ $`M_{}`$ supernovae and $`23`$ $`M_{}`$ AGB stars, respectively. It is clear that the conventional model predicts a much smaller \[Ba/Mg\] than is consistent with the observed $`i`$-branch at \[Mg/H\] $`\stackrel{<}{}`$$`2.7`$.
To summarize our hypothesis: The earliest stars, stars that were formed from the shells swept up by SNe II with $`M_{\mathrm{ms}}=2025M_{}`$, first create the vertical boundary seen in Figure 2a. Thereafter, the $`y`$-branch develops at higher \[Mg/H\] and converges to the plateau level of \[Ba/Mg\] at higher metallicities, as described in detail in TSY. At the same time, the subsequently formed stars, arising from the SNRs with $`M_{\mathrm{ms}}=1220M_{}`$, do not produce Ba, and therefore have the lowest value of \[Ba/Mg\] at \[Mg/H\]$`\stackrel{<}{}2.5`$. These stars populate the lower envelope of the $`i`$-branch, then fill in this branch at lower \[Mg/H\], but gradually increasing \[Ba/Mg\]. As the Mg abundance in the ISM becomes dominant over that ejected from SNe II, the $`i`$-branch develops at higher abundances, increasing both \[Mg/H\] and \[Ba/Mg\].
Because the upper mass limit in our present model for the Ba production site, $`25M_{}`$, is set below the canonical upper limit $`50M_{}`$ for stable stellar masses, we predict the existence of another $`i`$-branch in the range of $`2.3\stackrel{<}{}`$\[Mg/H\]$`\stackrel{<}{}1.5`$, which is made up of stars arising from the SNRs with $`M_{\mathrm{ms}}=2550M_{}`$. While this branch does not show up in the sample of McWilliam (1998), the data from Luck & Bond (1981,1985), interestingly, support such a prediction (see filled and open circles in Fig.3a). If we increase the upper mass limit to $`50M_{}`$ for the Ba production site, a small amount of Ba is produced in SNe with $`M_{\mathrm{ms}}\stackrel{>}{}25M_{}`$, which only slightly contaminates the $`i`$-branch stars. The results of our calculations with an upper mass limit of $`M_{\mathrm{ms}}=25M_{}`$ remain almost unchanged, because of the assumed declining IMF at high masses.
The value of $`x_{\mathrm{III}}=2.5\times 10^4`$ adopted here can be converted into a probability of finding Pop III stars, $`p_{\mathrm{III}}=2.5\times 10^3`$, defined as the expected number of Pop III stars divided by the total number of the long-lived stars in a sample (see TSY). However, we regard this $`x_{\mathrm{III}}`$-value as an upper bound, because the predicted \[Ba/Mg\] vs. \[Mg/H\] distribution with smaller $`x_{\mathrm{III}}`$ provides an adequate fit to the observed distribution. A much more strict constraint on $`p_{\mathrm{III}}`$ could be obtained if the sample of extremely metal-poor stars with measured Ba abundances is increased by a factor of ten.
As a final exercise, we also calculate the Eu evolution in the early Galaxy using the same models and assumptions we have employed for the prediction of Ba evolution. Figure 3b shows the color-coded frequency distribution of stars in the \[Eu/Mg\]$``$\[Mg/H\] plane, after convolving with $`\sigma =0.2`$ dex for \[Eu/Mg\] and $`\sigma =0.1`$ dex for \[Mg/H\]. The Eu production site must be the same as for Ba, or at least share a part of SNe II sites where Ba is produced, because the rapid $`n`$-capture process cannot synthesize Eu without producing Ba. Therefore, we assume that Eu is also produced by SNe II with $`M_{\mathrm{ms}}=2025M_{}`$. We scale the Ba yield by the pure $`r`$-process value \[Ba/Eu\]=$`0.72`$ (Wisshak, Voss, & Käpeler 1996) and obtain a predicted Eu yield of $`M_{\mathrm{Eu}}=1.3\times 10^6M_{}`$ for $`M_{\mathrm{ms}}=20M_{}`$, and $`M_{\mathrm{Eu}}=7.0\times 10^9M_{}`$ for $`M_{\mathrm{ms}}=25M_{}`$. Our calculation predicts the existence of a Eu $`i`$-branch at \[Mg/H\]$`\stackrel{<}{}2.5`$, which must be confirmed by future measurements of Eu abundances for extremely metal-poor stars.<sup>2</sup><sup>2</sup>2 In TSY, we used the data of Ryan et al. (1996) for determining the Eu yield. These observations exhibit a much lower \[Eu/Mg\] ratio than that of McWilliam et al. (1995). As a consequence, the entire population of SNe II with $`M_{\mathrm{ms}}=1250M_{}`$ may be the Eu production site.
## 4 CONCLUSION
We have suggested the existence of two distinct correlations of \[Ba/Mg\] with \[Mg/H\] for metal-poor stars in the Galaxy, the $`y`$\- and $`i`$-branches, separated at \[Mg/H\] $`2.7`$. These branches cross one another almost perpendicularly, and form an arrow-like frequency distribution of stars in the \[Ba/Mg\]$``$\[Mg/H\] plane. The vertical boundary to the $`y`$-branch extending from \[Ba/Mg\]$`=2.0`$ to $`+0.6`$ reflects the different nucleosynthesis contributions of \[Ba/Mg\] for each SN II of progenitor mass in a narrow range, $`M_{\mathrm{ms}}=2025M_{}`$. The horizontal $`i`$-branch extending from \[Mg/H\]$`=2.7`$ to $`3.7`$ is populated by stars born from the shells swept up by SNe II with $`M_{\mathrm{ms}}=1220M_{}`$ that do not synthesize the $`r`$-process elements. Therefore, the Ba abundances of the $`i`$-branch stars reflect those in the ISM at the time when such stars form.
The $`i`$-branch provides strong evidence for the existence of a lower mass limit of the SN progenitors for the site of $`r`$-process nucleosynthesis. The upper mass limit for the $`r`$-process site cannot be determined well from the current observations. Because of the declining IMF, SNe II with $`M_{\mathrm{ms}}\stackrel{>}{}25M_{}`$ contributes little to the total Ba production in the early Galaxy. In other words, in our formulation, the $`r`$-process site is dominated by SNe II with $`M_{\mathrm{ms}}2025M_{}`$.
The $`i`$-branch should not exist for $`\alpha `$\- and Fe-peak heavy elements formed in the early Galaxy, because the entire range of SNe II progenitor masses will contribute to substantial amounts of such elements. However, the $`i`$-branch is expected to exist, though is not confirmed as of yet, in the case of the $`r`$-process elements such as La, Ce, Nd, and Eu. It is therefore important to measure the abundances of such elements for the stars CS 22878$``$101, CS 22897$``$008, CS 22891$``$200, CS 22891$``$209, CS 22948$``$066, CS 22952$``$015, and CS 22950$``$046 in the sample of McWilliam et al. (1995), which actually form the Ba $`i`$-branch.
We are grateful to Timothy C. Beers for his careful reading and helpful comments on our Letter. This work has been partially supported by COE research (07CE2002) and a Grant-in-Aid for Scientific Research (11640229) of the Ministry of Education, Science, Culture, and Sports in Japan. |
warning/0001/math0001164.html | ar5iv | text | # Bernstein–Gelfand–Gelfand sequences
## 1. Introduction
Our approach to geometries modeled on homogeneous spaces goes back to E. Cartan’s notion of an ‘espace generalisé’. The central objects for such geometries are suitably normalized Cartan connections in the sense commonly adopted, see e.g. . The models for the geometries considered in this paper are homogeneous spaces of the type $`G/P`$, where $`G`$ is real or complex semisimple and $`PG`$ is a parabolic subgroup. In this case, there is a close link to the project of parabolic invariant theory suggested by Ch. Fefferman in and in view of this context we talk about the (real and complex) parabolic geometries.
We explore the semi–holonomic jet modules and we use implicitly the cohomological information given by Kostant’s version of the Bott–Borel–Weil theorem in order to construct sequences of homomorphisms between jet–modules, which in turn give rise to sequences of invariant differential operators expressed in terms of the invariant derivatives with respect to Cartan connections, on all (curved) geometries in question. These sequences are differential complexes if certain twisted de Rham sequences are complexes, and then they compute the same cohomology. In particular, this always happens for the homogeneous models themselves and then our sequences specialize to the Bernstein–Gelfand–Gelfand resolutions well known from representation theory for complex $`G/P`$, while their real smooth analogues are provided for all real forms of this situation.
In spite of the fact that we have mentioned a few concepts from representation theory, we want to underline that no deeper aspects of representation theory are used in the construction of our new sequences of invariant operators and in the discussion of their basic properties. In particular, no infinite dimensional representation theory is needed. It is rather the language and the way of thinking of representation theory that is essential (in a similar way as the categorical language is useful in mathematics even if no deep results of category theory are used). In order to stress this feature, we have postponed the more detailed analysis of the structure of the sequences to a forthcoming second part of the article and we hope that the first part is accessible for differential geometers without a deeper background in representation theory. We also provide a quite detailed exposition of the necessary algebraic background. In particular we have included two appendices covering some material which is rather well known in representation theory.
The first general geometric theory close to our needs had been worked out in the series of papers by N. Tanaka and his school aiming at the original equivalence problem of E. Cartan, see and the references therein. Our inspiration comes, however, rather from the interest in the links between twistor theory and representation theory, as explained in the book . In the generality we need, the normalized Cartan connections were constructed in first. We have been also influenced by the translation principle in representation theory (see for example) and, in particular, by some ideas in the second part of Baston’s paper . Some arguments and proofs in the latter paper seem very unclear to us, however.
There are also many treatments of specific examples of parabolic geometries in the literature, including e.g. projective, conformal, almost Grassmannian, and CR–geometries. Most of these well known geometries correspond to the so called $`|1|`$–graded Lie algebras $`𝔤`$ which can be equivalently expressed by the requirement that the tangent spaces correspond to irreducible representations of the parabolic subgroup $`P`$. Our theory of semi–holonomic jet–modules is in fact a generalization of the approach worked out for all real $`|1|`$–graded algebras in our former papers (and this paper could be also viewed as a fourth part of this series expanded to the full generality of parabolic geometries). On the other hand, there are only few explicit examples of curved analogues of the Bernstein–Gelfand–Gelfand resolutions available in the literature, see e.g. , and in fact only the case of conformal Riemannian geometries has been studied systematically, see and for two different approaches. For an introduction addressed to wide audience, see the forthcoming paper .
Let us indicate the structure of the paper. In the next section, we first collect the necessary information on $`|k|`$-graded Lie algebras and the structure of the corresponding Lie groups, and then real and complex parabolic geometries are introduced (cf. 2.7). Our point of view is that the geometry on a manifold $`M`$ is given by a *choice* of a Cartan connection (with possible further normalization) and we are interested in the general calculus which such a choice offers. In a certain sense, this is similar to the rôle of the general calculus for linear connections in Riemannian geometry by application to the Levi–Civita connection. Thus we only briefly discuss the more classical underlying geometrical information on the manifolds $`M`$ themselves and the question of constructing a (normalized) Cartan connection from these more basic data, cf. 2.10. See for more information on this aspect. We also introduce the concepts of natural bundles and operators for parabolic geometries in the end of Section 2.
The third section deals with our basic algebraic tool, the semi–holonomic jet modules. The *invariant derivative* with respect to Cartan connections then leads to the notion of *strongly invariant differential operators* which are defined by means of $`P`$–module homomorphisms. As a first application, we introduce the twisted exterior derivatives which are certain torsion adjusted versions of the covariant exterior derivatives induced by the Cartan connections on certain bundles.
The main results are stated and proved in Section 4. Referring implicitly to the structure of the Lie algebra cohomologies, we first embed the natural vector bundles corresponding to cohomologies into exterior forms by means of distinguished differential operators $`L`$, see Theorem 4.8. Then we use the twisted exterior derivatives in order to construct explicitly many $`P`$–module homomorphisms of the semi–holonomic jet modules, cf. Proposition 4.9. The corresponding invariant differential operators build the *Bernstein–Gelfand–Gelfand sequences*. Finally we discuss the conditions under which these sequences form differential complexes, and we discuss their cohomologies, cf. 4.134.15.
Finally, we illustrate briefly the achievements on at least one non–trivial parabolic geometry and this is done in Section 5.
Throughout the paper, we discuss the real and complex manifolds and groups at the same time. We should point out however, that the relation between the real and complex settings deserves more attention. In fact, we are able to present both smooth and holomorphic results in one line of arguments, because our point is to use the $`P`$–module homomorphisms in order to construct the sequences of operators. The distinction is hidden in the explicit structure of the Lie algebra cohomologies, which we use only implicitly. One should say, however, this does not mean that working out the details for one real form gives explicit results for all other real or complex forms of the group in question. This ambiguity disappears only if we restrict ourselves to complex representations of the real forms.
A more detailed discussion of our Bernstein–Gelfand–Gelfand sequences requires a deeper study of the cohomological information. Essentially, the non–trivial operators between the irreducible bundles in the sequence correspond to arrows in the Hasse diagram of the parabolic subalgebras and the knowledge of this structure leads to quite explicit information on the individual operators. We have preferred to postpone all considerations which need more involved information from representation theory to a prospective continuation in order to keep the flavor of this article.
Acknowledgements. The research evolved during a stay of the first two authors at the University of Adelaide supported by the Australian Research Council, and during the meetings of all three authors at the Erwin Schrödinger Institute for Mathematical Physics in Vienna, the Masaryk University in Brno, and the Charles University in Prague. The institutional support by GAČR, Grant Nr. 201/99/0675 has been essential, too. Our particular thanks are due to Michael Eastwood who explained to us several aspects of the Bernstein–Gelfand–Gelfand resolutions.
## 2. Parabolic geometries
In this section we review basic facts about $`|k|`$–graded Lie algebras and we give basic definitions on parabolic geometries and invariant differential operators on manifolds equipped with geometries of that type. Most of the facts on the algebras go back to , see also which is fully compatible in notation.
### 2.1. Definition
Let $`𝕂`$ be $``$ or $``$. A $`|k|`$–graded Lie algebra over $`𝕂`$, $`k`$ is a Lie algebra $`𝔤`$ over $`𝕂`$ together with a decomposition
$$𝔤=𝔤_k\mathrm{}𝔤_1𝔤_0𝔤_1\mathrm{}𝔤_k$$
such that $`[𝔤_i,𝔤_j]𝔤_{i+j}`$ and such that the subalgebra $`𝔤_{}:=𝔤_k\mathrm{}𝔤_1`$ is generated by $`𝔤_1`$. In the whole paper, we will only deal with semisimple $`|k|`$–graded Lie algebras.
By $`𝔭`$ we will denote the subalgebra $`𝔤_0\mathrm{}𝔤_k`$ of $`𝔤`$, and by $`𝔭_+`$ the subalgebra $`𝔤_1\mathrm{}𝔤_k`$ of $`𝔭`$.
There is always a unique element $`E𝔤`$ whose adjoint action is given by $`[E,X]=\mathrm{}X`$ for $`X𝔤_{\mathrm{}}`$. The element $`E`$ is contained in the center of the subalgebra $`𝔤_0`$, which is always reductive. Using this, one shows that any ideal of $`𝔤`$ is homogeneous. Thus, a semisimple $`|k|`$–graded Lie algebra is always a direct sum of simple $`|k_i|`$–graded Lie algebras, where all $`k_ik`$. Hence, one usually can reduce most discussions to the simple case. When dealing with the semisimple case, we have to assume that none of the simple factors is contained in $`𝔤_0`$, for technical reasons. Since basically we are interested in homogeneous spaces $`G/P`$, where $`G`$ is a Lie group with Lie algebra $`𝔤`$ and $`P`$ an appropriate subgroup with Lie algebra $`𝔭`$, and their curved analogs, this is not really a restriction.
For each $`i=1,\mathrm{},k`$, the Killing form of $`𝔤`$ induces an isomorphism $`𝔤_i𝔤_i^{}`$ of $`𝔤_0`$–modules. Finally, the powers of $`𝔭_+`$ are given by $`𝔭_+^i=𝔤_i\mathrm{}𝔤_k`$, for $`i=1,\mathrm{},k`$. See e.g. \[35, Section 3\] for details.
### 2.2.
In the complex case, the meaning of a $`|k|`$–grading is particularly simple to describe. One can show that there always exists a Cartan subalgebra $`𝔥𝔤`$ which contains the element $`E`$ from above, and a choice of positive roots $`\mathrm{\Delta }_+`$ for $`𝔥`$ such that all root spaces corresponding to simple roots are either contained in $`𝔤_0`$ or in $`𝔤_1`$. Denoting by $`\mathrm{\Sigma }`$ the set of those simple roots, whose root spaces are contained in $`𝔤_1`$, one sees that the grading on $`𝔤`$ is given by the $`\mathrm{\Sigma }`$–height of roots. That is, if $`\alpha `$ is a root, then the root space $`𝔤_\alpha `$ is contained in $`𝔤_i`$, where $`i`$ is the sum of all coefficients of elements of $`\mathrm{\Sigma }`$ in the expansion of $`\alpha `$ as a linear combination of simple roots. In particular, this implies that the subalgebra $`𝔭`$ is always a parabolic subalgebra of $`𝔤`$, and $`𝔭=𝔤_0𝔭_+`$ is exactly the Levi decomposition of $`𝔭`$ into a reductive and a nilpotent part.
Conversely, if $`𝔤`$ is complex and semisimple and $`𝔭𝔤`$ is a parabolic subalgebra, then one can find a Cartan subalgebra and a set of positive roots such that $`𝔭`$ is the standard parabolic corresponding to a set $`\mathrm{\Sigma }`$ of simple roots. But then the $`\mathrm{\Sigma }`$–height as defined above gives a $`|k|`$–grading on $`𝔤`$, where $`k`$ is the $`\mathrm{\Sigma }`$–height of the maximal root of $`𝔤`$, such that $`𝔭=𝔤_0\mathrm{}𝔤_k`$. See e.g. \[22, p. 88\] or \[2, Section 2\] for more details.
Thus, in the complex case giving a $`|k|`$–grading on $`𝔤`$ is the same thing as giving a parabolic subalgebra $`𝔭`$ of $`𝔤`$. Therefore, complex $`|k|`$–graded semisimple Lie algebras can be conveniently denoted by Dynkin diagrams with crossed nodes. That is, given a $`|k|`$–graded semisimple complex Lie algebra we may assume that $`𝔭`$ is the standard parabolic subalgebra corresponding to a set $`\mathrm{\Sigma }`$ of simple roots. Then we denote the $`|k|`$–graded Lie algebra $`𝔤`$ by crossing out the nodes corresponding to the simple roots contained in $`\mathrm{\Sigma }`$ in the Dynkin diagram of $`𝔤`$. See the book for a detailed discussion of the Dynkin diagram notation for parabolic subalgebras.
Finally note that for a $`|k|`$–graded Lie algebra $`𝔤`$ over $``$ the complexification $`𝔤^{}`$ of $`𝔤`$ is $`|k|`$–graded, too. So in general we deal with certain real forms of pairs $`(𝔤,𝔭)`$, where $`𝔤`$ is complex and semisimple and $`𝔭`$ is a parabolic in $`𝔤`$. The classification of all these real forms is provided in \[35, Section 4\].
### 2.3.
Suppose that $`𝔤`$ is $`|k|`$–graded and semisimple over $`𝕂=`$ or $``$, and let $`G`$ be any Lie group with Lie algebra $`𝔤`$. (We do not assume that $`G`$ is connected.) Then we can define subgroups $`G_0PG`$ as follows: $`G_0`$ consists of all elements of $`G`$ such that the adjoint action $`\mathrm{Ad}(g):𝔤𝔤`$ of $`g`$ preserves the grading of $`𝔤`$. By $`P`$ we denote the subgroup of all elements $`gG`$ such that $`\mathrm{Ad}(g)`$ preserves the filtration by right ends induced by the grading of $`𝔤`$, i.e. $`\mathrm{Ad}(g)(𝔤_i)𝔤_i\mathrm{}𝔤_k`$. By definition $`G_0`$ is a subgroup of $`P`$, and one easily verifies that $`G_0`$ and $`P`$ have Lie algebras $`𝔤_0`$ and $`𝔭`$, respectively, see e.g. \[7, 2.9\]. Moreover, it can be shown that if $`𝔤`$ is simple, then $`P`$ equals the normalizer $`N_G(𝔭)`$ of $`𝔭`$ in $`G`$, so it is the usual parabolic subgroup associated to the parabolic subalgebra $`𝔭`$.
The following proposition clarifies the structure of the group $`P`$:
###### Proposition.
Let $`gP`$ be any element. Then there exist unique elements $`g_0G_0`$ and $`X_i𝔤_i`$ for $`i=1,\mathrm{},k`$, such that
$$g=g_0\mathrm{exp}(X_1)\mathrm{}\mathrm{exp}(X_k).$$
###### Proof.
See \[7, 2.10\]. ∎
### 2.4.
For $`i=1,\mathrm{},k`$ we define a subgroup $`P_+^iP`$ as the image under the exponential map of $`𝔤_i\mathrm{}𝔤_k`$, and we write $`P_+`$ for $`P_+^1`$. Then we have $`PP_+P_+^2\mathrm{}P_+^k`$. The subgroup $`P_+P`$ is obviously normal and by Proposition 2.3 we have $`P/P_+G_0`$, so $`P`$ is the semidirect product of $`G_0`$ and the normal nilpotent subgroup $`P_+`$. More generally, for each $`i>1`$ we see that $`P/P_+^i`$ is the semidirect product of $`G_0`$ and the normal nilpotent subgroup $`P_+/P_+^i`$.
The adjoint action of $`P`$ on $`𝔤`$ by definition preserves any of the subspace $`𝔤_i\mathrm{}𝔤_k`$ for $`i=k,\mathrm{},k`$. Thus for each $`i=k,\mathrm{},k`$ and $`j>i`$ we get an induced action of $`P`$ on the quotient $`𝔤_i\mathrm{}𝔤_k/(𝔤_j\mathrm{}𝔤_k)`$. With a slight abuse of notation, we will denote this $`P`$–module by $`𝔤_i\mathrm{}𝔤_{j1}`$. Again by Proposition 2.3, the action of $`P_+^{ji}`$ on $`𝔤_i\mathrm{}𝔤_{j1}`$ is trivial, so the action of $`P`$ on this space is induced by an action of $`P/P_+^{ji}`$. In particular, we get an action of $`P`$ on $`𝔤_{}=𝔤/𝔭`$, which is induced by an action of $`P/P_+^k`$.
There is another important consequence of Proposition 2.3: Suppose that $`𝕍`$ and $`𝕎`$ are $`P`$–modules and that $`\mathrm{\Phi }:𝕍𝕎`$ is a linear mapping. Suppose that $`\mathrm{\Phi }`$ is equivariant for the action of $`G_0`$ and for the (infinitesimal) action of $`𝔤_1`$. Since $`𝔭_+`$ is generated by $`𝔤_1`$ this implies equivariancy with respect to $`𝔭_+`$ and thus also with respect to $`P_+`$, so using Proposition 2.3 we see that $`\mathrm{\Phi }`$ is actually a homomorphism of $`P`$–modules. This will be technically very important in the sequel.
### 2.5.
For a Lie group $`G`$ with $`|k|`$–graded semisimple Lie algebra $`𝔤`$ and the subgroup $`P`$ defined in 2.3 above, consider the homogeneous space $`G/P`$. This homogeneous space is the flat model for the parabolic geometry of the type $`(G,P)`$ that we are going to study. It is well known that the canonical projection $`GG/P`$ is a principal fiber bundle with structure group $`P`$.
If $`G`$ is a complex Lie group, then $`P`$ is a parabolic subgroup, so $`G/P`$ is a generalized flag manifold, and thus in particular a compact complex manifold. In the real case, $`G/P`$ need not be compact in general, as the example of the conformal spheres in indefinite signature shows.
Next suppose that $`\lambda :P\text{GL}(𝕍)`$ is a representation of $`P`$ on a finite dimensional vector space $`𝕍`$. Then we can form the associated bundle $`V:=G\times _P𝕍G/P`$. This is a *homogeneous vector bundle*, that is the canonical left action of $`G`$ on $`G/P`$ lifts to a left action of $`G`$ on $`V`$ by vector bundle homomorphisms. Conversely, given a homogeneous vector bundle $`EG/P`$, consider the fiber $`𝔼`$ of $`E`$ over the canonical base point $`oG/P`$. Since the action of any element of $`P`$ on $`G/P`$ maps $`o`$ to itself, the action of $`G`$ on $`E`$ induces a representation of $`P`$ on $`𝔼`$ and one easily verifies that $`G\times _P𝔼`$ and $`E`$ are isomorphic homogeneous vector bundles (i.e. there is a $`G`$–equivariant isomorphism of vector bundles between them). Consequently, there is a bijective correspondence between finite dimensional representations of $`P`$ and homogeneous vector bundles over $`G/P`$. In the case where $`G`$ is a complex Lie group, the action of $`G`$ on $`G/P`$ is holomorphic and there is a bijective correspondence between holomorphic finite dimensional representations of $`P`$ and holomorphic homogeneous vector bundles over $`G/P`$ (that is holomorphic bundles with holomorphic $`G`$–actions).
In particular, the tangent and cotangent bundles of $`G/P`$ are homogeneous vector bundles. One easily verifies that they correspond to the representations of $`P`$ on $`𝔤_{}𝔤/𝔭`$ and $`𝔭_+`$ induced by the adjoint action, respectively. In the complex case, these representations induce the holomorphic tangent and cotangent bundle.
For a homogeneous vector bundle $`EG/P`$ consider the space $`\mathrm{\Gamma }(E)`$ of smooth sections of $`E`$. There is an induced action of $`G`$ on this space given by $`(gs)(x)=g(s(g^1x))`$ for $`xG/P`$. In the complex case, we can deal similarly with the spaces of holomorphic sections.
### Definition
Let $`E`$ and $`F`$ be homogeneous vector bundles over $`G/P`$. A (linear) invariant differential operator $`D:\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ is a linear differential operator $`D`$ which is equivariant for the $`G`$–actions constructed above.
### 2.6.
If $`D`$ is of order $`r`$, then it is induced by a vector bundle homomorphism $`\stackrel{~}{D}:J^r(E)F`$, where $`J^r(E)`$ is the $`r`$–th jet prolongation of $`E`$. Now simply by functoriality of the $`r`$–th jet prolongation, $`J^r(E)`$ is again a homogeneous vector bundle, and the invariance of $`D`$ is equivalent to the fact that $`\stackrel{~}{D}`$ is equivariant for the $`G`$–actions on $`J^r(E)`$ and $`F`$. Since $`G`$ acts transitively on $`G/P`$, the homomorphism $`\stackrel{~}{D}`$ is actually determined by its restriction $`\stackrel{~}{D}:J^r(E)_oF_o`$ to the fiber over $`oG/P`$, and by invariance of $`D`$, this map is $`P`$–equivariant.
Conversely, a $`P`$–homomorphism $`J^r(E)_oF_o`$ extends uniquely to a $`G`$–homomorphism $`J^r(E)F`$ and thus gives rise to an invariant differential operator. Thus, invariant differential operators $`\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ of order $`r`$ are in bijective correspondence with $`P`$–homomorphisms $`J^r(E)_oF_o`$. To avoid the restriction on the order, one can simply pass to infinite jets and we obtain that invariant differential operators $`\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ are in bijective correspondence with $`P`$–homomorphisms $`J^{\mathrm{}}(E)_oF_o`$, which factorize over some $`J^r(E)`$.
Surprisingly, the problem of determining all such homomorphisms has a nice reformulation in term of (infinite–dimensional) representation theory, which has led to a complete solution in several cases. Namely, suppose that $`E`$ and $`F`$ correspond to representations $`𝔼`$ and $`𝔽`$ of $`P`$, respectively. For the dual representation $`𝔼^{}`$, one can form the induced module $`𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$, which is a $`(𝔤,P)`$–module, i.e. it admits compatible actions of $`𝔤`$ and $`P`$. In the case where $`𝔭𝔤`$ is the Borel subalgebra (i.e. the minimal parabolic) and $`𝔼`$ is irreducible, these are the Verma–modules while for general $`𝔭`$ and irreducible $`𝔼`$, they are called generalized Verma–modules. By a dualization argument and Frobenius reciprocity one shows that for $`𝔼`$ and $`𝔽`$ irreducible, the space of all $`P`$–module homomorphisms $`J^{\mathrm{}}(E)_oF_o`$, which factorize over some $`J^r(E)_o`$ is isomorphic to the space of all $`(𝔤,P)`$–homomorphisms $`𝒰(𝔤)_{𝒰(𝔭)}𝔽^{}𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$. Since these considerations are essential for understanding of the links of our development to the standard Bernstein–Gelfand–Gelfand resolutions, we provide some more details in Appendix Appendix A.
Let us remark however that while there is a complete classification of homomorphisms of Verma–modules in the complex case in , the classification of homomorphisms of generalized Verma modules is a very difficult problem, which is unsolved in general (even in the complex case). There is a complete classification in the case of real rank one for one dimensional representations in and for general representations in and . The problem in the case of generalized Verma modules is the following: One has a class of homomorphisms which are induced by homomorphisms of Verma modules, the so–called standard homomorphisms. These are exactly the homomorphisms which occur in Bernstein–Gelfand–Gelfand resolutions. But it may happen that a homomorphism of Verma modules induces the zero–homomorphism between generalized Verma modules, and in this situation there may still be nonzero homomorphisms (the so called non–standard homomorphisms).
### 2.7. Parabolic geometries
Some geometries can be viewed as curved analogs of the homogeneous spaces $`G/P`$ considered above. For the purpose of this paper, the best way to define them is simply as generalized spaces in the sense of E. Cartan.
Let $`𝔤=𝔤_k\mathrm{}𝔤_k`$ be a real $`|k|`$–graded Lie algebra and let $`G`$ be a Lie group with Lie algebra $`𝔤`$. Let $`G_0`$ and $`P`$ be the subgroups of $`G`$ defined in 2.3 above. Then we define a (real) parabolic geometry of type $`(G,P)`$ on a smooth manifold $`M`$ to be a principal $`P`$–bundle $`𝒢M`$ equipped with a Cartan connection of type $`(G,P)`$, i.e. a differential form $`\omega \mathrm{\Omega }^1(𝒢,𝔤)`$ such that
1. $`\omega (\zeta _X)=X`$ for all $`X𝔭`$
2. $`(r^b)^{}\omega =\text{Ad}(b^1)\omega `$ for all $`bP`$
3. $`\omega |_{T_u𝒢}:T_u𝒢𝔤`$ is a linear isomorphism for all $`u𝒢`$.
Here $`\zeta _X`$ denotes the fundamental vector field generated by $`X𝔭`$ and $`r^b`$ denotes the principal right action of $`bP`$. Thus, $`\omega `$ gives a smooth $`P`$–equivariant trivialization of the tangent bundle of $`𝒢`$, which reproduces the generators of fundamental fields. Each $`X𝔤`$ defines the *constant vector field* $`\omega ^1(X)`$ given by $`\omega ^1(X)(u)=\omega _u^1(X)T_u𝒢`$. Clearly, a parabolic geometry of type $`(G,P)`$ on $`M`$ can only exist if $`M`$ has the same dimension as $`G/P`$.
In the complex setting, the Lie algebras and groups, as well as the manifold $`M`$ are complex and the above definition remains unchanged except for the replacement of smooth by holomorphic. Thus a complex parabolic geometry of the type $`(G,P)`$ on a complex manifold $`M`$ is given by a holomorphic principal fiber bundle equipped with a holomorphic absolute parallelism $`\omega `$ with the three properties listed above.
The (real or complex) homogeneous space $`G/P`$ always carries a canonical parabolic geometry, namely $`𝒢=G`$ and the Cartan connection is given by the left Maurer Cartan form. Then the constant vector fields are exactly the left invariant fields on $`G`$.
It is fairly easy to make the parabolic geometries as defined above into a category. Let $`(𝒢,\omega )`$ be a real parabolic geometry on $`M`$ and $`(𝒢^{},\omega ^{})`$ be a parabolic geometry on $`M^{}`$, and suppose that $`\mathrm{\Phi }:𝒢𝒢^{}`$ is a smooth homomorphism of principal $`P`$–bundles, such that the induced map $`\underset{¯}{\mathrm{\Phi }}:MM^{}`$ is a local diffeomorphism. Then for any point $`u𝒢`$ the tangent map $`T_u\mathrm{\Phi }:T_u𝒢T_{\mathrm{\Phi }(u)}𝒢^{}`$ is a linear isomorphism, and using this, one immediately verifies that $`\mathrm{\Phi }^{}\omega ^{}:=\omega ^{}T\mathrm{\Phi }`$ is a Cartan connection on $`𝒢`$. Now we define a morphism from $`(𝒢,\omega )`$ to $`(𝒢^{},\omega ^{})`$ to be a homomorphism $`\mathrm{\Phi }`$ of principal bundles such that the induced map $`\underset{¯}{\mathrm{\Phi }}:MM^{}`$ is a local diffeomorphism and such that $`\omega =\mathrm{\Phi }^{}\omega ^{}`$. For complex parabolic geometries we additionally require all maps to be holomorphic.
Note that any homomorphism $`\mathrm{\Phi }:𝒢𝒢^{}`$ of principal bundles which lies over a local diffeomorphism can be viewed as a morphism $`(𝒢,\mathrm{\Phi }^{}\omega ^{})(𝒢^{},\omega ^{})`$. More generally, if $`(𝒢^{},\omega ^{})`$ is a parabolic geometry on $`M^{}`$ and $`f:MM^{}`$ is a local diffeomorphism, then we can form the pullback bundle $`f^{}𝒢^{}M`$. Then there is an induced homomorphism $`\mathrm{\Phi }:f^{}𝒢^{}𝒢^{}`$ of principal bundles which lies over $`f`$, and we get an induced morphism $`(f^{}𝒢^{},\mathrm{\Phi }^{}\omega ^{})(𝒢^{},\omega ^{})`$.
### 2.8.
For some purposes, the category of parabolic geometries as defined above is too large, and one has to impose certain restrictions. Usually, these restrictions are on the curvature of the Cartan connection. Initially, the curvature of a Cartan connection $`\omega `$ is defined as the $`𝔤`$–valued two–form $`K\mathrm{\Omega }^2(𝒢,𝔤)`$ defined by the structure equation
$$K(\xi ,\eta )=d\omega (\xi ,\eta )+[\omega (\xi ),\omega (\eta )],$$
where $`\xi `$ and $`\eta `$ are vector fields on $`𝒢`$ and the bracket is in $`𝔤`$. Using the properties of $`\omega `$ one immediately verifies that $`K`$ is horizontal and equivariant. In particular, this implies that $`K`$ is uniquely determined by the curvature–function $`\kappa :𝒢\mathrm{\Lambda }^2𝔤_{}^{}𝔤`$ defined by $`\kappa (u)(X,Y)=K(u)(\omega _u^1(X),\omega _u^1(Y))`$. There are two natural ways to split $`\kappa `$ into components. First, the splitting of $`𝔤`$ induces a splitting of $`\kappa `$ according to the values in $`𝔤`$. In particular, we can split $`\kappa =\kappa _{}\kappa _𝔭`$ according to the splitting $`𝔤=𝔤_{}𝔭`$. Following the classical terminology for affine connections, $`\kappa _{}`$ is called the torsion of $`\omega `$. The other possibility is to split $`\kappa `$ according to homogeneous components. We denote the homogeneous component of degree $`i`$ of $`\kappa `$ by $`\kappa ^{(i)}`$. So $`\kappa ^{(i)}`$ maps $`𝔤_j𝔤_k`$ to $`𝔤_{i+j+k}`$.
Another important point is that the space $`\mathrm{\Lambda }^2𝔤_{}^{}𝔤`$ is the second chain group $`C^2(𝔤_{},𝔤)`$ in the standard complex for the Lie algebra cohomology $`H^{}(𝔤_{},𝔤)`$ of the nilpotent Lie algebra $`𝔤_{}`$ with coefficients in the $`𝔤_{}`$–module $`𝔤`$. As we shall recall in detail in Section 4, there is the adjoint $`^{}`$ to the Lie algebra differential $``$ in this complex, so in particular, we have $`^{}:\mathrm{\Lambda }^2𝔤_{}^{}𝔤𝔤_{}^{}𝔤`$.
### Definition
Let $`(𝒢,\omega )`$ be a (real or complex) parabolic geometry on a manifold $`M`$, and let $`\kappa `$ be the curvature of $`\omega `$. Then the parabolic geometry is called
1. normal if $`^{}\kappa =0`$.
2. regular if it is normal and $`\kappa ^{(i)}=0`$ for all $`i0`$.
3. torsion–free if $`\kappa _{}=0`$.
4. flat if $`\kappa =0`$.
Note that forming the curvature of a Cartan connection is a natural operation. This means that if $`\mathrm{\Phi }:𝒢𝒢^{}`$ is a homomorphism of principal bundles and $`\omega ^{}`$ is a Cartan connection with curvature $`K^{}`$ and curvature–function $`\kappa ^{}`$ then the curvature $`K`$ and curvature function $`\kappa `$ of the pullback $`\mathrm{\Phi }^{}\omega ^{}`$ are given by $`K=\mathrm{\Phi }^{}K^{}`$ and $`\kappa =\kappa ^{}\mathrm{\Phi }`$, respectively. Since all the subclasses of parabolic geometries defined above are given by restricting the values of the curvature–function, morphisms into a parabolic geometry from one of the four subclasses can only come from geometries from the same subclass. Clearly, for any of the four subclasses the geometries belonging to the class form a full subcategory of the category of all parabolic geometries of fixed type.
### 2.9. Examples
Before we review the construction of parabolic geometries from underlying data, we present two well known examples.
### Conformal structures
Consider $`^n`$ with coordinates $`x_1,\mathrm{},x_n`$ and the standard inner product $`,`$ of signature $`(p,q)`$, and $`^{n+2}`$ with coordinates $`x_0,x_1,\mathrm{},x_n,x_{\mathrm{}}`$ and the inner product associated to the quadratic form $`2x_0x_{\mathrm{}}+(x_1,\mathrm{},x_n),(x_1,\mathrm{},x_n)`$, which has signature $`(p+1,q+1)`$. Let $`G=SO_0(p+1,q+1)`$ be the connected component of the special orthogonal group of this metric. Then the Lie algebra $`𝔤`$ of $`G`$ admits a $`|1|`$–grading by decomposing matrices into blocks of sizes $`1`$, $`n`$, and $`1`$, see e.g. \[8, 3.3(2)\]. The construction of the canonical Cartan connection for manifolds endowed with a conformal structure of signature $`(p,q)`$, originally due to E. Cartan (see ), shows that conformal structures of this signature are precisely the same thing as normal parabolic geometries corresponding to that choice of $`G`$ and $`P`$. See for a construction of the canonical Cartan connection on conformal manifolds in a style similar to the approach of this paper. In this special situation, normal Cartan connections turn out to be automatically regular and torsion free, so three of the four subclasses defined in 2.8 above coincide. The flat parabolic geometries in this case are exactly the locally conformally flat manifolds.
### Partially integrable almost CR–structures
The complex analog of the above construction leads to the partially integrable almost CR–structures which present another example of real parabolic geometries. Here we have to consider the complex vector space $`^n`$ with the standard Hermitian inner product of signature $`(p,q)`$ and $`^{n+2}`$ with the Hermitian inner product associated to $`z_0\overline{z}_{\mathrm{}}+\overline{z}_0z_{\mathrm{}}+(z_1,\mathrm{},z_n),(z_1,\mathrm{},z_n)`$. Now we put $`G=PSU(p+1,q+1)`$ the quotient of the special unitary group corresponding to this Hermitian inner product by its center. Splitting the matrices in the Lie algebra $`𝔤`$ of $`G`$ into blocks of sizes $`1`$, $`n`$, and $`1`$ this time gives rise to a $`|2|`$–grading. The construction of canonical Cartan connections in shows that partially integrable almost CR–structures with non–degenerate Levi–form of signature $`(p,q)`$ are exactly the same thing as regular parabolic geometries corresponding to $`G`$ (see \[7, 4.14\]). In this case, three of the four subclasses of geometries defined in 2.8 above are really different: The torsion free parabolic geometries in this case are precisely the CR–structures (see \[7, 4.16\]), and the flat ones are those which are locally isomorphic to the homogeneous model. The only coincidence in this case is that normal parabolic geometries are automatically regular.
### 2.10. Underlying structures
These two examples already show that identifying a geometrical structure on a manifold as a parabolic geometry should be rather the result of a theorem than a definition. In fact one can show in a fairly general setting that certain parabolic geometries are determined by underlying structures. This is the subject of the paper which generalizes , see also and . To review the results, we first describe the underlying structures we have in mind.
Suppose that $`(𝒢,\omega )`$ is a regular parabolic geometry on a manifold $`M`$. The first thing we get out of this is a filtration $`TM=T^kMT^{k+1}M\mathrm{}T^1M`$ of the tangent bundle of $`M`$. This is given by defining $`T^iM`$ to be the set of those tangent vectors $`\xi `$ on $`M`$ for which there is a tangent vector $`\stackrel{~}{\xi }`$ in $`T𝒢`$ lying over $`\xi `$ with $`\omega (\stackrel{~}{\xi })𝔤_i\mathrm{}𝔤_k`$. The latter condition is independent of the choice of $`\stackrel{~}{\xi }`$ since changing the vector with fixed footpoint adds a vertical vector whose image under $`\omega `$ lies in $`𝔭`$, while changing the footpoint leads to the adjoint action of an element of $`P`$, which by definition preserves the subspace $`𝔤_i\mathrm{}𝔤_k`$. Clearly, this filtration has the property that the rank of $`T^iM/T^{i+1}M`$ equals the dimension of $`𝔤_i`$ for all $`i=k,\mathrm{},1`$.
Now the underlying structures basically are given by considering the bundles $`𝒢/P_+^iM`$ for $`i=1,\mathrm{},k`$ and the “traces” of the Cartan connection that remain on these bundles. This “trace” on the bundle $`𝒢/P_+^iM`$ is a frame form of length $`i`$ in the sense of \[7, 3.2\]. For the case $`i=1`$ the geometric meaning of such a frame form is particularly easy to describe: It is exactly a reduction to the structure group $`G_0`$ of the associated graded vector bundle
$$\mathrm{gr}TM=T^kM/T^{k+1}M\mathrm{}T^2M/T^1MT^1M$$
to the tangent bundle $`TM`$. The fact that the curvature–function $`\kappa `$ of the regular Cartan connection $`\omega `$ has the property that $`\kappa ^{(i)}=0`$ for all $`i0`$ is reflected in a property of the underlying frame forms called the structure equation, see \[7, 3.4\]. The bundle $`𝒢/P_+^i`$ together with the frame form of length $`i`$, which satisfies the structure equations is called the underlying $`P`$–frame bundle of degree $`i`$. Again, for $`i=1`$ this condition can be easily understood geometrically. It is equivalent to the fact that the algebraic Lie bracket on $`\mathrm{gr}TM`$ which comes from the reduction to the group $`G_0`$ is induced by the Lie bracket of vector fields, that is it is given by a (generalized) Levi–form.
Now the main result of can be stated (with the help of the language of Dynkin diagrams for the pairs $`(𝔤,𝔭)`$ mentioned in 2.2 above) as follows:
Let $`(𝔤,𝔭)`$, $`G`$, $`P`$, and $`G_0`$ be as in 2.3 and suppose throughout that no simple factor of $`𝔤`$ is contained in $`𝔤_0`$ and $`𝔤`$ does not contain a simple factor of type $`A_1`$. Then:
(1) If $`(𝔤,𝔭)`$ does not contain any simple factor of one of the types
or
then any regular parabolic geometry can be reconstructed from the underlying $`P`$–frame bundle of degree one, and any $`P`$–frame bundle of degree one comes from a regular parabolic geometry. Thus, in all these cases regular parabolic geometries are the same thing as manifolds with filtered tangent bundle plus reductions of $`\mathrm{gr}TM`$ to the group $`G_0`$ such that the resulting algebraic bracket is induced by the Lie bracket.
(2) If $`𝔤`$ contains simple factors of one of the two above types, then any regular parabolic geometry can be reconstructed from the underlying $`P`$–frame bundle of degree two and any such bundle comes from a regular parabolic geometry. Moreover, any $`P`$–frame bundle of degree one can be extended (in various ways) to a $`P`$–frame bundle of degree two.
The classical examples of the second case are the projective structures where the $`P`$–frame bundle of degree one is simply the full frame bundle and all the structure is contained in the choice of an extension to a $`P`$–frame bundle of degree two. The other exceptional examples are the so called projective contact structures.
### 2.11. Natural bundles and operators
We will not go into much detail in the generalities about natural bundles and natural operators, but just outline the basic facts. We do not want to compare the various notions of naturality (this will be taken up elsewhere) but just show that the operators we are going to construct are natural (or invariant) in any reasonable sense.
Given a representation of $`P`$ on a vector space $`𝕍`$ and a parabolic geometry $`(𝒢M,\omega )`$ we can form the associated bundle $`VM=𝒢\times _P𝕍M`$. If $`\mathrm{\Phi }:𝒢𝒢^{}`$ is a homomorphism of principal bundles which covers a local diffeomorphism $`\underset{¯}{\mathrm{\Phi }}:MM^{}`$, then we get an induced homomorphism of vector bundles $`VMVM^{}`$ which lies over the same map $`\underset{¯}{\mathrm{\Phi }}`$ and restricts to a linear isomorphism in each fiber. To put it in another way, we get a functor from the category of parabolic geometries to the category of vector bundles over manifolds of the same dimension as $`G/P`$ and vector bundle homomorphisms which cover local diffeomorphisms and induce linear isomorphisms in each fiber such that the composition of the base functor with the given functor equals the base functor. Thus, we get a special case of a gauge natural bundle as defined in \[24, Chapter XII\].
Consider next a fixed category of real parabolic geometries, and two representations $`𝕍`$ and $`𝕎`$ of $`P`$. Let $`V`$ and $`W`$ be the corresponding natural vector bundles. A natural linear operator mapping sections of $`V`$ to sections of $`W`$ is defined to be a system of linear operators $`D_{(𝒢,\omega )}:\mathrm{\Gamma }(VM)\mathrm{\Gamma }(WM)`$, where $`M`$ is the base of $`𝒢`$ such that for any morphism $`\mathrm{\Phi }:(𝒢,\omega )(𝒢^{},\omega ^{})`$ we have
$$\mathrm{\Phi }^{}D_{(𝒢^{},\omega ^{})}=D_{(𝒢,\omega )}\mathrm{\Phi }^{}.$$
This definition implies immediately, that each of the operators is local both in the section and in the Cartan connection: Suppose that $`s\mathrm{\Gamma }(VM)`$ vanishes identically on an open subset $`UM`$. Then there is an obvious inclusion morphism $`i:(𝒢|_U,\omega |_U)(𝒢,\omega )`$ and $`i^{}s=0`$. Thus also $`i^{}(D_{(𝒢,\omega )}(s))=0`$, i.e. $`D_{(𝒢,\omega )}(s)`$ is identically zero on $`U`$. Similarly, assume that $`\omega `$ and $`\omega ^{}`$ are two Cartan connections which coincide on $`𝒢|_U`$. Then for any section $`s\mathrm{\Gamma }(VM)`$ we have $`D_{(𝒢,\omega )}(s)|_U=D_{(𝒢,\omega ^{})}(s)|_U`$. In particular, the classical Peetre theorem implies that each of the operators $`D_{(𝒢,\omega )}`$ is locally over $`M`$ a finite order differential operator with respect to the arguments in the vector bundles and the Cartan connection.
For complex parabolic geometries, we deal with holomorphic representations of $`P`$, the natural vector bundles are holomorphic, and the natural operators act on holomorphic sections. Let us also remark that all these concepts extend to non-linear objects without essential changes.
### 2.12.
The natural operators on the category of flat parabolic geometries are particularly easy to describe: It is a classical result on Cartan connections that any flat parabolic geometry is locally isomorphic to the homogeneous model $`G/P`$ (see \[7, 4.12\] for a proof in the setting of parabolic geometries). This immediately implies that any natural operator on the category of flat parabolic geometries is uniquely determined by its value on the homogeneous model $`G/P`$, i.e. the parabolic geometry $`(GG/P,\omega )`$. Moreover, an operator on the flat model extends to a natural operator on the category of flat parabolic geometries if and only if it is natural with respect to all automorphisms of $`(G,\omega )`$. The left multiplication by any element of $`G`$ induces an automorphism of the principal bundle $`GG/P`$ and by left invariance of the Maurer Cartan form this actually is an automorphism of the parabolic geometry $`(G,\omega )`$. On the other hand, by \[31, Theorem 3.5.2\] the only smooth functions $`GG`$ which pull back the Maurer Cartan form to itself are the constant left translations. Thus $`G`$ is exactly the group of all automorphisms of $`(G,\omega )`$. But this immediately implies that an operator on the homogeneous model extends to a natural operator on the category of flat parabolic geometries if and only if it is invariant in the sense of definition 2.5. Thus for the flat case, the description of natural operators is equivalent to a problem in representation theory.
Usually, the question on more general natural operators is then posed (in the special cases that have been studied so far) as the question of the existence of curved analogs of invariant operators. This should be viewed as follows: As we discussed in 2.6, an invariant operator of order $`r`$ is induced by a $`P`$–module homomorphism $`J^r(E)_oF_o`$, which does not factor over $`J^{r1}(E)_o`$. Now the kernel of the projection $`J^r(E)_oJ^{r1}(E)_o`$ is the bundle $`S^rT^{}(G/P)E`$, so it corresponds to the representation $`S^r𝔭_+𝔼`$. Thus the invariant operator gives rise to a $`P`$–module homomorphism $`S^r𝔭_+𝔼𝔽`$, which in turn gives a $`G`$–equivariant homomorphism between the corresponding homogeneous vector bundles which is precisely the symbol of the operator we started with. But this $`P`$–module homomorphism induces a homomorphism of associated bundles on any parabolic geometry, so for any parabolic geometry $`(𝒢,\omega )`$ over a manifold $`M`$, we get the corresponding homomorphism $`S^rT^{}MEMFM`$. Now a curved analog of an invariant operator is a natural operator such that for each $`(𝒢,\omega )`$ the symbol of $`D_{(𝒢,\omega )}`$ is the above homomorphism. Otherwise put, the question is whether we can extend a given natural operator from the category of flat parabolic geometries to some larger category of parabolic geometries without changing its symbol, which, as a natural transformation, makes sense on any parabolic geometry.
### 2.13.
We conclude this introductory section with some more remarks on the beautiful geometric structure underlying each parabolic geometry. This topic deserves much more attention than we could pay here and it will be studied in detail elsewhere. Some first steps have been done in .
Suppose that $`(𝒢,\omega )`$ is a real parabolic geometry on a manifold $`M`$. Then we have the tower of principal fiber bundles $`𝒢𝒢/P_+M`$ and the top level has the structure group $`P_+`$. Now using the Baker–Campbell–Hausdorff formula, Proposition 2.3 can be restated in the form that for any $`gP`$ there is a unique $`g_0G_0`$ and a unique $`Z𝔭_+`$ such that $`g=g_0\mathrm{exp}(Z)`$. But using this, one easily shows that the bundle $`𝒢𝒢/P_+`$ admits global $`G_0`$–equivariant smooth sections. Namely, one can use a local trivialization of $`𝒢M`$ to construct equivariant sections over the preimage in $`𝒢/P_+`$ of appropriate open subsets of $`M`$. Such local sections can then be glued to a global section using a partition of unity (compare with the proof of \[8, Lemma 3.6\]). As in this last reference one also proves that the space of all these sections is an affine space modeled on the space $`\mathrm{\Omega }^1(M)`$ of one–forms on $`M`$.
Each such global section $`\sigma `$ reduces the structure group of the tangent space $`TM`$ to $`G_0`$ and induces an affine connection $`\gamma ^\sigma =\sigma ^{}(\omega _𝔤_{}+\omega _{𝔤_0})`$ on $`TM`$. This affine connection is $`\sigma `$–related to another Cartan connection $`\omega ^\sigma `$ on $`𝒢`$, which differs from $`\omega `$ only in the $`𝔭_+`$–component. The class of all connections $`\gamma ^\sigma `$ is a straightforward generalization of Weyl structures on conformal geometries and all differential operators built of the Cartan connection $`\omega `$ can be expressed by uniform formulae in terms of these affine connections and their torsions and curvatures. The technique based on this general framework was developed systematically for all $`|1|`$-graded algebras $`𝔤`$ in .
## 3. Semi–holonomic jet modules and strongly invariant operators
Semi–holonomic jet prolongations of modules were first introduced in the context of AHS–structures in . Here we develop the concept in the more general setting of parabolic geometries and we discuss how the homomorphisms of semi–holonomic jet prolongations give rise to natural operators. Throughout this section, there will be essentially no differences in the arguments for the real and complex parabolic geometries. Thus we shall not mention the field of scalars explicitly, and one has to think of the proper real or complex modules in the applications below.
### 3.1. The absolutely invariant derivative
Suppose that $`(𝒢,\omega )`$ is a parabolic geometry on a manifold $`M`$. We mentioned in 2.5, that the tangent and cotangent bundles on the homogeneous spaces are homogeneous vector bundles. The Cartan connection $`\omega `$ extends this identification to all parabolic geometries as follows:
We identify $`𝔤_{}`$ (as a $`P`$–module) with $`𝔤/𝔭`$, and consider the map $`𝒢\times 𝔤_{}TM`$ defined by mapping $`(u,X)`$ to $`Tp\omega _u^1(X)`$, where $`p:𝒢M`$ is the projection. The equivariancy of the Cartan connection immediately implies that this factors to a vector bundle homomorphism $`𝒢\times _P𝔤_{}TM`$. Since this is immediately seen to be surjective, it must be an isomorphism of vector bundles by dimensional reasons. Thus we have identified $`TM`$ with the natural bundle associated to the $`P`$–module $`𝔤_{}`$. Now, the invariance of the Killing form on $`𝔤`$ implies that $`𝔤/𝔭`$ and $`𝔭_+`$ with the actions induced by the adjoint action are dual $`P`$–modules. Thus, similarly as above the cotangent bundle $`T^{}M`$ of $`M`$ can be identified with the bundle $`𝒢\times _P𝔭_+`$ (implicitly, this has been used in 2.13 above).
There is a nice way to encode the action of vector fields on functions (or equivalently the exterior derivative of functions) using the identifications made above. As we have seen, a typical tangent vector on $`M`$ can be written as $`Tp\omega _u^1(X)`$ for an element $`X𝔤_{}`$. Acting with this tangent vector on a smooth function $`fC^{\mathrm{}}(M,)`$, we get $`\omega _u^1(X)(fp)`$. Now, smooth functions on $`M`$ are in bijective correspondence with smooth $`P`$–invariant functions on $`𝒢`$, the correspondence given by mapping $`f`$ to $`fp`$. To any smooth, $`P`$–invariant function $`f`$ on $`𝒢`$ we associate a function $`^\omega f:𝒢L(𝔤_{},)`$ defined by $`^\omega f(u)(X):=\omega _u^1(X)f`$. The equivariancy properties of $`\omega `$ imply that the map $`^\omega f`$ is $`P`$–equivariant. Taking into account the above identification of $`T^{}M`$ with an associated bundle and of $`L(𝔤_{},)𝔭_+`$, we see that $`^\omega f`$ is a one form on $`M`$, which by definition coincides with $`df`$.
The above procedure immediately suggests a generalization. Let $`𝕍`$ be any representation of $`P`$ and let $`VM=𝒢\times _P𝕍`$ be the corresponding associated bundle. Then we can identify smooth sections of $`VM`$ with smooth maps $`𝒢𝕍`$, which are $`P`$–equivariant. Now to any smooth function $`s:𝒢𝕍`$ we associate a smooth function $`^\omega s:𝒢L(𝔤_{},𝕍)`$ defined by
$$^\omega s(u)(X):=\omega _u^1(X)s.$$
Obviously, this defines a differential operator
$$C^{\mathrm{}}(𝒢,𝕍)C^{\mathrm{}}(𝒢,L(𝔤_{},𝕍))$$
and these operators (for all $`(𝒢,\omega )`$) form a natural operator on all parabolic geometries in the sense of 2.11. This operation is called the universal covariant derivative in the book \[31, p. 194\]. In \[8, 2.3\] we have chosen to call it the absolutely invariant derivative. The reason for the latter name also shows the main drawback of this operation: It is not really covariant, i.e. if one starts with an equivariant map $`s`$ (i.e. a section of $`VM`$) the result is not equivariant in general. Thus in general, if we start with a section, the result of the invariant derivative is not a section of a bundle anymore.
### 3.2.
There is a way, however, to make a section of an associated bundle out of a section of an associated bundle and its absolutely invariant derivative. This is called the invariant one–jet of the section. To describe it, we first have to analyze the action of $`G`$ on one–jets in the homogeneous case. Thus, let us consider a representation $`𝕍`$ of $`P`$, the corresponding homogeneous bundle $`V(G/P)=G\times _P𝕍`$ and its first jet prolongation $`J^1(V(G/P))G/P`$. As we noted in 2.6 this is again a homogeneous bundle, and we want to describe the corresponding action of $`P`$ on its standard fiber $`𝒥^1(𝕍):=J^1(V(G/P))_o`$. As we noticed in 2.4 it suffices to understand this space as a module over $`G_0`$ and over $`𝔭_+`$ (in fact, already $`𝔤_1`$ would be sufficient).
If we think of sections in $`\mathrm{\Gamma }(V(G/P))`$ as $`P`$–equivariant functions $`sC^{\mathrm{}}(G,𝕍)^P`$, then the 1–jets of sections at the distinguished point $`oG/P`$ are identified with 1–jets of equivariant functions at the unit $`eG`$ and the action is given by $`g.(j_e^1s)=j_e^1(s\mathrm{}_{g^1})`$ for all $`gG`$. Thus, the induced action of $`Z𝔭`$ on the section $`s`$ is given by the differentiation in the direction of the right invariant vector field $`R_Z`$ on $`G`$, $`Z.j_e^1s=j_e^1(R_Zs)`$.
Now we can identify a one–jet $`j_e^1(s)`$ with $`(s(e),ds(e))`$ and as we saw in 3.1 above, $`ds(e)=^\omega s(e)`$. As a vector space we can thus write
$$𝒥^1(𝕍)=𝕍(𝔤_{}^{}𝕍)$$
and we have to understand the induced actions of $`G_0`$ and $`𝔭_+`$ on this space. Let us first assume that $`gG_0`$. Then $`(s\mathrm{}_{g^1})(e)=s(g^1)=gs(e)`$ by equivariancy of $`s`$. On the other hand, we have to evaluate $`\omega _e^1(X)(s\mathrm{}_g^1)`$. This can be computed as
$$\begin{array}{c}\frac{d}{dt}|_{t=0}s(g^1\mathrm{exp}(tX))=\frac{d}{dt}|_{t=0}s(g^1\mathrm{exp}(tX)gg^1)=\hfill \\ \hfill =\omega _e^1(\mathrm{Ad}(g^1)X)(gs)=g(\omega _e^1(\mathrm{Ad}(g^1)X)s).\end{array}$$
Now since $`gG_0`$, we have $`\mathrm{Ad}(g^1)X𝔤_{}`$ for all $`X𝔤_{}`$ (the adjoint action on $`𝔤_{}`$ coincides with the induced action on $`𝔤/𝔭`$ in this case), so we see that $`𝒥^1(𝕍)=𝕍(𝔤_{}^{}𝕍)`$ even as a $`G_0`$–module.
For $`Z𝔭_+`$ we have $`(R_Zs)(e)=Z(s(e))`$ by the infinitesimal version of equivariancy of $`s`$. On the other hand, for the derivative component we have to compute the linear mapping $`𝔤_{}X\omega ^1(X)R_Zs(e)`$. Since $`\omega ^1(X)`$ is left invariant, it commutes with $`R_Z`$ and the resulting expression depends only on $`R_Z(e)=Z=\omega ^1(Z)(e)`$, and we get
$`\omega ^1(X)R_Zs(e)`$ $`=\omega ^1(Z)\omega ^1(X)s(e)`$
$`=\omega ^1(X)\omega ^1(Z)s(e)[\omega ^1(Z),\omega ^1(X)]s(e).`$
The infinitesimal version of equivariancy of $`s`$ shows that the first term in the last expression gives $`Z(\omega _e^1(X)s(e))`$. Since $`\omega ^1(\mathrm{\_})`$ is just the left invariant vector field, the second term gives $`\omega _e^1([Z,X])s`$. Now let us split $`\mathrm{ad}(Z)=\mathrm{ad}_{}(Z)\mathrm{ad}_𝔭(Z)`$ according to the splitting $`𝔤=𝔤_{}𝔭`$. Then the $`\mathrm{ad}_𝔭(Z)(X)`$–part acts algebraically by equivariancy of $`s`$ while the rest simply produces $`\omega _e^1(\mathrm{ad}_{}(Z)(X))s`$.
Thus, if we denote elements of $`𝒥^1(𝕍)`$ as pairs $`(v,\phi )`$, where $`v𝕍`$ and $`\phi `$ is a linear map from $`𝔤_{}`$ to $`𝕍`$, then the appropriate action of $`Z𝔭_+`$ is given by
$$Z(v,\phi )=(Zv,XZ(\phi (X))\phi (\mathrm{ad}_{}(Z)(X))+\mathrm{ad}_𝔭(Z)(X)v),$$
i.e. we get the tensorial action plus one additional term mapping the value–part to the derivative–part.
This action can also be nicely written in a tensorial notation. To do this let us choose a basis $`\{\eta _\alpha \}`$ of $`𝔭_+`$ such that each element $`\eta _\alpha `$ is homogeneous of degree $`|\eta _\alpha |`$, and let $`\{\xi _\alpha \}`$ be the dual basis of $`𝔤_{}`$ (with respect to the Killing form $`B`$). Now consider an element $`(v_0,Z_1v_1)𝒥^1(𝕍)`$, where $`v_0,v_1𝕍`$ and $`Z_1𝔭_+𝔤_{}^{}`$. Then by definition $`Z_1v_1`$ maps $`X𝔤_{}`$ to $`B(Z,X)v_1`$. Thus $`[Z,X]_{}:=\mathrm{ad}_{}(Z)(X)`$ is mapped to $`B(Z_1,[Z,X]_{})v_1`$. Since the Killing form vanishes on $`𝔭_+\times 𝔭`$, this can be rewritten as $`B(Z_1,[Z,X])v_1=B([Z_1,Z],X)v_1`$. Moreover, we can write $`\mathrm{ad}_Z`$ as an element of $`L(𝔤_{},𝔤)𝔭_+𝔤`$ in the form $`_\alpha \eta _\alpha [Z,\xi _\alpha ]`$. This implies that for $`Z`$ homogeneous of degree $`|Z|`$, we may rewrite the action on $`𝒥^1𝕍`$ as
$$Z(v_0,Z_1v_1)=(Zv_0,Z_1Zv_1+[Z,Z_1]v_1+\underset{|\eta _\alpha ||Z|}{}\eta _\alpha [Z,\xi _\alpha ]v_0).$$
A simple computation shows that $`𝒥^1(\mathrm{\_})`$ can be made into a functor on the category of $`P`$–modules by defining
$$𝒥^1(f)(v,\phi ):=(f(v),f\phi )$$
for each $`P`$–module homomorphism $`f:𝕍𝕎`$.
### 3.3.
Surprisingly, the first jet prolongation of representations introduced above leads for any parabolic geometry to a natural identification of the first jet prolongation of any natural bundle with an associated bundle, i.e. with another natural bundle. Let $`(𝒢,\omega )`$ be a parabolic geometry on $`M`$, let $`𝕍`$ be a representation of $`P`$, and let $`VM`$ be the corresponding associated bundle over $`M`$.
###### Proposition.
The invariant differential $`^\omega `$ defines the mapping
$$\iota :C^{\mathrm{}}(𝒢,𝕍)^PC^{\mathrm{}}(𝒢,𝒥^1𝕍)^P,\iota (s)(u)=(s(u),(X^\omega s(u)(X)))$$
which yields an isomorphism $`J^1VM𝒢\times _P𝒥^1𝕍`$.
For each fiber bundle map $`VMWM`$ induced by a $`P`$–module homomorphism $`f:𝕍𝕎`$, the first jet prolongation of the bundle map is induced by the $`P`$–module homomorphism $`𝒥^1(f)`$.
###### Proof.
Let us recall that $`^\omega s(u)(X)=\omega ^1(X)(u)s`$. Thus the mapping $`\iota :s(s,^\omega s)`$ is well defined and depends on first jets only, so we only have to check that the values are actually equivariant. First, for $`gG_0`$ we have to compute $`(s(ug),^\omega s(ug))`$. Equivariancy of $`s`$ implies $`s(ug)=g^1(s(u))`$. The second component maps $`X𝔤_{}`$ to $`\omega _{ug}^1(X)s`$. Now the equivariancy of $`\omega `$ immediately implies that $`\omega _{ug}^1(X)=Tr^g\omega _u^1(\mathrm{Ad}(g)X)`$. Since $`gG_0`$ we see that $`\mathrm{Ad}(g)X𝔤_{}`$ and using equivariancy of $`s`$ again, we see that $`^\omega s(ug)`$ maps $`X`$ to $`g^1(\omega _u^1(\mathrm{Ad}(g)X)s)`$, and thus $`(s(ug),^\omega s(ug))=g^1(s(u),^\omega s(u))`$.
On the other hand, we have to check equivariancy for the infinitesimal action of $`Z𝔭_+`$. Thus, we have to compute $`((\zeta _Zs)(u),\zeta _Z(^\omega s)(u))`$. Equivariancy of $`s`$ implies that the first component equals $`Z(s(u))`$. The second component maps $`X𝔤_{}`$ to $`(\zeta _Z\omega ^1(X)s)(u)`$. Now $`\zeta _Z=\omega ^1(Z)`$ and we can rewrite the expression as
$$(\omega ^1(X)\omega ^1(Z)s)(u)+[\omega ^1(Z),\omega ^1(X)]s(u).$$
Since the curvature of $`\omega `$ is horizontal and $`\omega ^1(Z)`$ is vertical, we may rewrite the second term in this expression as $`(\omega ^1([Z,X])s)(u)`$. Now we can split $`[Z,X]`$ into a $`𝔤_{}`$ and a $`𝔭`$–component and conclude as in 3.2 above that $`((\zeta _Zs)(u),\zeta _Z(^\omega s)(u))=Z(s(u),^\omega s(u))`$.
Clearly, this construction gives a smooth injective homomorphism of vector bundles $`J^1VM𝒢\times _P𝒥^1𝕍`$, which covers the identity map on $`M`$. Since both bundles clearly have the same rank, this must be an isomorphism.
Finally, consider a homomorphism $`f:𝕍𝕎`$. The corresponding bundle map $`VMWM`$ is induced by $`(u,v)(u,f(v))`$, and so the induced action on sections is induced by
$$s(x(u(x),fs(u(x)))).$$
Taking 1–jet of this expression we obtain just the homomorphism $`𝒥^1(f)`$. ∎
### 3.4. Semi–holonomic jets
Since we posed no conditions on the representation $`𝕍`$ above, we can iterate the functors $`J^1`$ on the associated vector bundles as well as the functors $`𝒥^1`$ on the $`P`$–modules. Proposition 3.3 then implies that the $`r`$–th iteration $`J^1\mathrm{}J^1VM`$ is an associated bundle to $`𝒢`$ corresponding to the $`P`$–module $`𝒥^1\mathrm{}𝒥^1𝕍`$. Let us look more carefully at $`𝒥^1𝒥^1𝕍`$ and $`J^1J^1VM`$. There are two obvious $`P`$–module homomorphisms $`𝒥^1𝒥^1𝕍𝒥^1𝕍`$, the first one given by the projection $`p_{𝒥^1𝕍}`$ defined on each first jet prolongation by projection to the first component, and the other one obtained by the action of $`𝒥^1`$ on $`p_𝕍`$. Thus there is the submodule $`\overline{𝒥}^2𝕍`$ in $`𝒥^1𝒥^1𝕍`$ on which these two projections coincide. As a vector space and a $`G_0`$–module we have
$$\overline{𝒥}^2𝕍=𝕍(𝔤_{}^{}𝕍)(𝔤_{}^{}𝔤_{}^{}𝕍).$$
The two $`P`$–module homomorphisms $`𝒥^1(p_𝕍)`$ and $`p_{𝒥^1𝕍}`$ give rise to vector bundle homomorphisms $`J^1J^1VMJ^1VM`$ which are just the two standard projections on the second non–holonomic jet prolongation. So we conclude that the second semi–holonomic prolongation $`\overline{J}^2VM`$ is naturally isomorphic to $`𝒢\times _P\overline{𝒥}^2𝕍`$.
Iterating this procedure, we obtain the $`r`$–th semi–holonomic jet prolongations and $`𝒥^1(\overline{𝒥}^r𝕍)`$ equipped with two natural projections onto $`𝒥^1(\overline{𝒥}^{r1}𝕍)`$, which correspond to the usual projections on the first jet prolongation of semi–holonomic jets. Their equalizer is then the submodule $`\overline{𝒥}^{r+1}𝕍`$. As a $`G_0`$–module
$$\overline{𝒥}^r𝕍=\underset{i=0}{\overset{r}{}}(^i𝔤_{}^{}𝕍).$$
###### Proposition.
For each positive integer $`r`$, the $`r`$–th semi–holonomic jet prolongation $`\overline{J}^rVM`$ carries the natural structure of associated vector bundle $`𝒢\times _P\overline{𝒥}^r𝕍`$. Moreover, there is the natural embedding
$$J^rVM\overline{J}^rVM𝒢\times _P\overline{𝒥}^r𝕍$$
$$j^rs(u)\{u,(s(u),^\omega s(u),\mathrm{},(^\omega )^rs(u))\}.$$
###### Proof.
The first part of the statement has been already shown. What remains is to discuss the equivariancy properties of the invariant differentials. However also this follows from the first order case easily by induction, using only the definition of the semi–holonomic prolongations. ∎
### 3.5. Strongly invariant operators
The problem, why we cannot work with true (holonomic) $`r`$–jets but have to use the semi–holonomic ones, is that absolutely invariant derivatives commute only for flat Cartan connections. More precisely, from the definition of the absolutely invariant derivative and the properties of the curvature, one immediately concludes the so called general Ricci–identity
$`(^\omega ^\omega s)(u)(XYYX)`$ $`=^\omega s(u)([X,Y])+\kappa _𝔭(X,Y)(s(u))`$
$`^\omega s(u)(\kappa _{}(X,Y))`$
for all $`X,Y𝔤_{}`$. This also shows that the torsion–part of $`\kappa `$ has a quite different geometric meaning than the component valued in $`𝔭`$. Thus, the identification from proposition 3.4 has values in the $`P`$–submodule $`𝒥^r(𝕍)`$ of symmetric elements $`_{i=0}^r(S^i𝔤_{}^{}𝕍)`$ in the flat case. Consequently we have recovered the standard identification of the $`r`$–th holonomic jet prolongation of a homogeneous bundle with an associated bundle for flat geometries, but this does not work in the curved case.
Nevertheless, one can well use the semi–holonomic jet prolongations to generate invariant operators. Suppose that $`𝕍`$ and $`𝕎`$ are representations of $`P`$ and suppose that $`\mathrm{\Phi }:\overline{𝒥}^r(𝕍)𝕎`$ is a homomorphism of $`P`$–modules. Then for any parabolic geometry $`(𝒢,\omega )`$ we can define a differential operator $`\mathrm{\Gamma }(VM)\mathrm{\Gamma }(WM)`$ as follows: For a section $`s`$ viewed as an equivariant function $`𝒢𝕍`$ define
$$D_{(𝒢,\omega )}(s)(u)=\mathrm{\Phi }(s(u),^\omega s(u),\mathrm{},(^\omega )^rs(u)).$$
From Proposition 3.4 above it follows that this gives a section of the bundle $`WM`$ and that each $`D_{(𝒢,\omega )}`$ is a differential operator of order $`r`$. Moreover, by construction the operators $`D_{(𝒢,\omega )}`$ form a natural operator on the category of all parabolic geometries in the sense of 2.11. Operators arising in this way will be called strongly invariant operators in the sequel. We will often not distinguish carefully between a strongly invariant operator and the corresponding homomorphism $`\overline{𝒥}^r(𝕍)𝕎`$. Thus, the semi–holonomic jet modules give a possibility to construct natural operators for a parabolic geometry in a completely algebraic way, since one only has to construct a homomorphism between two finite dimensional $`P`$–modules.
There is a slight problem about strongly invariant operators, however. Namely, even if a homomorphism $`\overline{𝒥}^r(𝕍)𝕎`$ does not factor over $`\overline{𝒥}^{r1}(𝕍)`$, the corresponding operators may be of order strictly less than $`r`$ or even identically zero. To see this, note that we can easily compute the symbol of a strongly invariant operator. This symbol is a vector bundle homomorphism $`S^rT^{}MVMWM`$, which is induced by a homomorphism $`S^r𝔤_{}^{}𝕍𝕎`$. Using Proposition 3.4 it is clear that this homomorphism is given by restricting $`\mathrm{\Phi }`$ to $`S^r𝔤_{}^{}𝕍`$, viewed as a submodule of $`^r𝔤_{}^{}𝕍`$, which in turn can be viewed as a submodule of $`\overline{𝒥}^r(𝕍)`$. Thus, if a homomorphism restricts to zero on the symmetric part of the top component of the jet–module, then the corresponding operator actually is of lower order (and contains terms involving the curvature of the Cartan connection).
There is an important situation in which this problem does not play any role. Suppose that we have an operator of order $`r`$ in the flat case with nontrivial symbol, and suppose that we can find a homomorphism $`\overline{𝒥}^r(𝕍)𝕎`$ which induces this operator (in the flat case). Then this gives a curved analog of the operator in question, and there is no problem with the symbol at all. This will always be the case for the operators we are going to study. In particular, since $`\overline{𝒥}^1(𝕍)=𝒥^1(𝕍)`$, any first order invariant operator on the category of flat parabolic geometries is automatically strongly invariant, and thus has a canonical curved analog.
###### 3.6 Remark.
There are operators which are natural (invariant) in the sense of 2.11 but are not strongly invariant. Basically, there is only one example of such an operator known: It is shown in that on conformal manifolds of dimension $`2m`$ there exists a conformally invariant $`m`$–th power of the Laplacian on smooth functions. In it is shown that this operator is not strongly invariant. It can, however, be written in terms of absolutely invariant derivatives, and thus it is also natural. In fact, it is shown in that for AHS–structures, i.e. parabolic geometries corresponding to $`|1|`$–graded Lie algebras, naturality of (even non-linear) operators is equivalent to the possibility to express them by means of the absolute invariant derivative and curvature of the defining Cartan connection, and this, in turn, is equivalent to the existence of a universal formula in terms of all underlying affine connections, cf. 2.13.
The existence of invariant operators which are not strongly invariant is due to symmetries of the curvature of a Cartan connection. Suppose that we write an expression in terms of absolutely invariant derivatives and check whether the result is $`P`$–equivariant. Otherwise put, we can compute the obstruction against being equivariant which usually contains expressions involving the curvature of the Cartan connection and its derivatives. In the case of a strongly invariant operator, these obstructions vanish algebraically. But the jets of the curvature of any Cartan connection have certain symmetries, basically due to the Bianchi identity, see e.g. \[7, 4.9\]. This implies that expressions that do not vanish algebraically, still may vanish whenever the jet of the curvature of a Cartan connection is inserted, and this is precisely what happens in the case of the critical powers of the Laplacian.
### 3.7. Twisted invariant operators
Besides the completely reducible representations (which come from the reductive subgroup $`G_0`$) there is a second class of particularly simple representations of the group $`P`$. Namely one can take a representation of the full (semisimple) group $`G`$ and restrict it to $`P`$. These representations have particularly nice features in the case of the flat model since they give rise to trivial homogeneous bundles. There are many ways to see that, but the most appropriate one for our purposes is to associate to any element $`v`$ in a representation $`𝕍`$ of $`G`$ a global nonzero section of the associated bundle $`G\times _P𝕍`$. To do this, we just have to specify a $`P`$–equivariant map $`G𝕍`$, and we define this map simply by $`gg^1v`$. This map is even $`G`$–equivariant and not only $`P`$–equivariant.
There is a simple generalization of this result. Suppose that $`𝕎`$ is any representation of $`P`$. Then sections of $`W(G/P)`$ are in bijective correspondence with $`P`$–equivariant maps $`G𝕎`$. Now we define a map on sections of homogeneous bundles
$$\mathrm{\Gamma }(W(G/P))𝕍\mathrm{\Gamma }\left(W(G/P)V(G/P)\right)$$
$$sv(gGs(g)g^1v)$$
and one immediately verifies that this is an isomorphism of $`G`$–modules. In particular, this implies that if $`𝕎^{}`$ is another $`P`$–representation and $`D:\mathrm{\Gamma }(W(G/P))\mathrm{\Gamma }(W^{}(G/P))`$ is an invariant differential operator, then we can pull back
$$D\mathrm{id}_𝕍:\mathrm{\Gamma }(W(G/P))𝕍\mathrm{\Gamma }(W^{}(G/P))𝕍$$
along these isomorphisms to get an invariant operator
$$D_𝕍:\mathrm{\Gamma }\left(W(G/P)V(G/P)\right)\mathrm{\Gamma }\left(W^{}(G/P)V(G/P)\right).$$
This operator is called the twisted invariant operator corresponding to $`D`$ and $`𝕍`$.
Now, let us notice that the above isomorphism between the spaces of sections of the associated bundles induces a $`P`$–module isomorphism $`\overline{𝒥}^r(𝕎)𝕍\overline{𝒥}^r(𝕎𝕍)`$ for all $`P`$–modules $`𝕎`$ and $`G`$–modules $`𝕍`$ and all orders $`r`$. Thus, for strongly invariant operators $`D`$, we may extend the construction of the twisted invariant operators to natural operators $`D_𝕍`$ acting on all geometries $`(𝒢,\omega )`$ of the type $`(G,P)`$ and the resulting operators are again strongly invariant. Let us remark that a completely algebraic treatment of this construction has been worked out (in the special case of the AHS-structures) in .
In particular, we obtain the strongly invariant twisted operators $`D_𝕍`$ for all first order invariant operators $`D`$ on the homogeneous vector bundles and all $`G`$–modules $`𝕍`$.
### 3.8. Twisted exterior derivatives
The standard exterior derivatives $`d`$ on the differential forms on $`G/P`$ are first order invariant operators (since they are even invariant under the action of all diffeomorphisms of $`G/P`$), so we can apply the construction above to get the twisted exterior derivatives
$$d_𝕍:\mathrm{\Gamma }\left(\mathrm{\Lambda }^nT^{}(G/P)V(G/P)\right)\mathrm{\Gamma }\left(\mathrm{\Lambda }^{n+1}T^{}(G/P)V(G/P)\right)$$
for $`n=0,\mathrm{},\text{dim}(G/P)`$. Moreover, the operators $`d_𝕍`$ are strongly invariant, since they are of first order, and so there are the corresponding $`P`$–module homomorphisms on the semi–holonomic jet modules. Since we will need it later, we will compute these homomorphisms explicitly.
Let us start with the ordinary exterior derivative. We have already noted in 3.1 that the exterior derivative of functions equals the absolutely invariant derivative. To compute the exterior derivative for general differential forms, we first have to describe nicely the identification of $`n`$–forms with smooth equivariant functions $`G\mathrm{\Lambda }^n𝔭_+`$. Throughout, we are going to identify $`\mathrm{\Lambda }^n𝔭_+`$ with the space of $`n`$–linear alternating maps from $`𝔤_{}𝔤/𝔭`$ to $`𝕂`$. Now using the identification of the tangent bundle of $`G/P`$ with $`G\times _P𝔤_{}`$ described in 3.1, one easily verifies that the relation between a form $`\phi \mathrm{\Omega }^k(G/P)`$ and the corresponding function $`s:G\mathrm{\Lambda }^n𝔭_+`$ is given by
$$(p^{}\phi )(g)(\omega _g^1(X_1),\mathrm{},\omega _g^1(X_n))=s(g)(X_1,\mathrm{},X_n),$$
where $`p^{}\phi `$ is the pullback of $`\phi `$ along the projection $`p:GG/P`$, and the $`X_i`$ are in $`𝔤_{}`$. Note that this formula remains correct for $`X_i𝔤`$ if one interprets $`s(g)`$ as an $`n`$–linear map on $`𝔤`$ which vanishes if at least one argument lies in $`𝔭`$.
###### Lemma.
Let $`s`$ and $`ds`$ be the functions on $`G`$ corresponding to differential forms $`\phi `$ and $`d\phi `$ on $`G/P`$, respectively. Then the formula for the exterior derivative reads as
$`ds(X_0,\mathrm{},X_n)=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}(1)^i(^\omega s)(g)(X_i)(X_0,\mathrm{},\widehat{i},\mathrm{},X_n)+`$
$`{\displaystyle \underset{i<j}{}}(1)^{i+j}s(g)([X_i,X_j],X_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},X_n)`$
where $`\omega `$ is the left Maurer-Cartan form on $`G`$ and, as usual, the hat denotes omission.
###### Proof.
To compute the function corresponding to $`d\phi `$, we just have to evaluate $`p^{}(d\phi )(g)=d(p^{}\phi )(g)`$ on vector fields of the form $`\stackrel{~}{X}(g)=\omega _g^1(X)`$. We have
$`d(p^{}\phi )`$ $`(\stackrel{~}{X}_0,\mathrm{},\stackrel{~}{X}_n)={\displaystyle \underset{i=0}{\overset{n}{}}}(1)^i\stackrel{~}{X}_i(p^{}\phi )(\stackrel{~}{X}_0,\mathrm{},\widehat{i},\mathrm{}\stackrel{~}{X}_n)+`$
$`+{\displaystyle \underset{i<j}{}}(1)^{i+j}(p^{}\phi )([\stackrel{~}{X}_i,\stackrel{~}{X}_j],\stackrel{~}{X}_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},\stackrel{~}{X}_n).`$
Inserting $`p^{}\phi `$ from above and evaluating at $`g`$, we see directly that the first summand agrees with the first summand in the claimed formula, which clearly equals $`n+1`$ times the alternation of $`(^\omega s)(g)`$ evaluated at $`(X_0,\mathrm{},X_n)`$.
For the second summand, we just have to note that by the Maurer–Cartan equation for $`\omega `$ we have $`[\stackrel{~}{X}_i,\stackrel{~}{X}_j]=\stackrel{~}{[X_i,X_j]}`$. Thus, this summand gives exactly the other part of the required formula. ∎
Now let us pass to the general case of a $`V(G/P)`$–valued $`n`$–form, where $`𝕍`$ is a representation of the whole group $`G`$. Any such form can be written as a finite sum of expressions of the form $`\phi \stackrel{~}{v}`$, where $`\phi \mathrm{\Omega }^n(G/P)`$ and $`\stackrel{~}{v}`$ is the global section of $`V(G/P)`$ corresponding to $`v𝕍`$ as in 3.7 above. By definition, the twisted exterior derivative is given by $`d_𝕍(\phi \stackrel{~}{v})=(d\phi )\stackrel{~}{v}`$. Now let $`s`$ be the function corresponding to $`\phi `$ and denote by $`\stackrel{~}{v}`$ also the function corresponding to the global section. From above, we thus see that $`d_𝕍(\phi \stackrel{~}{v})`$ is represented by the function which maps $`(X_0,\mathrm{},X_n)`$ to
$$\begin{array}{c}\underset{i=0}{\overset{n}{}}(1)^i(^\omega s)(g)(X_i)(X_0,\mathrm{},\widehat{i},\mathrm{},X_n)\stackrel{~}{v}(g)+\hfill \\ \hfill +\underset{i<j}{}(1)^{i+j}s(g)([X_i,X_j],X_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},X_n)\stackrel{~}{v}(g).\end{array}$$
By definition of the absolutely invariant derivative, we have
$$^\omega (s\stackrel{~}{v})(X)=^\omega s(X)\stackrel{~}{v}+s(^\omega \stackrel{~}{v}(X))$$
and the infinitesimal version of $`G`$–invariance of $`\stackrel{~}{v}`$ says that
$$^\omega \stackrel{~}{v}(g)(X)=X(\stackrel{~}{v}(g)).$$
Thus we may rewrite the first summand in $``$ as
(1)
$$\begin{array}{c}\underset{i=0}{\overset{n}{}}(1)^i^\omega (s\stackrel{~}{v})(g)(X_i)(X_0,\mathrm{},\widehat{i},\mathrm{},X_n)+\hfill \\ \hfill +\underset{i=0}{\overset{n}{}}(1)^iX_i(s(g)(X_0,\mathrm{},\widehat{i},\mathrm{},X_n)\stackrel{~}{v}(g)).\end{array}$$
Finally note that the second term in $``$ adds up with the second term in $``$ to the value of the standard Lie algebra differential $`:C^n(𝔤_{},𝕍)=\mathrm{\Lambda }^n𝔤_{}^{}𝕍C^{n+1}(𝔤_{},𝕍)`$ (cf. 4.1 for the explicit formula) applied to the map $`s(g)\stackrel{~}{v}(g)`$ evaluated on $`(X_0,\mathrm{},X_n)`$. Thus we may summarize:
###### 3.9 Proposition.
The twisted exterior derivative $`d_𝕍`$ on $`G/P`$ is a strongly invariant operator induced by the $`P`$–module homomorphism $`\overline{𝒥}^1(\mathrm{\Lambda }^n𝔭_+𝕍)\mathrm{\Lambda }^{n+1}𝔭_+𝕍`$, which is given by the formula
$$(f_0,Zf_1)(f_0)+(n+1)Zf_1,$$
where we view elements of $`\mathrm{\Lambda }^n𝔭_+𝕍`$ as $`n`$–linear alternating maps from $`𝔤_{}`$ to $`𝕍`$ and $`Zf_1`$ denotes the alternation of the map $`(X_0,\mathrm{},X_n)B(Z,X_0)f_1(X_1,\mathrm{},X_n)`$.
###### 3.10 Corollary.
The Lie algebra differential $``$ satisfies
$$(W(f)(Wf))=(n+1)\underset{|\eta _\alpha ||W|}{}\eta _\alpha [W,\xi _\alpha ]f$$
for $`f\mathrm{\Lambda }^n𝔭_+𝕍`$ and $`W𝔭_+`$, where $`\xi _\alpha `$ and $`\eta _\alpha `$ are homogeneous dual bases of $`𝔤_{}`$ and $`𝔭_+`$ with respect to the Killing form.
###### Proof.
The claim can be verified by a nice and elementary, but tedious algebraic computation. However, the previous proposition offers the following simple argument:
We know that the formula for the strongly invariant operator
$$d_𝕍(f_0,Zf_1)=(f_0)+(n+1)Zf_1$$
is $`P`$–equivariant. Thus for all $`f_0`$, $`f_1𝕍`$, $`Z𝔭_+`$, $`W𝔭_+`$ we obtain the equality of the following two expressions
$`d_𝕍(W(f_0,`$ $`Zf_1))=d_𝕍((Wf_0,W(Zf_1)+{\displaystyle }\eta _\alpha [W,\xi _\alpha ]f_0)=`$
$`=(Wf_0)+(n+1)W(Zf_1)+(n+1){\displaystyle \eta _\alpha }[W,\xi _\alpha ]f_0`$
$`W((f_0)`$ $`+(n+1)Zf_1)=W(f_0)+(n+1)W(Zf_1).`$
This yields the required formula. ∎
### 3.11. The covariant exterior derivatives
Proposition 3.9 offers a canonical curved analog of the twisted exterior derivatives on all manifolds with a parabolic geometry of the type $`(G,P)`$. It should be remarked that we may obtain another curved analog as follows. For any parabolic geometry $`(𝒢,\omega )`$ on $`M`$, we consider the extended bundle $`\stackrel{~}{𝒢}=𝒢\times _PG`$, which is a principal $`G`$–bundle over $`M`$. It is a classical observation that the Cartan connection $`\omega `$ induces a principal connection $`\stackrel{~}{\omega }`$ on $`\stackrel{~}{𝒢}`$. Now if $`𝕍`$ is a representation of $`G`$, then we can view the corresponding natural bundle $`VM=𝒢\times _P𝕍`$ also as $`VM=\stackrel{~}{𝒢}\times _G𝕍`$, and thus we have the induced linear connection on this bundle. The covariant exterior derivative with respect to this connection gives a natural operator on $`VM`$–valued forms on $`M`$. If $`s:\stackrel{~}{𝒢}\mathrm{\Lambda }^k𝔭_+𝕍`$ is the equivariant function corresponding to a $`k`$-form $`\phi `$ on $`M`$, then the value of the latter operator is a $`(k+1)`$-form on $`M`$, given by the formula
$`d^{\stackrel{~}{\omega }}s(u)(`$ $`X_0,\mathrm{},X_n)={\displaystyle }_{i=0}^k(1)^i^{\stackrel{~}{\omega }}_{X_i}s(u)(X_0,\mathrm{},\widehat{i},\mathrm{},X_k)+`$
$`+{\displaystyle \underset{i<j}{}}(1)^{i+j}s(u)([X_i,X_j],X_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},X_k)`$
where $`X_0,\mathrm{},X_k𝔤_{}`$, $`u\stackrel{~}{𝒢}`$, $`_{X_i}^{\stackrel{~}{\omega }}s(u)`$ means the derivative of $`s`$ in the direction of the horizontal vector at $`u`$ determined by $`X_i`$, and there are the standard omissions of arguments in the expressions on the right hand side. Indeed, $`d^{\stackrel{~}{\omega }}`$ is defined as the pullback of the standard $`d`$ on $`\stackrel{~}{𝒢}`$ by the horizontal projection of $`\stackrel{~}{\omega }`$, applied to the pullback of the $`k`$-form $`\phi `$ on $`M`$ by the projection $`p:\stackrel{~}{𝒢}M`$. Since the curvature of $`\stackrel{~}{\omega }`$ produces vertical fields on $`\stackrel{~}{𝒢}`$, the above formula equals to the standard evaluation of $`d(p^{}\phi )`$ on the horizontal lifts of vector fields on $`M`$.
These operators coincide with the twisted exterior derivatives on the homogeneous space but they differ in general. The explicit general comparison is as follows:
###### Lemma.
Let $`𝕍`$ be a $`G`$-module, $`VM`$ the corresponding natural vector bundle over a manifold $`M`$ equipped with a parabolic geometry $`(𝒢,\omega )`$. The covariant exterior derivative $`d^{\stackrel{~}{\omega }}`$ on $`\mathrm{\Lambda }^kT^{}MVM`$, $`k>0`$, and the twisted exterior derivative $`d_𝕍`$ on the same space satisfy
$$d^{\stackrel{~}{\omega }}\phi =d_𝕍\phi +i_\kappa _{}\phi $$
where $`\kappa _{}`$ is the torsion–component of the curvature of $`\omega `$ and $`i_\kappa _{}\phi `$ is the usual insertion operator evaluated on $`\kappa _{}`$ and $`\phi `$, i.e. the alternation of $`\phi (\kappa _{}(X_0,X_1),X_2,\mathrm{},X_k)`$ over the arguments.
###### Proof.
The key to the required formula is in the expressions $``$ and $``$ in 3.8. Namely, the latter expressions which were derived on the homogeneous spaces describe also the twisted exterior derivatives in general, but we have to be aware that instead of the bracket $`[X_i,X_j]`$ in $``$ we have to plug in
$$\omega (u)([\omega ^1(X_i),\omega ^1(X_j)])=[X_i,X_j]\kappa (u)(X_i,X_j).$$
At the same time, for all $`u𝒢\stackrel{~}{𝒢}`$, the covariant derivative $`^{\stackrel{~}{\omega }}`$ of a section $`s:\stackrel{~}{𝒢}𝕍`$ relates to the absolute invariant derivative as
$$^{\stackrel{~}{\omega }}s(u)(X)=^\omega s(u)(X)+Xs(u)$$
(since the horizontal fields given by $`\stackrel{~}{\omega }`$ equal to $`\omega ^1(X)`$ minus the fundamental field $`\zeta _X`$).
Combining the latter two facts, we see that exactly the expression
$$\begin{array}{c}i_\kappa _{}\phi (u)(X_0,\mathrm{},X_k)=\hfill \\ \hfill \underset{i<j}{}(1)^{i+j}\phi (u)(\kappa _{}(X_i,X_j),X_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},X_k)\end{array}$$
has to be added to $`d_𝕍(u)\phi (X_0,\mathrm{},X_k)`$ in order to obtain the covariant derivative. This is exactly the evaluation of the insertion operator, cf. \[24, 8.2\]. ∎
The latter lemma shows that our twisted exterior differentials $`d_𝕍`$ are certain torsion adjusted versions of the standard covariant exterior derivatives. In particular, even in the case $`𝕍=`$ the twisted derivative $`d_{}`$ equals to the usual exterior derivative $`d`$ if and only if the geometry is torsion–free.
### 3.12. Remarks
(1) As we saw in 3.8, the isomorphism
$$\mathrm{\Gamma }(W(G/P))𝕍\mathrm{\Gamma }(W(G/P)V(G/P))$$
of $`G`$–modules induces an isomorphism of $`P`$–modules $`\overline{𝒥}^r(𝕎)𝕍\overline{𝒥}^r(𝕎𝕍)`$ for any $`P`$–representation $`𝕎`$ and $`G`$–representation $`𝕍`$. This can also be proved algebraically along the lines of . This isomorphism can then be used to define twisted versions of any strongly invariant operators in a completely algebraic way. Using this picture, the subsequent developments in this paper can be viewed as a curved analog of the Jantzen–Zuckermann translation principle in representation theory. The first version of such a curved translation procedure appeared in the context of 4–dimensional conformal geometry in , see also .
(2) The twisted exterior derivatives give a sequence
$$\mathrm{\Gamma }(VM)\mathrm{\Omega }^1(M;VM)\mathrm{}\mathrm{\Omega }^{\text{max}}(M;VM)0,$$
of invariant differential operators, where sections and forms are smooth in the real case and holomorphic in the complex case. In the case of the flat model, this sequence is just the pullback of the tensor product of the (smooth or holomorphic) de Rham sequence with $`𝕍`$, so it is a resolution of the constant sheaf $`𝕍`$. In the case of a general parabolic geometry, it fails to be a complex. Actually, it is easy to verify that the composition $`d_𝕍d_𝕍`$ is just given by the action of the curvature of $`\omega `$. Thus, in the case of a flat parabolic geometry, we still get a complex, which by Lemma 3.11 coincides with the complex given by the covariant exterior derivative with respect to the flat linear connection induced by the Cartan connection. Note however, that on a flat parabolic geometry bundles corresponding to representation of $`G`$ are no more trivial in general.
(3) As a $`G_0`$–module, one can split any representation $`𝕎`$ of $`P`$ as $`𝕎_j`$ according to eigenvalues of the grading element $`E𝔤_0`$. Clearly, the action of $`𝔭_+`$ maps $`𝔤_i𝕎_j`$ to $`𝕎_{j+i}`$. In particular, we can apply this to $`\mathrm{\Lambda }^n𝔭_+𝕍`$ to split the space $`\mathrm{\Omega }^n(M;VM)`$ into homogeneous components, and analyze how the twisted exterior derivative behaves with respect to this splitting. From the formula in Proposition 3.9 it is obvious that $`d_𝕍`$ never lowers homogeneous degree and the component of the same homogeneous degree as the input is just the Lie algebra differential $``$ composed with the given form. Thus, the homogeneous component of degree zero of $`d_𝕍`$ is algebraic and equals $``$. This observation is crucial for the subsequent development. Using the fact that the Lie algebra cohomology of $`𝔤_{}`$ with coefficients in $`𝔤`$ admits a Hodge theory (which we will discuss in the next section), we will show that we can replace the sequence of remark (2) above by a different sequence in which only sections of completely reducible bundles occur, and which is a complex computing the same cohomology if the original sequence was a complex.
## 4. Curved analogs of Bernstein–Gelfand–Gelfand resolutions
In this section, we first discuss the Hodge–structure on the standard complex for the cohomology $`H^{}(𝔤_{},𝕍)`$ for a $`𝔤`$–module $`𝕍`$. Then we come to the core of the paper, the construction of a huge class of distinguished natural operators on all parabolic geometries.
### 4.1.
We have already mentioned the standard complex for the cohomology $`H^{}(𝔤_{},𝕍)`$ in 3.8. The chain groups in this complex are the groups $`C^n(𝔤_{},𝕍)=\mathrm{\Lambda }^n𝔤_{}^{}𝕍`$, which are viewed as the spaces of $`n`$–linear alternating maps from $`𝔤_{}`$ to $`𝕍`$. The differential
$$:C^n(𝔤_{},𝕍)C^{n+1}(𝔤_{},𝕍)$$
is defined by
$`(f)(X_0,\mathrm{},X_n)=`$ $`{\displaystyle \underset{i=0}{\overset{n}{}}}(1)^iX_if(X_0,\mathrm{},\widehat{i},\mathrm{},X_n)+`$
$`+`$ $`{\displaystyle \underset{i<j}{}}(1)^{i+j}f([X_i,X_j],X_0,\mathrm{},\widehat{i},\mathrm{},\widehat{j},\mathrm{},X_n),`$
where the hats denote omission. Clearly, if we start with a representation $`𝕍`$ of the group $`G`$, then $``$ is a homomorphism of $`G_0`$–modules, and it is well known that $`=0`$.
The crucial fact for us is that on this standard complex there is a Hodge theory, which was first introduced for complex simple Lie algebras in . The most conceptual way to describe this Hodge structure is to use the natural duality between $`𝔤_{}`$ and $`𝔭_+`$ via the Killing form. This is a duality of $`G_0`$–modules, but if we consider $`𝔤_{}`$ as a $`P`$–module via the adjoint action and the identification with $`𝔤/𝔭`$, then it even is a duality of $`P`$–modules by invariance of the Killing form. Thus, given a representation $`𝕍`$ of $`𝔤`$ and its dual $`𝕍^{}`$, we can naturally identify $`C^n(𝔭_+,𝕍^{})`$ with the dual $`P`$–module of $`C^n(𝔤_{},𝕍)`$. Thus, the dual map to the Lie algebra differential $`:C^n(𝔭_+,𝕍^{})C^{n+1}(𝔭_+,𝕍^{})`$ can be viewed as a linear map
$$^{}:C^{n+1}(𝔤_{},𝕍)C^n(𝔤_{},𝕍)$$
which is called the codifferential. From the above, it is obvious that the codifferential is a $`G_0`$–homomorphism and $`^{}^{}=0`$. Moreover, one immediately verifies that the Lie algebra differential for $`𝔭_+`$ is even a $`P`$–homomorphism and thus the same is true for $`^{}`$.
A formula for $`^{}`$ can be easily computed for elements of the form $`Z_0\mathrm{}Z_nv`$, where the $`Z_i`$ are in $`𝔭_+`$ and $`v`$ is in $`𝕍`$. Pairing this element with a multilinear map $`\psi C^{n+1}(𝔭_+,𝕍^{})`$, we simply get $`\psi (Z_0,\mathrm{},Z_n)(v)`$. Using this, one immediately computes that
$`^{}(Z_0\mathrm{}Z_nv)`$ $`={\displaystyle \underset{i=0}{\overset{n}{}}}(1)^{i+1}Z_0\mathrm{}\widehat{i}\mathrm{}Z_nZ_iv+`$
$`+{\displaystyle \underset{i<j}{}}(1)^{i+j}[Z_i,Z_j]\mathrm{}\widehat{i}\mathrm{}\widehat{j}\mathrm{}Z_nv.`$
From this formula, it is again obvious that $`^{}`$ is a $`P`$–homomorphism.
Using Lie theory, one constructs an inner product on the spaces of cochains, with respect to which $``$ and $`^{}`$ are adjoint operators. The proof for this fact in the generality we need it is only a rather simple extension of results available in the literature, see e.g. . For the sake of completeness and the convenience of the reader, we give a complete proof in Appendix Appendix B.
### 4.2.
This adjointness result has a number of important consequences: First of all one gets a harmonic theory for the cohomology $`H^{}(𝔤_{},𝕍)`$. We define the Laplacian
$$\mathrm{}=^{}+^{}.$$
Then for each $`n`$ this is a $`G_0`$–endomorphism of $`C^n(𝔤_{},𝕍)`$. Moreover, the adjointness implies that $`\mathrm{ker}(\mathrm{})=\mathrm{ker}()\mathrm{ker}(^{})`$ and we have a $`G_0`$–invariant splitting
$$C^n(𝔤_{},V)=\mathrm{im}()\mathrm{ker}(\mathrm{})\mathrm{im}(^{}).$$
This implies then that the cohomology group $`H^n(𝔤_{},𝕍)`$ is isomorphic (as a $`G_0`$–module) to the subspace $`\mathrm{ker}(\mathrm{})C^n(𝔤_{},𝕍)`$. Moreover, the situation between $``$ and $`^{}`$ is completely symmetric, so we can as well compute the cohomology groups $`H^{}(𝔤_{},𝕍)`$ as $`\mathrm{ker}(^{})/\mathrm{im}(^{})`$. This is more suitable for our purposes, since, as we have noticed above, $`^{}`$ is even a $`P`$–homomorphism. This also implies that (even as a $`G_0`$–module) the cohomology group $`H^n(𝔤_{},𝕍)`$ is dual to $`H^n(𝔭_+,𝕍^{})`$.
Thus, we get a canonical action of $`P`$ on the cohomology groups $`H^n(𝔤_{},𝕍)`$. We claim, that this module is completely reducible, i.e. a direct sum of irreducibles. To prove this, we only have to show that $`𝔭_+`$ acts trivially on the cohomology groups. Fortunately, there is the following simple observation
###### Lemma.
Let $`Z𝔭_+`$ and $`fC^n(𝔤_{},𝕍)\mathrm{\Lambda }^n𝔭_+𝕍`$. Consider $`ZfC^n(𝔤_{},𝕍)`$ and $`ZfC^{n+1}(𝔤_{},𝕍)`$, as defined in 3.9. Then
$$^{}(Zf)=ZfZ^{}(f).$$
###### Proof.
This is a direct consequence of the general formula for $`^{}`$ in 4.1. ∎
Now, one immediately concludes that the $`𝔭_+`$–action maps $`\mathrm{ker}(^{})`$ to $`\mathrm{im}(^{})`$, and thus in particular the induced action on the cohomology groups is trivial.
In , B. Kostant computed the cohomology groups $`H^{}(𝔭_+,𝕍)`$ in the case when $`𝔤`$ is complex and simple and $`𝕍`$ is a complex irreducible representation. The basic idea in the proof is to analyze the action of the Laplacian $`\mathrm{}`$ in terms of Casimir operators.
In fact, our construction of the sequences of natural operators will not need the explicit knowledge of the cohomologies. On the other hand, the full power of Kostant’s theorem is necessary in order to compare the resulting sequences with the standard BGG–resolutions in representation theory.
Let us also remark here, that the knowledge of the second cohomologies with values in $`𝔤`$ determines nicely the structure of the curvature of normal parabolic geometries, see e.g. .
### 4.3. A sketch of the construction
Let us return to the twisted de Rham sequence
$$\mathrm{\Gamma }(VM)\mathrm{\Omega }^1(M;VM)\mathrm{}\mathrm{\Omega }^{\text{max}}(M;VM)0$$
on a manifold $`M`$ equipped with a parabolic geometry $`(𝒢,\omega )`$. For each $`i`$, $`\mathrm{\Omega }^i(M;VM)`$ is the space of sections (smooth in the real case, holomorphic in the complex setting) of the natural bundle associated to the representation $`\mathrm{\Lambda }^i𝔭_+𝕍`$. The maps $``$, $`^{}`$, and $`\mathrm{}`$ defined above induce maps on the spaces of sections, which we denote by the same symbols. Moreover, since these maps are induced by pointwise operations the Hodge decomposition of $`\mathrm{\Lambda }^i𝔭_+𝕍`$ gives rise to a Hodge decomposition
$$\mathrm{\Omega }^i(M;VM)=\mathrm{im}()\mathrm{ker}(\mathrm{})\mathrm{im}(^{}).$$
One has to be careful, however, that this decomposition is not $`P`$–invariant but just $`G_0`$–invariant, since $`^{}`$ is a $`P`$–homomorphism but $``$ and $`\mathrm{}`$ are not. Thus the latter decomposition makes explicit geometrical sense only after a reduction of $`𝒢`$ to $`G_0`$, cf. the discussion in 2.13.
Since $`^{}`$ is a $`P`$–homomorphism, the kernel $`\mathrm{ker}(^{})`$ and the image $`\mathrm{im}(^{})`$ are the spaces of sections of natural subbundles of $`\mathrm{\Lambda }^nT^{}MVM`$. Moreover, from 4.2 we know that the quotient $`\mathrm{ker}(^{})/\mathrm{im}(^{})`$ can be identified with the space of sections of the bundle associated to the (completely reducible) representation $`_𝕍^n=H^n(𝔤_{},𝕍)`$ of $`P`$, so we get an algebraic natural operator from the subset $`\mathrm{ker}(^{})`$ of $`\mathrm{\Omega }^n(M;VM)`$ to the space of smooth sections of the natural bundle corresponding to the representation $`_𝕍^n`$. If $`𝔼`$ is an irreducible component of $`_𝕍^n`$, then we can further project onto this component to get an algebraic natural operator $`\mathrm{ker}(^{})\mathrm{\Gamma }(EM)`$.
On the other hand, $`_𝕍^n`$ can be identified with $`\mathrm{ker}(\mathrm{})\mathrm{\Lambda }^n𝔭_+𝕍`$ as a $`G_0`$–module, so we may view any section of the corresponding bundle as a $`VM`$–valued $`n`$–form, but this is not a natural operator. The main point of the following will be that one can construct a natural differential operator $`L`$ from sections of the bundle corresponding to $`_𝕍^n`$ to $`VM`$–valued $`n`$–forms in $`\mathrm{ker}(^{})`$, which has this inclusion as the lowest homogeneous component. Otherwise put, one can split the algebraic projection $`\pi `$ constructed above by a natural differential operator $`L`$. Moreover, it will turn out that this operator is fully determined by the following surprising fact: For each section $`\alpha \mathrm{\Gamma }(H_𝕍^nM)`$ there exists the unique section $`L(\alpha )\mathrm{ker}(^{})\mathrm{\Omega }^n(M;VM)`$ such that $`\pi L(\alpha )=\alpha `$ and $`d_𝕍(L(\alpha ))\mathrm{ker}(^{})\mathrm{\Omega }^{n+1}(M;VM)`$.
Summarizing the prospective achievement, the twisted exterior derivatives will produce plenty of natural differential operators indicated by the dotted arrows in the diagram.
The idea for the construction of this natural differential operator $`L`$ is fairly simple. Recall from 3.9 that the lowest homogeneous component of $`d_𝕍`$ equals the Lie algebra differential $``$. Suppose we have a section $`s`$ in the bundle corresponding to $`_𝕍^n`$, which is homogeneous of some degree $`i`$. Then it lies in $`\mathrm{ker}(\mathrm{})`$ and thus in particular in $`\mathrm{ker}()`$, so the homogeneous component of degree $`i`$ of $`d_𝕍(s)`$ is automatically zero. The idea is now to extend $`s`$ to $`\stackrel{~}{s}`$ in such a way that $`d_𝕍(\stackrel{~}{s})`$ is as small as possible. The homogeneous component of degree $`i+1`$ of $`d_𝕍(s)`$ can be split into components in $`\mathrm{im}()`$, $`\mathrm{ker}(\mathrm{})`$, and $`\mathrm{im}(^{})`$, and the best we can do to kill it is to add to $`s`$ an element $`s_{i+1}`$ which is homogeneous of degree $`i+1`$ such that $`(s_{i+1})`$ is the negative of the $`\mathrm{im}()`$–component of $`d_𝕍(s)`$ in degree $`i+1`$. There is a freedom in the choice of $`s_{i+1}`$ which can be fixed by requiring that $`s_{i+1}\mathrm{im}(^{})`$ (which is a complement to $`\mathrm{ker}()`$ by the adjointness). But this allows us already to compute $`s_{i+1}`$: Since $`^{}(s_{i+1})=0`$ we see that $`\mathrm{}(s_{i+1})=^{}((s_{i+1}))`$. But $`(s_{i+1})`$ is just the negative of the $`\mathrm{im}()`$–part of the homogeneous component of degree $`i+1`$ of $`d_𝕍(s)`$, so this is known. Moreover, by definition $`\mathrm{}`$ commutes both with $``$ and $`^{}`$, and $`\mathrm{ker}(\mathrm{})\mathrm{im}(^{})=\{0\}`$. Thus $`\mathrm{}`$ restricts to an isomorphism $`\mathrm{im}(^{})\mathrm{im}(^{})`$. Hence we can compute $`s_{i+1}`$ by applying $`\mathrm{}^1^{}`$ to the homogeneous component of degree $`i+1`$ of $`d_𝕍(s)`$. Similarly we can continue to add an appropriate homogeneous component of degree $`i+2`$ and so on.
From this description it is not at all obvious that this construction produces a natural operator, since the map $`\mathrm{}^1`$ involved in the construction is not a $`P`$–homomorphism, and the subsequent steps of the construction use $`d_𝕍`$ which is not natural either. Below we will manage, however, to work out the procedure sketched above within the framework of homomorphisms between semi–holonomic jet modules. Thus the resulting operators $`L`$ will be even strongly invariant.
### 4.4.
Each $`P`$–module $`𝕍`$ enjoys a decomposition
$$𝕍=𝕍_{i_0}𝕍_{i_0+1}\mathrm{}𝕍_{i_0+k}$$
as a $`G_0`$–module, where the submodules $`𝕍_i`$ are distinguished by the requirement that the grading element $`E𝔤_0`$ (cf. 2.1) acts by scalar multiplication by $`i`$. The action of the elements $`Z𝔤_j`$ then maps $`𝕍_i`$ into $`𝕍_{i+j}`$ and so for each $`j=0,\mathrm{},k`$ the subspace $`𝕍^j:=𝕍_{i_0+j}\mathrm{}𝕍_{i_0+k}`$ is a $`P`$–submodule of $`𝕍`$. In particular, this decomposition of an irreducible $`G`$–module $`𝕍`$, viewed as $`P`$-module, runs from $`𝕍_k`$ to $`𝕍_k`$, where $`𝕍_k`$ is the $`P`$–submodule generated by the highest weight of $`𝕍`$.
Now, let $`𝔼_{i_0}`$ be an irreducible component of $`H^n(𝔤_{},𝕍)`$, on which the grading element acts by multiplication by $`i_0`$. Then we can view $`𝔼_{i_0}`$ as a $`G_0`$ submodule of $`\mathrm{ker}(\mathrm{})\mathrm{\Lambda }^n𝔭_+𝕍`$ and we write $`𝔼`$ for the $`P`$–submodule in $`\mathrm{\Lambda }^n𝔭_+𝕍`$ generated by $`𝔼_{i_0}`$. Let
$$𝔼=𝔼_{i_0}\mathrm{}𝔼_{i_0+r}$$
be the above $`G_0`$–module decomposition according to eigenvalues of the grading element. Then the action of $`𝔤_{\mathrm{}}`$ maps each $`𝔼_{i_0+i}`$ to $`𝔼_{i_0+i+\mathrm{}}`$. For each $`i=1,\mathrm{},r+1`$ we have the $`P`$–submodule $`𝔼^i`$ as above, so we can form the quotient $`𝔼/𝔼^i`$, which is as a $`G_0`$–module isomorphic to $`𝔼_{i_0}\mathrm{}𝔼_{i_0+i1}`$. In particular, $`𝔼/𝔼^1`$ is again the irreducible module $`𝔼_{i_0}`$ we started with but now viewed as a $`P`$–module, and $`𝔼/𝔼^{r+1}=𝔼`$.
###### Lemma.
(1) $`𝔼\mathrm{ker}(^{})`$ and $`𝔼^1\mathrm{im}(^{})`$.
(2) The Laplacian $`\mathrm{}`$ restricts to a $`G_0`$–isomorphism $`𝔼_{i_0+i}𝔼_{i_0+i}`$ for each $`i=1,\mathrm{},r`$.
###### Proof.
(1) The first part is clear, since $`\mathrm{ker}(^{})`$ is a $`P`$–submodule which by construction contains $`𝔼_{i_0}`$. Since we have already seen in Lemma 4.2 that the action of $`𝔭_+`$ maps $`\mathrm{ker}(^{})`$ to $`\mathrm{im}(^{})`$, the second part is also clear.
(2) We have already noted in 4.3 above that $`\mathrm{}`$ restricts to an automorphism on $`\mathrm{im}(^{})`$. Hence it suffices to prove that $`\mathrm{}(𝔼_{i_0+i})𝔼_{i_0+i}`$. By Corollary 3.10, we have for all $`e𝔼`$, $`Z𝔤_1`$
$$(Ze)=Z(e)(n+1)\underset{|\eta _\alpha |=1}{}\eta _\alpha [Z,\xi _\alpha ]e.$$
Applying $`^{}`$ to the first term we get $`Z\mathrm{}(e)`$.
Let us first take $`e_0𝔼_{i_0}`$, and consider $`\mathrm{}(Ze_0)=^{}((Ze_0))`$. Then the first term vanishes while each summand in the second term is contained in $`^{}(𝔤_1𝔤_0𝔼_{i_0})^{}(𝔤_1𝔼_{i_0})`$. Since $`𝔼_{i_0}\mathrm{ker}(^{})`$, Lemma 4.2 implies that $`^{}(𝔤_1𝔼_{i_0})𝔤_1𝔼_{i_0}𝔼_{i_0+1}`$. Thus, we see that $`\mathrm{}(𝔼_{i_0+1})𝔼_{i_0+1}`$. Now one can proceed inductively in the same way to show that $`\mathrm{}(𝔼_{i_0+i})𝔼_{i_0+i}`$. ∎
### 4.5.
The actual construction of the splitting operators is a little tricky. The problem is that the individual steps in the construction sketched in 4.3 are induced by maps on jet–modules which are not $`P`$–module homomorphisms but only restrict to $`P`$–module homomorphisms on appropriate submodules, which also have to be constructed during the procedure.
For $`ji0`$ we denote by $`\pi _i^j`$ the canonical projection $`𝔼/𝔼^j𝔼/𝔼^i`$, which is a homomorphism of $`P`$–modules. Clearly, $`\pi _i^i`$ is the identity and $`\pi _i^j\pi _j^k=\pi _i^k`$ for $`ijk`$. By $`p_i:𝒥^1(𝔼/𝔼^i)𝔼/𝔼^i`$ we denote the footpoint projection, which is a $`P`$–homomorphism, too. For any element $`\psi `$ in a general $`G_0`$–module, we denote by $`\psi _i`$ the component of $`\psi `$ on which the grading element $`E`$ acts by multiplication by $`i_0+i`$. Note that the mapping $`\psi \psi _i`$ is only a $`G_0`$–homomorphism and not a $`P`$–homomorphism. Finally, let us denote by $`j_i:𝔼/𝔼^i𝔼/𝔼^{i+1}`$ the $`G_0`$–homomorphism given by the inclusion $`𝔼_{i_0}\mathrm{}𝔼_{i_0+i1}𝔼_{i_0}\mathrm{}𝔼_{i_0+i}`$. Again, this is obviously not a $`P`$–homomorphism. Finally, let $`\mathrm{Alt}:𝔭_+\mathrm{\Lambda }^n𝔭_+𝕍\mathrm{\Lambda }^{n+1}𝔭_+𝕍`$ denote the alternation mapping. This is a $`P`$–homomorphism preserving homogeneous degrees.
For $`i=1,\mathrm{},r+1`$ consider now the module $`𝒥^1(𝔼/𝔼^i)`$. A typical element of this module is a pair $`(e,\psi )`$, with $`e𝔼/𝔼^i`$ and
$$\psi 𝔭_+𝔼/𝔼^i𝔭_+\mathrm{\Lambda }^n𝔭_+𝕍.$$
Now we define a mapping $`L_i:𝒥^1(𝔼/𝔼^i)𝔼/𝔼^{i+1}`$ by
$$L_i(e,\psi )=j_i(e)(n+1)\mathrm{}^1^{}((\mathrm{Alt}(\psi ))_i).$$
In particular, if $`\psi =Zf`$ for $`Z𝔭_+`$ and $`f𝔼/𝔼^i`$, then $`L_i(e,Zf)=j_i(e)(n+1)\mathrm{}^1^{}((Zf)_i)`$. Now the main technical step in the construction is the following
###### 4.6 Proposition.
The maps $`L_i:𝒥^1(𝔼/𝔼^i)𝔼/𝔼^{i+1}`$ have the following properties:
(1) $`L_i`$ is a $`G_0`$–homomorphism and $`\pi _i^{i+1}L_i=p_i`$.
(2) For $`\mathrm{\Psi }𝒥^1(𝔼/𝔼^i)`$ and $`W𝔤_1`$, we have
$$L_i(W\mathrm{\Psi })WL_i(\mathrm{\Psi })=\mathrm{}^1\left(W(\mathrm{}j_i(L_{i1}𝒥^1(\pi _{i1}^i)p_i)(\mathrm{\Psi }))\right).$$
In particular, $`L_1`$ is a $`P`$–homomorphism.
###### Proof.
(1) The fact that $`L_i`$ is a $`G_0`$–homomorphism follows immediately from the fact that $`𝒥^1(𝔼/𝔼^i)𝔼/𝔼^i(𝔭_+𝔼/𝔼^i)`$ as a $`G_0`$–module, see 3.2 and the definition of $`L_i`$. Moreover, since $`\mathrm{Alt}`$, $`^{}`$, and $`\mathrm{}`$ all preserve homogeneities, the last term in the definition of $`L_i`$ is homogeneous of degree $`i_0+i`$, so it lies in the kernel of $`\pi _i^{i+1}`$ and the second part follows.
(2) Clearly, it suffices to check this for elements $`\mathrm{\Psi }`$ of the form $`(e,Zf)`$ with $`e,f𝔼/𝔼^i`$ and $`Z𝔭_+`$. By definition of the action on jets, see 3.2, we see that $`W(e,Zf)`$ has footpoint $`We`$, while the homogeneous part of degree $`i_0+i`$ of the second component is given by
$$\underset{|\eta _\alpha |=1}{}\eta _\alpha [W,\xi _\alpha ]e_{i1}+W(Zf)_{i1}.$$
Consequently,
$`L_i(W(e,Zf))=`$ $`j_i(We)(n+1)\mathrm{}^1^{}({\displaystyle \underset{|\eta _\alpha |=1}{}}\eta _\alpha [W,\xi _\alpha ]e_{i1})`$
$`(n+1)\mathrm{}^1^{}(W(Zf)_{i1}).`$
By Corollary 3.10 the second term on the right hand side of this equation just gives
$$\mathrm{}^1^{}((We_{i1})W(e_{i1}))=We_{i1}\mathrm{}^1(W\mathrm{}(e_{i1})),$$
where we have used that $`^{}`$ is a $`P`$–homomorphism, $`e_{i1}`$ and $`We_{i1}`$ lie in the kernel of $`^{}`$, and that we are in a region where the Laplacian is invertible. On the other hand, we clearly have $`j_i(We)+We_{i1}=Wj_i(e)`$, since $`W𝔤_1`$ and $`e_{i1}`$ is the highest nonzero homogeneous component of $`e`$. Finally, we clearly have $`WL_i(e,Zf)=Wj_i(e)`$, since the rest lies in the component of maximal homogeneity, on which $`𝔭_+`$ acts trivially. Thus, we have arrived at
$$\begin{array}{c}L_i(W(e,Zf))WL_i(e,Zf)=\hfill \\ \hfill =\mathrm{}^1(W\mathrm{}(e_{i1}))(n+1)\mathrm{}^1(W^{}((Zf)_{i1})),\end{array}$$
where we have used once more the fact that $`^{}`$ is a $`P`$–homomorphism.
On the other hand, consider $`𝒥^1(\pi _{i1}^i)(e,Zf)`$. The footpoint of this element is just $`\pi _{i1}^i(e)`$, while in the jet part, the component of maximal homogeneity must coincide with $`(Zf)_{i1}`$. Consequently, we get
$$L_{i1}(𝒥^1(\pi _{i1}^i)(e,Zf))=j_{i1}(\pi _{i1}^i(e))(n+1)\mathrm{}^1^{}((Zf)_{i1}).$$
Subtracting $`e=p_i(e,Zf)`$ from this, we get
$$e_{i1}(n+1)\mathrm{}^1^{}((Zf)_{i1}),$$
and the formula follows. In the case $`i=1`$, we get $`L_1(W\mathrm{\Psi })WL_1(\mathrm{\Psi })=\mathrm{}^1(W(\mathrm{}j_1p_1)(\mathrm{\Psi }))`$, which vanishes, since $`𝔼_{i_0}\mathrm{Ker}(\mathrm{})`$. Hence, $`L_1`$ is equivariant for the action of $`𝔤_1`$ and thus a $`P`$–homomorphism by (1) and 2.4. ∎
### 4.7.
Now we inductively define subspaces $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)𝒥^1(𝔼/𝔼^i)`$ by $`\stackrel{~}{𝒥}^1(𝔼/𝔼^1)=𝒥^1(𝔼/𝔼^1)`$ and
$$\stackrel{~}{𝒥}^1(𝔼/𝔼^{i+1})=𝒥^1(\pi _i^{i+1})^1(\stackrel{~}{𝒥}^1(𝔼/𝔼^i))\mathrm{Ker}(L_i𝒥^1(\pi _i^{i+1})p_{i+1}).$$
###### Proposition.
For each $`i=1,\mathrm{},r+1`$ the space $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)𝒥^1(𝔼/𝔼^i)`$ is a $`P`$–submodule and $`L_i`$ restricts to a homomorphism $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)𝔼/𝔼^{i+1}`$ of $`P`$–modules. Moreover, for each $`k<i`$ we have
$$𝒥^1(\pi _k^i)\left(\stackrel{~}{𝒥}^1(𝔼/𝔼^i)\right)\stackrel{~}{𝒥}^1(𝔼/𝔼^k),$$
and on $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)`$ we have $`\pi _{k+1}^ip_i=L_k𝒥^1(\pi _k^i)`$.
###### Proof.
For $`i=1`$ the first two properties are satisfied by definition of $`\stackrel{~}{𝒥}^1(𝔼/𝔼^1)`$ and Proposition 4.6(2), while the last two properties are trivially satisfied. If we inductively assume that the result has been proved for $`i1`$, then $`𝒥^1(\pi _{i1}^i)^1(\stackrel{~}{𝒥}^1(𝔼/𝔼^{i1}))`$ is a $`P`$–submodule of $`𝒥^1(𝔼/𝔼^i)`$, and $`L_{i1}𝒥^1(\pi _{i1}^i)p_i`$ defines a $`P`$–homomorphism from this submodule to $`𝔼/𝔼^i`$, so $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)`$ is a $`P`$–submodule. Moreover, Proposition 4.6(2) immediately implies that the restriction of $`L_i`$ to this submodule is equivariant under the action of $`𝔤_1`$ and thus $`L_i`$ restricts to a $`P`$–homomorphism on that submodule by Proposition 4.6(1) and 2.4. Moreover, we get the last two properties in the case $`k=i1`$.
For $`k<i1`$, note first that $`\pi _k^i=\pi _k^{i1}\pi _{i1}^i`$ immediately implies that $`𝒥^1(\pi _k^i)\left(\stackrel{~}{𝒥}^1(𝔼/𝔼^i)\right)\stackrel{~}{𝒥}^1(𝔼/𝔼^k)`$ by induction. Finally, we compute
$`L_k𝒥^1(\pi _k^i)`$ $`=L_k𝒥^1(\pi _k^{i1})𝒥^1(\pi _{i1}^i)=\pi _{k+1}^{i1}p_{i1}𝒥^1(\pi _{i1}^i)=`$
$`=\pi _{k+1}^{i1}\pi _{i1}^ip_i=\pi _{k+1}^ip_i,`$
by functoriality of $`𝒥^1`$, induction, and the definition of the jet prolongation of a homomorphism. ∎
For $`k2`$ and $`i=1,\mathrm{},r+1`$ we inductively define
$$\stackrel{~}{𝒥}^k(𝔼/𝔼^i):=𝒥^1(\stackrel{~}{𝒥}^{k1}(𝔼/𝔼^i))\overline{𝒥}^k(𝔼/𝔼^i).$$
By Proposition 4.7 and 3.4 it follows inductively that $`\stackrel{~}{𝒥}^k(𝔼/𝔼^i)`$ is a $`P`$–submodule in both modules on the right hand side of the definition. For $`i=1`$, we obtain $`\stackrel{~}{𝒥}^k(𝔼/𝔼^1)=\overline{𝒥}^k(𝔼/𝔼^1)`$, so we simply get the full semiholonomic jet–module in this case. Moreover, a simple inductive argument shows for all $`\mathrm{}<k`$, and $`i`$
$$\stackrel{~}{𝒥}^k(𝔼/𝔼^i)\overline{𝒥}^{\mathrm{}}(\stackrel{~}{𝒥}^k\mathrm{}(𝔼/𝔼^i))\overline{𝒥}^k(𝔼/𝔼^i).$$
For each of the homomorphisms $`L_i:\stackrel{~}{𝒥}^1(𝔼/𝔼^i)𝔼/𝔼^{i+1}`$ we can now restrict the semiholonomic jet prolongation $`\overline{𝒥}^k(L_i)`$ to the submodule $`\stackrel{~}{𝒥}^{k+1}(𝔼/𝔼^i)\overline{𝒥}^k(\stackrel{~}{𝒥}^1(𝔼/𝔼^i))`$ to obtain a $`P`$–homomorphism
$$\overline{𝒥}^k(L_i):\stackrel{~}{𝒥}^{k+1}(𝔼/𝔼^i)\overline{𝒥}^k(𝔼/𝔼^{i+1}).$$
###### 4.8 Theorem.
Let $`𝔼_{i_0}`$ be an irreducible component in the cohomology $`_𝕍^n`$ which generates the $`P`$–submodule $`𝔼`$ in $`\mathrm{\Lambda }^n𝔭_+𝕍`$, cf. 4.3. For each $`k1`$ and $`i=1,\mathrm{},r+1`$ we have
$$\overline{𝒥}^k(L_i)\left(\stackrel{~}{𝒥}^{k+1}(𝔼/𝔼^i)\right)\stackrel{~}{𝒥}^k(𝔼/𝔼^{i+1}).$$
In particular, the composition
$$L:=L_r\overline{𝒥}^1(L_{r1})\mathrm{}\overline{𝒥}^{r1}(L_1)$$
defines a $`P`$–homomorphism $`L:\overline{𝒥}^r(𝔼/𝔼_1)𝔼`$. Since by definition $`𝔼`$ is a $`P`$–submodule of $`\mathrm{\Lambda }^n𝔭_+𝕍`$, this homomorphism induces a strongly invariant operator $`\mathrm{\Gamma }(E_{i_0}M)\mathrm{Ker}(^{})\mathrm{\Omega }^n(M;VM)`$, which splits the algebraic projection $`\mathrm{Ker}(^{})\mathrm{\Gamma }(E_{i_0}M)`$ described in 4.3.
###### Proof.
Let us first consider the case $`k=1`$. So we have to show that $`𝒥^1(L_i)\left(\stackrel{~}{𝒥}^2(𝔼/𝔼^i)\right)\stackrel{~}{𝒥}^1(𝔼/𝔼^{i+1})`$. By definition of $`\stackrel{~}{𝒥}^1(𝔼/𝔼^{i+1})`$, we first have to consider the composition $`𝒥^1(\pi _i^{i+1})𝒥^1(L_i)=𝒥^1(\pi _i^{i+1}L_i)`$. By Proposition 4.6(1), this equals $`𝒥^1(p_i)`$. Since $`\stackrel{~}{𝒥}^2(𝔼/𝔼^i)\overline{𝒥}^2(𝔼/𝔼^i)`$, this projection coincides with the restriction of the canonical projection $`\overline{𝒥}^2(𝔼/𝔼^i)𝒥^1(𝔼/𝔼^i)`$, and since $`\stackrel{~}{𝒥}^2(𝔼/𝔼^i)𝒥^1(\stackrel{~}{𝒥}^1(𝔼/𝔼^i))`$, this canonical projection has values in $`\stackrel{~}{𝒥}^1(𝔼/𝔼^i)`$. Thus, we have verified that $`𝒥^1(L_i)\left(\stackrel{~}{𝒥}^2(𝔼/𝔼^i)\right)𝒥^1(\pi _i^{i+1})^1\left(\stackrel{~}{𝒥}^1(𝔼/𝔼^i)\right)`$.
But then it also follows that $`L_i𝒥^1(\pi _i^{i+1})𝒥^1(L_i)`$ coincides with the composition of $`L_i`$ with the canonical projection $`\stackrel{~}{𝒥}^2(𝔼/𝔼^i)\stackrel{~}{𝒥}^1(𝔼/𝔼^i)`$, which by definition of the jet prolongation of a homomorphism (see 3.2) coincides with $`p_{i+1}𝒥^1(L_i)`$ and the proof in the case $`k=1`$ is complete.
The case $`k2`$ now immediately follows from the definitions by induction. Thus, also the existence of $`L`$ and the corresponding strongly invariant operator is clear. The fact that this operator splits the algebraic projection follows from the fact that by Lemma 4.4(1) this algebraic projection is induced by the canonical projection $`𝔼𝔼/𝔼^1`$ and the fact that $`\pi _i^{i+1}L_i=p_i`$ from Proposition 4.6(1). ∎
Next, we consider the composition of $`d_𝕍`$ with the operator corresponding to $`L`$. The corresponding homomorphism on jet modules can be computed as the restriction to $`\overline{𝒥}^{r+1}(𝔼/𝔼^1)`$ of $`d_𝕍𝒥^1(L)`$.
###### 4.9 Proposition.
For each irreducible $`G`$-module $`𝕍`$, and irreducible $`G_0`$–submodule $`𝔼_{i_0}_𝕍^n=H^n(𝔤_{},𝕍)`$, the composition
$$d_𝕍𝒥^1L:\overline{𝒥}^{r+1}(𝔼/𝔼^1)\mathrm{\Lambda }^{n+1}𝔭_+𝕍$$
has values in $`\mathrm{ker}^{}`$. The composition with the projection to the cohomology $`\pi _H:(\mathrm{\Lambda }^n𝔭_+𝕍)(\mathrm{ker}^{})_𝕍^{n+1}=H^{n+1}(𝔤_{},𝕍)`$ gives the $`P`$-module homomorphism
$$\pi _Hd_𝕍𝒥^1L:\overline{𝒥}^{r+1}(𝔼/𝔼^1)_𝕍^{n+1}.$$
For each $`n=0,\mathrm{},\mathrm{dim}M1`$, there is the strongly invariant differential operator
$$D^𝕍:\mathrm{\Gamma }(H_𝕍^nM)\mathrm{\Gamma }(H_𝕍^{n+1}M)$$
whose restrictions to the subbundles $`E_0M`$ are determined by the above $`P`$–module homomorphisms $`\overline{𝒥}^{r+1}(𝔼/𝔼^1)_𝕍^{n+1}`$.
###### Proof.
Consider first the map $`^{}d_𝕍:𝒥^1(𝔼)\mathrm{\Lambda }^n𝔭_+𝕍`$. By definition of $`d_𝕍`$, Lemma 4.2, and using the fact that $`𝔼\mathrm{Ker}(^{})`$, we see that this maps $`(e,Zf)𝒥^1(𝔼)`$ to $`^{}(e)+(n+1)^{}(Zf)=\mathrm{}(e)(n+1)Zf`$, so $`^{}d_𝕍:𝒥^1(𝔼)𝔼`$. Now Theorem 4.8 applied to $`𝒥^1(L_r)`$ shows, that $`𝒥^1(L)`$ has values in the submodule $`\stackrel{~}{𝒥}^1(𝔼)𝒥^1(𝔼)`$, and we claim that $`^{}d_𝕍`$ restricts to zero on that submodule.
To simplify notations, let us write $`p:𝒥^1(𝔼)𝔼`$ for the footpoint projection $`p_{r+1}`$ and $`\pi _i`$ for $`\pi _i^{r+1}`$. For $`ir+1`$ consider the $`P`$–homomorphism $`\pi _i^{}d_𝕍:𝒥^1(𝔼)𝔼/𝔼^i`$. By definition, this maps $`(e,Zf)`$ to $`\pi _i(\mathrm{}(e))+(n+1)\pi _i(^{}(Zf))`$. Since the Laplacian and $`^{}`$ both preserve homogeneous degrees, we may rewrite the first summand as $`\mathrm{}(\pi _i(e))`$ and the second summand as $`(n+1)\pi _i(^{}(Z\pi _{i1}(f)))`$.
On the other hand, consider $`𝒥^1(\pi _{i1}):𝒥^1(𝔼)𝒥^1(𝔼/𝔼^{i1})`$. This maps $`(e,Zf)`$ to $`(\pi _{i1}(e),Z\pi _{i1}(f))`$, and applying $`L_{i1}`$ to this element, we get $`j_{i1}(\pi _{i1}(e))(n+1)\mathrm{}^1^{}((Z\pi _{i1}(f))_i)`$. Finally, $`\pi _ip`$ maps $`(e,Zf)`$ to $`\pi _i(e)`$. Consequently, $`\mathrm{}(\pi _ipL_{i1}𝒥^1(\pi _{i1}))`$ maps $`(e,Zf)`$ to
$$\mathrm{}(\pi _i(e))j_{i1}(\mathrm{}(\pi _{i1}(e)))+(n+1)^{}((Z\pi _{i1}(f))_i),$$
and the last summand in this expression equals
$$(\pi _ij_{i1}\pi _{i1})((n+1)^{}(Zf)),$$
since $`^{}`$ preserves homogeneous degrees. Hence, we see that on $`𝒥^1(𝔼)`$ we get the equation
$$\pi _i^{}d_𝕍j_{i1}\pi _{i1}^{}d_𝕍=\mathrm{}\left(\pi _ipL_{i1}𝒥^1(\pi _{i1})\right).$$
In fact, this equation is exactly what we were aiming at in the motivation for the whole construction in 4.3. But on the submodule $`\stackrel{~}{𝒥}^1(𝔼)`$, the right hand side of the above formula vanishes identically by Proposition 4.7. Thus, iterated application of this formula shows that on $`\stackrel{~}{𝒥}^1(𝔼)`$ we have
$$^{}d_𝕍=\pi _{r+1}^{}d_𝕍=j_r\pi _r^{}d_𝕍=\mathrm{}=j_1\pi _1^{}d_𝕍.$$
But $`\pi _1^{}d_𝕍`$ maps $`(e,Zf)`$ to $`\mathrm{}(\pi _1(e))`$, which vanishes since $`𝔼_{i_0}`$ is contained in the kernel of the Laplacian, so we have proved $`^{}d_𝕍𝒥^1(L)=0`$. All the rest is now an immediate consequence. ∎
### 4.10. Definition
Let $`(𝒢,\omega )`$ be a (real or complex) parabolic geometry on a manifold $`M`$. The construction above has given rise to a sequence of strongly invariant operators $`D^𝕍`$
$$\text{}.$$
which is called the *Bernstein–Gelfand–Gelfand sequence* or *BGG–sequence* determined by the $`G`$-module $`𝕍`$.
All bundles in this sequence correspond to completely reducible representations of $`P`$, so they all split into direct sums of bundles corresponding to irreducible representations. Let us also remark that the construction applies to both real and complex settings. Next, we will show that in the flat case this sequence is a resolution of the constant sheaf $`𝕍`$. Since by Kostant’s version of the Bott–Borel–Weil theorem, the bundles occurring in this resolution in the complex case are exactly the same bundles as in the Bernstein–Gelfand–Gelfand resolution, we have obtained curved analogs of this resolution even in the real case.
The main step towards the proof that we often get a resolution is formulated in the next lemma for the general real curved case. For the complex analog see below.
###### 4.11 Lemma.
Let $`(𝒢,\omega )`$ be a real parabolic geometry on a manifold $`M`$ and let $`s\mathrm{\Omega }^n(M;VM)`$ be a $`VM`$–valued $`n`$–form. Then:
(1) There is an element $`t\mathrm{\Omega }^{n1}(M;VM)`$ such that $`s+d_𝕍(t)`$ lies in $`\mathrm{ker}(^{})`$.
(2) If $`s`$ and $`d_𝕍(s)`$ both lie in $`\mathrm{ker}(^{})`$, then $`s=L(\pi _H(s))`$.
(3) If $`d_𝕍^2(\mathrm{ker}(^{}))\mathrm{ker}(^{})`$, then the diagram
is commutative. In particular, $`D^𝕍D^𝕍=0`$ whenever $`d_𝕍d_𝕍=0`$.
###### Proof.
(1) Put $`𝒢_0=𝒢/P_+`$ and choose a global $`G_0`$–equivariant section $`\sigma :𝒢_0𝒢`$ as indicated in 2.13. Then we get a smooth map $`\tau :𝒢P_+`$ characterized by $`u=\sigma (p(u))\tau (u)`$ for all $`u𝒢`$, and $`u(p(u),\tau (u))`$ is a diffeomorphism $`𝒢𝒢_0\times P_+`$. Using this, we get an isomorphism (depending on $`\sigma `$) between $`\mathrm{\Omega }^n(M;VM)`$ and the space of smooth $`G_0`$–equivariant functions $`𝒢_0\mathrm{\Lambda }^n𝔭_+𝕍`$. But $`\mathrm{}^1^{}`$ is a $`G_0`$–homomorphism $`\mathrm{\Lambda }^n𝔭_+𝕍\mathrm{\Lambda }^{n1}𝔭_+𝕍`$ such that $`e(\mathrm{}^1^{}(e))\mathrm{ker}(^{})`$ for all $`e\mathrm{\Lambda }^n𝔭_+𝕍`$.
Now, let $`\underset{¯}{s}:𝒢_0\mathrm{\Lambda }^n𝔭_+𝕍`$ be the $`G_0`$–equivariant map corresponding to the lowest homogeneous component $`s_j`$ of the given $`n`$–form $`s`$ such that $`^{}(s_j)0`$. Passing from $`\mathrm{}^1^{}\underset{¯}{s}`$ back to a $`P`$–equivariant map $`t:𝒢\mathrm{\Lambda }^{n1}𝔭_+𝕍`$, we see that the homogeneous components up to degree $`j`$ of $`^{}(s+d_𝕍(t))`$ vanish on the image of $`\sigma `$ and thus on the whole $`𝒢`$ by equivariancy. Inductively, we can find an element $`t`$ with the required properties.
(2) Put $`s^{}=\pi _H(s)`$. By construction of the operators $`L`$, we know that $`L(s^{})\mathrm{ker}(^{})`$, $`\pi _H(L(s^{}))=s^{}`$, and $`d_𝕍(L(s^{}))\mathrm{ker}(^{})`$. Thus, we see that $`sL(s^{})\mathrm{im}(^{})`$ and $`d_𝕍(sL(s^{}))\mathrm{ker}(^{})`$. Let $`a_j`$ be the lowest possibly nonzero homogeneous component of $`sL(s^{})`$. Then the lowest possibly nonzero component of $`d_𝕍(sL(s^{}))`$ is $`(a_j)`$. Since $`\mathrm{ker}(^{})`$ is complementary to $`\mathrm{im}()`$ we must have $`(a_j)=0`$. But on the other hand, we know that $`a_j\mathrm{im}(^{})`$ which is complementary to $`\mathrm{ker}()`$, so we must have $`a_j=0`$.
(3) For $`s\mathrm{\Gamma }(H_𝕍^nM)`$, consider the element $`d_𝕍(L(s))\mathrm{\Omega }^{n+1}(M;VM)`$. By Proposition 4.9, this lies in $`\mathrm{ker}(^{})`$. Moreover, since $`L(s)\mathrm{ker}(^{})`$, our assumption on $`d_𝕍^2`$ implies that $`d_𝕍(d_𝕍(L(s)))\mathrm{ker}(^{})`$. Hence from (2) we get $`d_𝕍(L(s))=L(\pi _H(d_𝕍(L(s))))=L(D^𝕍(s))`$.
The last claim is obvious. ∎
###### 4.12 Lemma.
Let $`(𝒢,\omega )`$ be a complex parabolic geometry on a complex manifold $`M`$. Then the second and third assertion in Lemma 4.11 remain valid with the same assumptions, while the claim 4.11(1) holds true under the additional assumption that the holomorphic bundle $`𝒢𝒢_0`$ admits a global holomorphic $`G_0`$-equivariant section. This additional requirement is always fulfilled locally.
###### Proof.
All arguments in the proof of (2) and (3) in 4.11 are on the level of the $`P`$-modules and so they go equally through for both real and complex settings. The only difference in (1) is the argument which constructs the global section by means of the smooth partition of unity. Once we assume the existence of the global section, the rest is clear again. Now, any point in $`M`$ has an open neighborhood $`UM`$ such that both $`𝒢`$ and $`𝒢_0`$ are trivial over $`U`$. Since $`G_0\times P_+`$ and $`P`$ are diffeomorphic, and the map in one direction is obviously holomorphic, they are biholomorphic. Thus, the complex parabolic geometry $`𝒢|_UU`$ admits appropriate global holomorphic $`G_0`$–equivariant section. ∎
###### 4.13 Theorem.
Let $`(𝒢,\omega )`$ be a real parabolic geometry of the type $`(G,P)`$ on a manifold $`M`$, $`𝕍`$ be a $`G`$–module. If the twisted de Rham sequence
$$\text{}.$$
is a complex, then also the Bernstein–Gelfand–Gelfand sequence
defined in 4.10 is a complex, and they both compute the same cohomology.
The same statement is true for complex parabolic geometries $`(𝒢,\omega )`$ under the additional requirement that $`𝒢𝒢_0=𝒢/P_+`$ admits a global holomorphic $`G_0`$–equivariant section.
### Remark
In particular, the complex version of the Theorem may be reformulated as follows: If the twisted de Rham sequence induces a complex on the sheaf level, then the same is true for the Bernstein–Gelfand–Gelfand sequence. In particular, if the twisted de Rham sequence induces a resolution of $`𝕍`$, then so does the BGG–sequence.
Now, the original representation theoretical version of the (generalized) BGG–resolution follows immediately by duality. Moreover, let us notice that the global $`G_0`$–equivariant section as required in the Theorem always exists over a dense open submanifold in the homogeneous space $`G/P`$ (the so called big cell).
###### Proof.
As we saw in Lemma 4.11, the BGG–sequence forms a complex whenever the twisted de Rham does. So let us assume, we deal with complexes.
Since $`d_𝕍^2=0`$, 4.11(3) implies that $`L`$ is a morphism of the corresponding complexes, hence the mapping
$$\mathrm{\Gamma }(H_𝕍^nM)s^{}L(s^{})\mathrm{\Omega }^n(M;VM)$$
induces a morphism between the cohomologies.
Next, suppose that $`s\mathrm{\Omega }^n(M;VM)`$, $`n1`$ is such that $`d_𝕍(s)=0`$. Then by 4.11(1) we find an element $`t\mathrm{\Omega }^{n1}(M;VM)`$ such that $`s+d_𝕍(t)\mathrm{ker}(^{})`$. But then $`d_𝕍(s+d_𝕍(t))=0`$ so by 4.11(2) we know that $`s+d_𝕍(t)=L(\pi _H(s+d_𝕍(t)))`$, and thus the mapping defined above is surjective.
Finally, let us assume that $`s^{}\mathrm{\Gamma }(H_𝕍^{n1}M)`$ is such that there exists a $`t\mathrm{\Omega }^{n1}(M;VM)`$ with $`d_𝕍(t)=L(s^{})`$. Then by 4.11(1) we may without loss of generality assume that $`t\mathrm{ker}(^{})`$. But by assumption $`d_𝕍(t)=L(s^{})`$, so this is also contained in $`\mathrm{ker}(^{})`$, and hence $`t=L(\pi _H(t))`$ by 4.11(2), and thus $`L(s^{})=d_𝕍(L(\pi _H(t)))`$ and applying $`\pi _H`$ on both sides we get $`s^{}=D^𝕍(\pi _H(t))`$, and so we get an isomorphism in the cohomology groups. ∎
###### 4.14 Corollary.
Let $`(𝒢,\omega )`$ be a torsion free real parabolic geometry of type $`(G,P)`$ on $`M`$. Then the de Rham cohomology of $`M`$ with coefficients in $`𝕂=`$ or $``$ is computed by the (much smaller) complex
$$\text{}.$$
###### Proof.
The covariant exterior differential corresponding to the choice of the trivial $`P`$–module $`𝕂`$ coincides with the usual exterior differential $`d`$. According to Lemma 3.11, the exterior covariant differential coincides with our twisted exterior differential for all torsion–free geometries. Thus the statement follows from 4.13. ∎
###### 4.15 Corollary.
Let $`(𝒢,\omega )`$ be a flat real parabolic geometry. Then for any representation $`𝕍`$ of $`G`$ the BGG–sequence
is a complex, which computes the twisted de Rham cohomology of $`M`$ with coefficients in the bundle $`VM`$, which is defined as the cohomology of the complex given by the covariant exterior derivative with respect to the linear connection on $`VM`$ induced by the Cartan connection $`\omega `$, see 3.11.
The importance of this corollary lies in the fact that while flat parabolic geometries are locally isomorphic to the homogeneous model $`G/P`$, they may be very different from $`G/P`$ from a global point of view. Just keep in mind the broad variety of smooth manifolds admitting a locally conformally flat Riemannian metric. In particular, the bundle $`VM`$ is not trivial in general, so the twisted de Rham cohomology is a less trivial object than in the homogeneous case.
On the other hand, we may always consider the obvious flat parabolic geometry on the trivial $`P`$–bundle over $`^{\mathrm{dim}(G/P)}𝔤_{}`$. In this case, the twisted de Rham cohomologies are obviously zero, so Corollary 4.15 provides global resolutions of the constant sheaf $`𝕍`$ in this case. Simple instances of such sequences are of basic importance in various areas of mathematics, see for example .
### 4.16. Remark
As we have seen already, the $`P`$–modules $`_𝕍^n`$ are completely reducible and so the natural bundles $`H_𝕍^nM`$ decompose into direct sums of irreducible bundles. Consequently, also the operators $`D^𝕍`$ split into sums of operators between the irreducible natural bundles. In the case of the homogeneous bundles, the latter operators (and sometimes also their non-trivial compositions) are usually referred to as *standard invariant operators*. In particular, our construction provides a distinguished curved analog for each of those standard operators.
As we have underlined already in the introduction, no deep representation theoretical results had to be applied in the construction of the BGG–sequences and in the proof of Theorem 4.13. On the other hand, the full information of the Kostant’s version of Bott–Borel–Weil theorem on the Lie algebra cohomologies is strictly necessary in order to get more explicit information about the individual standard operators and the overall structure of the BGG–sequence in the flat case. Moreover, further non-trivial operators with curvature contributions in their symbols may appear in general.
Let us also remark that the explicit formulae for the standard operators were given in closed form in the terms of the underlying linear connections on $`M`$ in for all parabolic geometries with irreducible tangent bundles, i.e. for all cases with $`|1|`$–graded Lie algebra $`𝔤`$. We believe that the technique developed there should be extendible to the general case, too.
### 4.17. Remark
In the flat case, the twisted de Rham complex can be viewed as a filtered complex with the filtration given by homogeneous degrees. The fact that the lowest homogeneous component of $`d_𝕍`$ is just $``$ implies that the differential on the associated graded complex is exactly $``$. Associated to this filtration there is a spectral sequence which obviously converges and computes the twisted de Rham cohomology. Now from the construction of the operators $`D^𝕍`$ it is obvious that when passing to the appropriate subquotients, they induce the higher differentials in this spectral sequence. Usually, these higher differentials are only well defined on the appropriate subquotients, but due to the fact that we have a (fairly simple) Hodge structure on the associated graded complex, we can get a global definition in our setting.
## 5. Example
We shall illustrate the power of our results in the simple case of 5–dimensional partially integrable almost CR–structures, cf. Example 2.9. We believe that this simple geometry reflects many of the general features of parabolic geometries and we can still write down the whole BGG–sequences very explicitly at the same time. We hope that based on this example, the reader is able to imagine the vast amount of invariant operators which our main theorems produce for all parabolic geometries.
Let $`M`$ be a smooth manifold of odd dimension $`2n+1`$ together with a distinguished rank $`n`$ complex subbundle $`T^{CR}M`$ of the tangent bundle $`TM`$. Then the Lie bracket of vector fields induces a skew–symmetric bundle map $`^{}:T^{CR}M\times T^{CR}MTM/T^{CR}M`$, the real Levi–Form. $`(M,T^{CR}M)`$ is called a partially integrable almost CR–manifold if and only if $``$ is non–degenerate and totally real, i.e. $`(J(\xi ),J(\eta ))=(\xi ,\eta )`$ for all $`\xi ,\eta T^{CR}M`$, where $`J`$ denotes the almost complex structure on $`T^{CR}M`$. In that case, choosing a local trivialization of $`TM/T^{CR}M`$, $``$ is the imaginary part of a Hermitian form. Here we consider the case where $`n=2`$, so $`M`$ has dimension 5 and this Hermitian form is positive definite (for an appropriate choice of the local trivialization).
Typical examples of such manifolds are smooth hypersurfaces in a six–dimensional smooth manifold $`N`$ endowed with an almost complex structure $`\stackrel{~}{J}`$, which satisfy a non–degeneracy and an integrability condition. In this case, we put $`T^{CR}M=TM\stackrel{~}{J}(TM)`$ and $`J=\stackrel{~}{J}|_{T^{CR}M}`$. To understand the non–degeneracy and integrability conditions, it is more convenient to pass to complexified tangent bundles. Since $`T^{CR}M`$ is a complex bundle, its complexification $`T_{}^{CR}M`$ splits into a direct sum $`T_{1,0}MT_{0,1}M`$ of a holomorphic and an antiholomorphic part. Moreover, mapping $`\xi ,\eta \mathrm{\Gamma }(T_{1,0}M)`$ to the class of $`i[\xi ,\overline{\eta }]`$ defines a bundle valued Hermitian form $`:T_{1,0}M\times T_{1,0}MT_{}M/T_{}^{CR}M=:QM`$, the Levi form. The partial integrability condition from above is then equivalent to the fact that $`[\xi ,\eta ]\mathrm{\Gamma }(T_{}^{CR}M)`$ for all sections $`\xi ,\eta `$ of $`T_{1,0}M`$, and the conditions of positive definiteness is equivalent to $``$ being positive definite in an appropriate local trivialization of $`QM`$. (Certainly, these conditions also make sense for abstract almost CR manifolds). A partially integrable almost CR manifold is called integrable or a CR–manifold if and only if the subbundle $`T_{1,0}M`$ is involutive. In particular, this is the case for hypersurfaces in complex manifolds.
By \[7, 4.14\], 5–dimensional partially integrable almost CR–manifolds are exactly the normal parabolic geometries corresponding to $`G=PSU(3,1)`$ and the parabolic subalgebra of $`𝔤=𝔰𝔲(3,1)`$ corresponding to the Dynkin diagram . Let us also consider $`\stackrel{~}{G}=SU(3,1)`$ and let $`P`$, $`G_0G`$, or $`\stackrel{~}{P}`$, $`\stackrel{~}{G}_0\stackrel{~}{G}`$ be the corresponding subgroups as in 2.3. Then the semisimple part of $`\stackrel{~}{G}_0`$ is $`SU(2)`$ and the center of $`G_0`$ is $``$.
In the Dynkin diagram notation, each (complex) irreducible $`\stackrel{~}{G}`$-module $`𝕍`$ is given by the choice of three non–negative integers $`a,b,c`$
$$𝕍=\text{}.$$
More explicitly, is the highest weight component in $`S^a^4S^b(\mathrm{\Lambda }^2^4)S^c(^4)`$, where $`S^i`$ denotes the $`i`$–th symmetric power, and so these representations integrate to representation of $`G`$ if and only if $`ac`$ is congruent to $`2b`$ modulo four (the center of $`\stackrel{~}{G}`$ consist of $`\pm 1`$ and $`\pm i`$ times the identity).
The irreducible $`\stackrel{~}{P}`$–modules correspond to choices with $`b`$ non–negative while $`a`$ and $`c`$ may be arbitrary integers. Now, $`b`$ determines the representation of $`SU(2)`$ while the other two parameters describe the action of the center of $`\stackrel{~}{G}_0`$. We adopt the convention used in , i.e. the parameters give the linear combination of the fundamental weights of $`\stackrel{~}{𝔤}`$ which is the highest weight of the dual module to $`𝕍`$. In this way, the resulting weights for our modules happen to be the same as those in the dual pictures known from representation theory. For our purposes, however, this has no importance and it is enough to say that the distinguished two subbundles $`T_{1,0}M`$ and $`T_{0,1}M`$ in the complexified tangent space and the complexified quotient $`QM=T_{}M/T_{}^{CR}M`$ have duals $`T_{1,0}^{}M`$, $`T_{0,1}^{}M`$ (quotients of the complexified cotangent bundle), and $`Q^{}M`$, which correspond to the modules
$$𝕋_{1,0}^{}=\text{},𝕋_{0,1}^{}=\text{},^{}=\text{}.$$
Now, all $`\stackrel{~}{P}`$–modules are tensor products of symmetric powers $`S^b(𝕋_{1,0}^{})`$ and suitable one-dimensional representations $`𝔼[a,c]`$ corresponding to the Dynkin diagram . We shall write $`S^b(𝕋_{1,0}^{})[a,c]`$ for these modules and use the shorthand $`𝒮_{[a,c]}^b`$ for the corresponding natural bundles. In particular,
$$𝒮_{[a,c]}^b=S^b(𝕋_{1,0}^{})[a,c]=\text{}$$
$$T_{0,1}^{}=T_{1,0}^{}[2,2]=𝒮_{[2,2]}^1$$
$$𝒮_{[1,1]}^0=E[1,1]=Q^{}$$
$$𝒮_{[4,0]}^0=\mathrm{\Lambda }^2T_{1,0}^{}Q^{}.$$
Another important bundle is the dual to the kernel of the bilinear Levi form $`(\mathrm{ker})^{}T_{1,0}^{}T_{0,1}^{}`$ which corresponds to $`𝒮_{[2,2]}^2`$.
All natural bundles $`𝒮_{[a,c]}^b`$ exist on manifolds $`M`$ with the so called *$`SU(3,1)`$–structures*, i.e. we have to choose coverings of the Cartan $`P`$–bundle $`𝒢`$ to the structure group $`\stackrel{~}{P}`$. This is clearly equivalent to the choice of a fixed line bundle $`E[1,0]`$ such that its fourth tensor power is $`\mathrm{\Lambda }^2T_{1,0}MQM`$. This is an analogous situation to natural bundles and natural operators in conformal Riemannian geometry which often depend on the choice of a spin structure.
Using the explicit description of the cohomology from Kostant’s Bott–Borel–Weil theorem we obtain explicitly all natural bundles appearing in our main theorems. The computations are done fairly simply in terms of the Dynkin diagram notation, see for the details. Furthermore, using elementary finite dimensional representation theory one easily shows that there are no homomorphisms between the semi–holonomic jet modules corresponding to the items in the neighboring columns of the BGG–sequences, except those which are indicated in Figure 1. Let us also notice that the orders of the operators are easily read off the homogeneities of the bundles with respect to the action of the grading element in $`G_0`$ and the homogeneity of $`𝒮_{[a,c]}^b`$ is $`a+cb`$. Thus we can summarize:
###### 5.1 Theorem.
For each $`SU(3,1)`$–module $`𝕍=\text{}`$, the BGG–sequence of invariant differential operators shown on Figure 1 exists on all 5–dimensional partially integrable almost CR–manifolds $`M`$ with a chosen $`SU(3,1)`$–structure. The orders of the operators are indicated by the labels over the arrows. Moreover, the sequence exists on all partially integrable CR–manifolds if $`a2b+c0\mathrm{mod}4`$, and then all bundles in question can be constructed from $`T_{1,0}^{}M`$ and $`Q^{}M`$. If $`M`$ is flat, then the BGG–sequence is a complex which computes the twisted de Rham cohomology of $`M`$ with coefficients in the bundle $`VM`$ corresponding to $`𝕍`$.
As a corollary, we immediately obtain
###### 5.2 Theorem.
For all (integrable) 5–dimensional CR–manifolds, there is the resolution of the sheaf of constant complex functions
which computes the de Rham cohomology with complex coefficients. The orders of the operators in the column in the middle of the diagram are two, while all the other ones are of first order.
This complex is a special instance of the so called Rumin complex on contact geometries, , see also for a refined version for the CR–structures. In the homogeneous case, this complex was also mentioned in . Similar questions were also studied by Lychagin earlier, see e.g. and the references therein. Notice that the dimensions of the individual columns are 1, 4, 5, 5, 4, 1 (opposed to dimensions 1, 5, 10, 10, 5, 1 in the de Rham complex).
## Appendix A. Infinite jets and Verma modules
The aim of this appendix is to provide differential geometers with basic information on the links between jets and Verma modules, and in particular to prove the correspondence between invariant differential operators and homomorphisms of generalized Verma modules used in 2.6.
### A.1.
We have seen in 2.6 that invariant operators $`\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ between homogeneous vector bundles over $`G/P`$ are in bijective correspondence with $`P`$–homomorphisms $`J^{\mathrm{}}(E)_oF_o`$, which factorize over some $`J^r(E)_o`$.
First note that sections of $`E`$ can be identified with smooth functions $`G𝔼`$, which are $`P`$–equivariant. Since this identification is purely algebraic, it gives an identification of infinite jets at $`o`$ of sections of $`E`$ with $`P`$–equivariant infinite jets of smooth functions $`G𝔼`$ at $`eG`$. Now it is easy to verify that in the picture of smooth equivariant functions, the action of $`G`$ is given by $`(gs)(g^{})=s(g^1g^{})`$. The corresponding infinitesimal action of $`𝔤`$ is given by $`(Xs)(g)=(R_Xs)(g)`$, where $`R_X`$ denotes the right–invariant vector field on $`G`$ generated by $`X𝔤=T_eG`$. For $`X𝔭`$, the infinitesimal version of equivariancy of $`s`$ implies that $`(Xs)(g)=X(s(g))`$, but for general $`X`$ the value $`(Xs)(g)`$ depends on the one–jet of $`s`$ at $`g`$. Thus we do not get an induced action of $`𝔤`$ on finite jets, but for infinite jets we get a well defined action of $`𝔤`$. Since this action is clearly compatible with the action of $`P`$, it makes $`J^{\mathrm{}}(E)_o`$ into a $`(𝔤,P)`$–module.
On the other hand, mapping each $`X𝔤`$ to the left invariant vector field $`L_X`$ generated by $`X`$ induces an isomorphism between the universal enveloping algebra $`𝒰(𝔤)`$ and the algebra of left invariant differential operators on $`G`$. Now we get a bilinear map $`J^{\mathrm{}}(E)_o\times (𝒰(𝔤)𝔼^{})𝕂`$ by mapping $`(j^{\mathrm{}}s(e),D\lambda )`$ to $`\lambda (D(s)(e))`$, where $`D`$ is a left invariant differential operator and $`\lambda `$ is an element of the dual representation $`𝔼^{}`$ to $`𝔼`$, and as above we view $`s`$ as an equivariant function on $`G`$. By equivariancy of $`s`$ this factors to a bilinear map $`J^{\mathrm{}}(E)_o\times (𝒰(𝔤)_{𝒰(𝔭)}𝔼^{})𝕂`$ because elements of $`𝒰(𝔭)`$ act algebraically and this can be expressed as an action on $`\lambda `$.
We claim that the above pairing is compatible with the actions of both $`𝔤`$ and $`P`$. For the action of $`𝔤`$, let us take a typical element $`X_1\mathrm{}X_n\lambda 𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$ and $`X𝔤`$. From above, we know that $`Xj^{\mathrm{}}s(e)=j^{\mathrm{}}(R_Xs)(e)`$. Pairing this with $`X_1\mathrm{}X_n\lambda `$, we get $`\lambda ((L_{X_1}\mathrm{}L_{X_n}R_Xs)(e))`$. Since left invariant vector fields commute with right invariant ones, this equals $`\lambda ((R_XL_{X_1}\mathrm{}L_{X_n}s)(e))`$. But this depends only on $`R_X(e)`$, so we may as well replace $`R_X`$ by $`L_X`$, so this coincides with $`XX_1\mathrm{}X_n\lambda `$ evaluated on $`j^{\mathrm{}}s(e)`$.
The action of $`bP`$ on $`𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$ is induced by mapping $`D\lambda `$ to $`bDb\lambda `$, where $`(bD)(s)=D(sr^{b^1})r^b`$ and $`r^b`$ denotes the right multiplication by $`b`$. This obviously maps the anihilator of the space of $`P`$–equivariant functions to itself and thus descends to an action on $`𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$. If $`s`$ is equivariant, then $`(sr^{b^1})(g)=b(s(g))`$, and thus $`(bD)(s)(g)=b(D(s)(gb))`$. But this implies that pairing $`j^{\mathrm{}}s(e)`$ with $`bDb\lambda `$ we get $`(b\lambda )((bD)(s)(e))=\lambda (D(s)(b))`$. On the other hand, the action of $`b`$ on $`J^{\mathrm{}}(E)_o`$ is given by $`bj^{\mathrm{}}s(e)=j^{\mathrm{}}(s\mathrm{}_{b^1})(e)`$, where $`\mathrm{}_b`$ denotes the left multiplication by $`b`$. Thus pairing $`b^1j^{\mathrm{}}s(e)`$ with $`D\lambda `$ we get $`\lambda (D(s\mathrm{}_b)(e))`$, which by left invariance of $`D`$ coincides with $`\lambda (D(s)(b))`$.
Now for any $`k`$, we have the natural projection $`J^{\mathrm{}}(E)_oJ^k(E)_o`$. On the other hand, the universal enveloping algebra $`𝒰(𝔤)`$ has a natural (infinite) filtration $`𝕂=𝒰^0(𝔤)𝒰^1(𝔤)\mathrm{}`$ such that $`𝒰(𝔤)=_i𝒰^i(𝔤)`$. In the picture of left invariant differential operators on $`G`$, this is just the filtration by the order of operators. This filtration clearly induces a filtration $`^i`$ on $`𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$, and each filtration component is a $`P`$–submodule (but not a $`𝔤`$–submodule). The pairing of an element of $`^k`$ with an element $`j^{\mathrm{}}(s)(e)J^{\mathrm{}}(E)_o`$ clearly depends only on $`j^ks(e)`$, so we get an induced paring between $`^k`$ and $`J^k(E)_o`$, and this induced pairing is obviously non–degenerate and still compatible with the $`P`$–actions, so since both sides are finite dimensional, they are dual $`P`$–modules.
Let us remark at this point that it is also possible to put locally convex topologies on the spaces in question, such that they become topologically dual $`(𝔤,P)`$–modules. Namely, one has to view $`J^{\mathrm{}}(E)_o`$ as the limit of the system $`\mathrm{}J^k(E)_oJ^{k1}(E)_o\mathrm{}`$, while $`𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$ has to be topologized as a direct sum of finite dimensional spaces.
### A.2.
Let $`𝔼`$ and $`𝔽`$ be $`P`$–representations, $`E`$ and $`F`$ the corresponding bundles and $`\phi :J^k(E)_oF_o=𝔽`$ a $`P`$–homomorphism. By the duality shown above, we can view the dual map $`\phi ^{}`$ as a $`P`$–homomorphism $`𝔽^{}^k𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$. Conversely, if we have a $`P`$–homomorphism $`𝔽^{}𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$, then this has values in some $`^i`$ since $`𝔽^{}`$ is finite dimensional, so dualizing it corresponds to a $`P`$–homomorphism $`J^i(E)_oF_o`$. Consequently, we see that the space of invariant operators $`\mathrm{\Gamma }(E)\mathrm{\Gamma }(F)`$ is isomorphic to $`\mathrm{Hom}_P(𝔽^{},𝒰(𝔤)_{𝒰(𝔭)}𝔼^{})`$.
By Frobenius reciprocity the latter space is isomorphic to
$$\mathrm{Hom}_{(𝔤,P)}(𝒰(𝔤)_{𝒰(𝔭)}𝔽^{},𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}).$$
This isomorphism is quite simple to prove: If $`\phi :𝔽^{}𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$ is a $`P`$–homomorphism, then $`\stackrel{~}{\mathrm{\Phi }}(A\lambda )=A\phi (\lambda )`$ defines a $`(𝔤,P)`$–homomorphism $`𝒰(𝔤)𝔽^{}𝒰(𝔤)_{𝒰(𝔭)}𝔼^{}`$, and since $`\phi `$ is a $`P`$–homomorphism, this factors to a $`(𝔤,P)`$–homomorphism $`\mathrm{\Phi }`$ between the required spaces. Conversely, we put $`\phi (\lambda )=\mathrm{\Phi }(1\lambda )`$ and this clearly is a $`P`$–homomorphism if $`\mathrm{\Phi }`$ is a $`(𝔤,P)`$–homomorphism.
## Appendix B. Adjointness of $``$ and $`^{}`$
### B.1.
As promised in the beginning of Section 4, we show that the operators $``$ and $`^{}`$ are adjoint operators with respect to a certain inner product on $`C^n(𝔤_{},𝕍)`$. To construct this inner product, we have to distinguish between the real and the complex case. Let us start with the case where $`𝔤`$ and $`𝕍`$ are complex. Since the grading element $`E𝔤_0`$ is semisimple, we can find a Cartan subalgebra $`𝔥𝔤`$ which contains $`E`$. Then each root space for this Cartan subalgebra is contained in some $`𝔤_i`$. Let $`𝔲`$ be a compact real form of $`𝔤`$ with a Cartan subalgebra $`𝔥_0`$ contained in $`𝔥`$, and let $`\sigma `$ be the complex conjugation with respect to this real form. By definition of $`E`$, the map $`\mathrm{ad}(E)\mathrm{ad}(E)`$ acts on $`𝔤_i`$ by multiplication by $`i^2`$, so for the Killing form we have $`B(E,E)>0`$. Consequently, we must have $`\sigma (E)=E`$, and thus $`\sigma (𝔤_i)=𝔤_i`$ for all $`i=k,\mathrm{},k`$. Now one immediately verifies directly that $`B^{}(X,Y):=B(X,\sigma (Y))`$ is a positive definite Hermitian inner product on $`𝔤`$, such that the decomposition $`𝔤=𝔤_k\mathrm{}𝔤_k`$ is an orthogonal direct sum. In particular, this induces a Hermitian inner product on $`𝔤_{}`$.
Next, since $`𝔲`$ is a compact real form, there is a positive definite Hermitian inner product $`,`$ on $`𝕍`$ such that the elements of $`𝔲`$ act as skew–Hermitian operators. But this immediately implies that for each $`X𝔤`$ and $`v_1,v_2𝕍`$, we have $`Xv_1,v_2=v_1,\sigma (X)v_2`$. Together with the inner product on $`𝔤_{}`$ constructed above we get a positive definite Hermitian inner product on $`C^n(𝔤_{},𝕍)`$ for each $`n`$.
In the real case, the situation is slightly more complicated. In this case we have to construct appropriate involutions $`\sigma `$ on the individual simple factors separately, and we have to distinguish between the case where the complexification of a simple factor is again simple and the case where it is not. Note that the simple factors of a $`|k|`$–graded Lie algebra are themselves $`|\mathrm{}|`$–graded for some $`\mathrm{}k`$ and that the grading element of $`𝔤`$ is just the sum of the grading elements of the simple factors.
If we have a real simple algebra $`𝔤`$ whose complexification is not simple, then it is well known that $`𝔤`$ is actually the underlying real Lie algebra of a complex simple Lie algebra. In this case, we can proceed as above to get a compact real form $`𝔲𝔤`$ and the corresponding involution $`\sigma `$.
In the case where both $`𝔤`$ and its complexification $`𝔤^{}`$ are simple, we choose a Cartan subalgebra $`𝔥𝔤^{}`$ which contains the element $`E𝔤`$. By \[29, Exposè 11, Théorème 3\] there is a compact real form $`𝔲𝔤^{}`$ with Cartan subalgebra $`𝔥_0𝔥`$ such that the complex conjugation $`\sigma `$ with respect to $`𝔲`$ commutes with the complex conjugation with respect to $`𝔤`$, and thus $`\sigma (𝔤)=𝔤`$.
The involutions on the simple factor together define an involution of $`𝔤`$ and as above one uses the Killing form on $`𝔤`$ and $`\sigma `$ to get a positive definite inner product on $`𝔤`$ and on $`𝔤_{}`$. If the representation $`𝕍`$ is not already complex, then we can pass to its complexification to get a Hermitian inner product such that $`Xv_1,v_2=v_1,\sigma (X)v_2`$ as above, an in both cases the real part of this Hermitian product gives a positive definite inner product on $`𝕍`$ which we use together with the inner product on $`𝔤_{}`$ to get a positive definite inner product on $`C^n(𝔤_{},𝕍)`$.
###### B.2 Proposition.
The differential $`:C^n(𝔤_1,𝕍)C^{n+1}(𝔤_1,𝕍)`$ and the codifferential $`^{}:C^{n+1}(𝔤_1,𝕍)C^n(𝔤_1,𝕍)`$ are adjoint operators with respect to the inner products constructed in B.1 above.
###### Proof.
The point about this is that in each case the inner product of $`f_1,f_2C^n(𝔤_{},𝕍)`$ can be computed as $`(f_2)(f_1)`$, where $``$ is a linear (over the reals) isomorphism $`C^n(𝔤_{},𝕍)C^n(𝔭_+,𝕍^{})`$. The map $``$ is defined by $`(f)(Z_1,\mathrm{},Z_n)(v):=f(\sigma (Z_1),\mathrm{},\sigma (Z_n)),v`$ for $`Z_i𝔭_+`$ and $`v𝕍`$, where $`\sigma `$ is the involution constructed in B.1 and the inner product is in $`𝕍`$. But then the compatibility of the inner product on $`𝕍`$ with the action of $`𝔤`$ implies that $`((f))=((f))`$. Thus we can compute:
$$\begin{array}{c}^{}(f_1),f_2=(f_2)(^{}(f_1))=((f_2))(f_1)\hfill \\ \hfill =((f_2))(f_1)=f_1,(f_2)\end{array}$$ |
warning/0001/hep-th0001167.html | ar5iv | text | # Discrete Symmetries (C,P,T) in Noncommutative Field Theories
## I Introduction
Recently it has been shown that the noncommutative spaces arise naturally when one studies the perturbative string theory in the presence of D-branes with non-zero B-field background, i.e. the low energy worldvolume theory on such branes is a noncommutative supersymmetric gauge theory(for a review of the field see ).
Besides the string theory arguments, the noncommutative field theories by themselves are very interesting. Generally, noncommutative version of a field theory is obtained by replacing the product of the fields appearing in the action by the \*-product:
$$(fg)(x)=exp(\frac{i}{2}\theta _{\mu \nu }_{x_\mu }_{y_\nu })f(x)g(y)|_{x=y},$$
(1)
where $`f`$ and $`g`$ are two arbitrary functions, which we assume to be infinitely differentiable. The ”Moyal Bracket” of two functions is
$$\{f,g\}_{M.B.}=fggf.$$
It is apparent that if we choose $`f`$ and $`g`$ to be the coordinates themselves we find
$$\{x_\mu ,x_\nu \}=i\theta _{\mu \nu },$$
(2)
and this is why these spaces are called noncommutative. Moreover, consistently we assume the derivatives to act trivially on this space:
$$\{x_\mu ,_\nu \}=\eta _{\mu \nu }\{_\mu ,_\nu \}=0.$$
(3)
Because of the nature of the \*-product, the noncommutative field theories for the slowly varying fields or low energies ($`\theta E^2\stackrel{<}{_{}}1`$), at classical level, effectively reduce to their commutative version. However, this is only the classical result and quantum corrections will show the effects of $`\theta `$ even at low energies . Since the derivatives are commuting after rewriting the noncommutative fields and their action in terms of the Fourier modes we find a commutative field theory in the momentum space, and this field theory has unfamiliar momentum dependent interactions. In this way we find a tool to study these theories perturbatively, like the usual commutative field theories.
It has been shown that the noncommutative version of $`\varphi ^4`$ theory in 4 dimensions is two loop renormalizable , moreover it is shown that the noncommutativity parameter, $`\theta `$, does not receive quantum corrections.
The pure noncommutative U(1) theory has been discussed and shown to be one loop renormalizable. The one loop beta function for noncommutative U(1) is negative (and hence the theory is asymptotically free). The interesting result is that this one loop beta function is not $`\theta `$ dependent . However, it is not clear whether this property remains at higher loops. It is believed that all of these one loop properties are a consequence of the fact that the planar degrees of freedom of noncommutative theories is the same as a commutative theory . The question of the renormalizability has also been addressed for noncommutative QED (noncommutative U(1)+ fermions).
In this paper we study another interesting question about noncommutative theories regarding their behaviour under discrete symmetries. Since in the noncommutative spaces we have missed the Lorentz symmetry, discrete symmetries and in particular the CPT invariance, in the context of noncommutative geometry in general, are usually non-trivial questions. This question has been very briefly discussed in . Hence, first we should build the noncommutative version of QED, NCQED. We show that there are two distinct choices for the fermion representations. We will show that these two are related by charge conjugation, so we may call them positively or negatively charged representations. We will give more intuitive explanations for these representations. In section 3, we study the behaviour of our theory under discrete symmetries. In this section we show the explicit calculations for the cases with $`\theta _{0i}=0`$ ($`x^0`$ is the time coordinate) and we only present the results for the non-zero $`\theta _{0i}`$ case in the last part of this section. For the $`\theta _{0i}=0`$ cases, we show that our theory, NCQED, is parity invariant, with the usual transformation of the fields; and studying the charge conjugation transformations we show that the the NCQED is not C-invariant and in order to make the theory invariant besides the usual field transformations we should also change $`\theta `$ by $`\theta `$. In addition we show that the same $`\theta `$ changing is needed for time reversal invariance. So, although our theory is CP violating, it is CPT invariant. For the general $`\theta `$ we show that though C , P and T are broken, the whole theory is again CPT invariant. The last section is devoted to conclusions and remarks.
## II Building the NCQED
i)Pure Gauge theory
The action for the pure gauge theory is
$$S=\frac{1}{4\pi }F_{\mu \nu }F^{\mu \nu }d^4x=\frac{1}{4\pi }F_{\mu \nu }F^{\mu \nu }d^4x$$
(4)
with
$$F_{\mu \nu }=_{[\mu }A_{\nu ]}+ig\{A_\mu ,A_\nu \}.$$
(5)
In the above $`g`$ is the gauge coupling constant. Let us consider the following transformations
$$\begin{array}{cc}A_\mu A_\mu ^{}=U(x)A_\mu U^1(x)+& \\ \frac{i}{g}U(x)_\mu U^1(x),& \\ U(x)=exp(i\lambda ),U^1(x)=exp(i\lambda ),& \end{array}$$
(6)
where
$$\begin{array}{cc}exp(i\lambda (x))1+i\lambda \frac{1}{2}\lambda \lambda \frac{i}{3!}\lambda \lambda \lambda +\mathrm{}& \\ U(x)U^1(x)=1.& \end{array}$$
(7)
Under the above transformations
$$F_{\mu \nu }F_{\mu \nu }^{}=U(x)F_{\mu \nu }U^1(x).$$
(8)
Hence due to the cyclic property of integration over space (see Appendix A) the action is invariant under (6). The above argument can be trivially generalized to U(N) cases by letting $`A`$ and $`\lambda `$ take values in the related algebra. However, in this paper we only consider the U(1) case.
ii)Fermionic Part
In order to write down the fermionic part of the action with the noncommutative U(1) symmetry mentioned above, first we need to find the proper ”fundamental” representations of the noncommutative U(1) group. There are two distinct choices for that:
a) The representation with
$$\{\begin{array}{cc}\psi _+(x)\psi _+^{}(x)=U(x)\psi _+(x),& \\ \overline{\psi }_+(x)\overline{\psi }_+^{}(x)=\overline{\psi }_+(x)U^1(x),& \end{array}$$
(9)
and
b) the other with
$$\{\begin{array}{cc}\psi _{}(x)\psi _{}^{}(x)=\psi _{}(x)U^1(x),& \\ \overline{\psi }_{}(x)\overline{\psi }_{}^{}(x)=U(x)\overline{\psi }_{}(x).& \end{array}$$
(10)
We will show that this two type of fermions are related by a charge conjugation transformation. The next step is finding a ”covariant” derivative, $`D_\mu `$. For these two types of fermions we have different ”covariant” derivatives
a)
$$D_+^\mu \psi _+(x)^\mu \psi _+(x)ig\psi _+(x)A^\mu (x),\mathrm{and}$$
(11)
and b)
$$D_{}^\mu \psi _{}(x)^\mu \psi _{}(x)+igA^\mu (x)\psi _{}(x).$$
(12)
One can show that with each of the covariant derivatives defined above (with the proper fermionic representation ), the action
$$S=d^4x\overline{\psi }(i\gamma ^\mu D_\mu \psi m\psi )$$
(13)
is invariant under the noncommutative U(1) transformations.
## III P, C and T Invariance
Having the form of the action we are ready to study the P, C, and T symmetries. For the sake of certainty up to the last part of this section we consider the noncommutative spaces, i.e. $`\theta _{0i}=0`$, and in the last paragraph we discuss the non-zero $`\theta _{0i}`$ and the most general noncommutative space-time.
Parity
Under the parity, $`x_ix_i`$, the $`\theta `$ parameter is not changed (see (2)). It is straightforward to show that for parity transformation given by:
$$\{\begin{array}{cc}A_0A_0& \\ A_iA_i& \\ \psi (x)\gamma ^0\psi & \\ x_ix_i,& \end{array}$$
(14)
the NCQED action is invariant for both of the fermionic choices.
Charge Conjugation
Let us first study the pure noncommutative U(1) case. Under the usual charge conjugation, C-, transformations,
$$A_\mu A_\mu ,$$
(15)
(4) is not invariant, because the first term in the $`F`$ will change the sign but the second term, $`\{A_\mu ,A_\nu \}`$, will remain the same. To make the theory C invariant we note that the Moyal bracket changes the sign if together with (15) we also change $`\theta `$ by
$$\theta \theta .$$
(16)
The above $`\theta `$ transformation has an intuitive explanation. As discussed in , the gauge symmetry of noncommutative U(1) is an infinite dimensional algebra which its Cartan subalgebra (the zero momentum sector) is a U(1) leading to a photon like state, and all the other gauge particles look like dipoles under this U(1), whose dipole moment is proportional to the $`\theta `$. So, in this picture we feel the necessity of (16), under charge conjugation.
Hence, the noncommutative U(1) theory with parameter $`\theta `$ is mapped into another noncommutative U(1) theory with $`\theta `$.
Now, we should consider the fermionic part. Since the kinetic part of the fermionic action is unchanged, we take the usual C-transformations:
$$\{\begin{array}{cc}\psi i\gamma ^0\gamma ^2\overline{\psi }^T=i\gamma ^2\psi ^{}& \\ \overline{\psi }i\psi ^T\gamma ^2\gamma ^0.& \end{array}$$
(17)
Let us first discuss the fermions in $`+`$ representation (type a) fermions). Under the above transformations, without changing $`\theta `$
$$\begin{array}{cc}d^4x\overline{\psi }(i\gamma ^\mu A_\mu (x)\psi )& \\ d^4x\overline{\psi }(i\gamma ^\mu \psi A_\mu (x)),& \end{array}$$
(18)
we see that this is exactly the form of interaction term for the type b) fermions. In other words, the types a) and b) fermions are charge conjugate of each other. Let us consider the $`\theta `$ transformation too. By using roles given in Appendix A, we see that forms of the interaction term for these two types of fermions are related by (16), which means that (17) together with (16) and (15) give proper charge conjugation transformations, a discrete symmetry of NCQED.
Time Reversal
First we consider the pure noncommutative U(1) and then we study fermions. Under the time reversal, in order to keep the kinetic part of our gauge field action, $`A_\mu `$ should transform as
$$\{\begin{array}{cc}A_0A_0& \\ A_iA_i.& \end{array}$$
(19)
Now let us look at the terms with Moyal brackets. Since time reversal operator involves a complex conjugation, for any two real fields, $`f`$ and $`g`$ we have
$$fg|_\theta f_Tg_T|_\theta =g_Tf_T|_\theta ,$$
(20)
where $`f_T`$ and $`g_T`$ show the time reversed $`f`$ and $`g`$ respectively, then we have
$$\{f,g\}_{M.B.}\{f_T,g_T\}_{M.B.}.$$
(21)
Since $`A_\mu `$’s are real fields,
$$ig\{A^\mu ,A^\nu \}ig\{A_T^\mu ,A_T^\nu \},$$
(22)
and
$$\begin{array}{cc}F_{0i}=_{[0}A_{i]}+ig\{A_0,A_i\}& \\ _{[0}A_{i]}ig\{A_0,A_i\},& \\ F_{ij}=_{[i}A_{j]}+ig\{A_i,A_j\}& \\ F_{ij}=_{[i}A_{j]}ig\{A_i,A_i\},& \end{array}$$
(23)
the only way to make the theory invariant under time reversal is changing $`\theta `$ as well as $`A_\mu `$:
$$\theta \theta .$$
(24)
So (19) together with (24) give the proper time reversal transformations.
The Fermionic Part
Since the kinetic term is quadratic in fields, $`\psi `$’s should obey the usual time reversal transformations:
$$\{\begin{array}{cc}\psi i\gamma ^1\gamma ^3\psi & \\ \overline{\psi }i\overline{\psi }\gamma ^1\gamma ^3.& \end{array}$$
(25)
As for the interaction term, for the sake of certainty let us consider the type a) case, without changing $`\theta `$. We find
$$\begin{array}{cc}d^4x\overline{\psi }(i\gamma ^\mu A_\mu (x)\psi )& \\ d^4x\overline{\psi }_T(i\gamma ^\mu A_\mu ^T(x)\psi _T)|_\theta ,& \end{array}$$
(26)
where $`\gamma ^\mu `$ is the complex conjugate of $`\gamma ^\mu `$. Replacing $`\psi _T`$ and $`A_T`$ from (25) and (19), we obtain
$$\begin{array}{cc}d^4x\overline{\psi }_T(i\gamma ^\mu A_\mu ^T(x)\psi _T)|_\theta =& \\ d^4x\overline{\psi }(i\gamma ^\mu A_\mu (x)\psi )|_\theta ,& \end{array}$$
(27)
which is exactly the interaction term for type b) fermions. As we see, in order to make the NCQED time reversal invariant, we should consider (24), (25) and (19) together.
CPT
Now that we have studied P, C, and T, it is interesting to consider the CP and CPT too. As we showed parity transformations remain the same as the commutative version, however each of C and T involves an extra $`\theta \theta `$. So altogether the NCQED (with parameter $`\theta `$) is CP violating, i.e. , it maps the theory into NCQED with $`\theta `$, and the theory is CPT invariant. We should note that although our system is not manifestly Lorentz invariant CPT, as an accidental symmetry, remains valid.
Non-zero $`\theta _{0i}`$ and general $`\theta _{\mu \nu }`$
Although a well-defined Hamiltonian for non-zero $`\theta _{0i}`$ cases is not found yet and hence the quantum theory for these cases is not understood in the same sense as $`\theta _{ij}`$ case, one can formally study the discrete symmetries for these cases. As it is readily seen from (1.2) under parity the $`\theta _{0i}`$ components should be replaced with $`\theta _{0i}`$, and we can show that the (3.1) transformations together with this $`\theta `$ change is the symmetry of NCQED.
For the charge conjugation to make the theory invariant the change in $`\theta `$ parameter, (3.3), should be extended to $`\theta _{0i}`$ components too.
It is straightforward to check that the time reversal invariance is achieved if $`\theta _{0i}`$ are unchanged while the $`\theta _{ij}`$ components should be transformed by (24). Having in mind that time reversal involves a complex conjugation, this result is expected from (1.2). Hence for general $`\theta _{\mu \nu }`$ the theory remains CPT invariant, although the theory violates P, C and T.
## IV Conclusions and Remarks
In this paper we have reviewed the noncommutative gauge theory and their gauge symmetry and shown that fermions can be added in two distinct fundamental representations of the gauge group. We have shown that these two representations are related by charge conjugation, so we called them positive or negative representations.
Studying the discrete symmetries for the $`\theta _{0i}=0`$ cases, we have shown that NCQED is parity invariant under the usual (commutative ) field transformations. For C and T transformations we showed that besides the usual field transformations we need an extra $`\theta \theta `$ transformation. In other words, NCQED with $`\theta `$ is charge conjugated (or time reversed) of NCQED with $`\theta `$. Therefore, despite being Lorentz non-invariant, in this case NCQED is CT invariant, and hence CPT invariant. In other words P and CPT is an accidental symmetry of the system.
For the general $`\theta _{\mu \nu }`$, we discussed that P, C and T invariance are broken, however the the theory is again CPT invariant.
Noncommutative gauge theories seem to provide a very good framework for the CP violating models, which are of great importance in particle physics phenomenology. The advantage of these theories is that the beta function is not $`\theta `$ dependent and furthermore $`\theta `$ does not receive quantum corrections. Therefore the amount of CP violation is completely under control.
Acknowledgements I would like to thank D.Demir and Y. Farzan for fruitful discussions. I would also like to thank D. Ployakov for reading the manuscript.
This research was partly supported by the EC contract no. ERBFMRX-CT 96-0090. Appendix: Some Useful Identities in \*-product calculus
Let $`f,g`$ be two arbitrary functions on noncommutative $`R^d`$:
$$f(x)=f(k)e^{ik.x}d^dk,g(x)=g(k)e^{ik.x}d^dk.$$
Then
$$(fg)(x)=f(k)g(l)e^{ik\theta l/2}e^{i(k+l).x}d^dkd^dl,$$
where $`k\theta l=k^\mu \theta _{\mu \nu }l^\nu `$. From the above relation it is straightforward to see:
1) $`gf=fg|_{\theta \theta }`$, and hence $`\{f,g\}_{M.B.}=fg|_\theta fg|_\theta `$.
2) $`(fg)(x)d^dx=(gf)(x)d^dx=fg(x)d^dx`$.
3) If we denote complex conjugation by $`c.c.`$, then
$`(fg)^{c.c.}=g^{c.c.}f^{c.c.}`$.
If $`h`$ is another arbitrary function:
4) $`(fg)h=f(gh)fgh`$.
5)$`(fgh)(x)d^dx=(hfg)(x)d^dx=(ghf)(x)d^dx`$.
6) $`(fgh)|_\theta =(hgf)|_\theta `$.
In other words the integration on the space coordinates, $`x`$, has the cyclic property, and it has all the properties of the $`Tr`$ in the matrix calculus.
From 2) we learn that the kinetic part of the actions (which are quadratic in fields) is the same as their commutative version. So the free field propagators in commutative and noncommutative spaces are the same. |
warning/0001/nlin0001053.html | ar5iv | text | # On the Drach superintegrable systems.
## 1 Introduction
In 1935 Jules Drach applied direct method for search of the integrable Hamiltonian systems of two degrees of freedom, which admit a cubic second integral . Recall, the direct approach leads to a complicated set of nonlinear equations, whose nonlinearity has no a priory restriction. We can attempt to solve these second invariant differential equations using various simplifying assumptions. The Drach ansatz for the Hamilton function $`H`$ and for the second cubic invariant $`K`$
$`H`$ $`=`$ $`p_xp_y+U(x,y),`$
$`K`$ $`=`$ $`6w(x,y)\left({\displaystyle \frac{H}{x}}p_yp_x{\displaystyle \frac{H}{y}}\right)P(p_x,p_y,x,y)`$
yields ten new integrable systems. Recall, for the the fixed polynomial $`P(p_x,p_y,x,y)`$ the potential $`U(x,y)`$ and function $`w(x,y)`$ be solution of some differential equations .
For the other natural Hamilton functions
$$H=p_x^2+p_y^2+V(x,y)$$
(1.2)
the similar approach has been used by Fokas and Lagerstrom , by Holt and by Thompson . Note, that the some Drach results have been rediscovered in these papers.
The most complete classifications of known results was later brought together by Hietarinta in 1987 . Recently , two-dimensional hamiltonian systems with the cubic integrals were investigated using the Jacobi change of the time. In a complete list of all known systems was extended in comparison with .
The aim of this note is to study the Drach systems and some other degenerate systems on the plane with the cubic in momenta integrals of motion. We prove that eight Drach hamiltonians belong to the Stackel family of integrals and, moreover, seven of them are degenerate systems.
Recall, the system is called superintegrable or degenerate if the Hamilton function $`H`$ is in the involution with two integrals of motion $`I`$ and $`K`$, such that
$$\{H,I\}=\{H,K\}=0,\{I,K\}=J(H,I,K).$$
(1.3)
Initial integrals and the constant of motion $`J(H,I,K)`$ are generators of the polynomial associative algebra , whose defining relations are polynomials of certain order in generators.
Below we shall consider two-dimensional systems with two quadratic integrals of motion $`I_1=H,I_2=I`$ and one qubic integral $`K`$. So, by the Bernard-Darboux theorem the system with integrals $`I_1,I_2`$ belong to the Stäckel family of integrable systems . Therefore, let us begin with remaining of some necessary results about the Stäckel systems .
## 2 The Stäckel systems
The systems associated with the name of Stäckel are holonomic systems on the phase space $`^{2n}`$ equipped with the canonical variables $`\{p_j,q_j\}_{j=1}^n`$. The nondegenerate $`n\times n`$ Stäckel matrix $`𝐒`$, with entries $`s_{kj}`$ depending only on $`q_j`$
$$det𝐒0,\frac{s_{kj}}{q_m}=0,jm$$
defines $`n`$ functionally independent integrals of motion
$$I_k=\underset{j=1}{\overset{n}{}}c_{jk}\left(p_j^2+U_j\right),c_{jk}=\frac{𝔰^{kj}}{det𝐒},$$
(2.1)
which are quadratic in momenta. Here $`𝐂=[c_{jk}]`$ denotes inverse matrix to $`𝐒`$ and $`𝔰^{kj}`$ be cofactor of the element $`s_{kj}`$. The common level surface of these integrals
$$M_\alpha =\{z^{2n}:I_k(z)=\alpha _k,k=1,\mathrm{},n\}$$
is diffeomorphic to the $`n`$-dimensional real torus and one immediately gets
$$p_j^2=\left(\frac{𝒮}{q_j}\right)^2=\underset{i=1}{\overset{n}{}}\alpha _is_{ij}(q_j)U_j(q_j).$$
(2.2)
Here $`𝒮(q_1\mathrm{},q_n)`$ is a reduced action function . For the rational entries of $`𝐒`$ and rational potentials $`U_j(q_j)`$ one gets
$$p_j^2=\frac{\underset{i=1}{\overset{k}{}}(q_je_i)}{\phi _j^2(q_j)},$$
(2.3)
where $`e_i`$ are constants of motion and functions $`\phi _j(q_j)`$ depend on coordinate $`q_j`$ and numerical constants . The Riemann surfaces (2.3) are isomorphic to the canonical hyperelliptic curves
$$𝒞_j:\mu _j^2=\underset{i=1}{\overset{k}{}}(\lambda e_i),\mu _j=\phi (q_j)p_j,$$
(2.4)
where the senior degree $`k`$ of polynomial fixes the genus $`g_j=[(k1)/2]`$ of the algebraic curve $`𝒞_j`$. Considered together, these curves determine an $`n`$-dimensional Lagrangian submanifold in $`^{2n}`$
$$𝒞^{(n)}:𝒞_1(p_1,q_1)\times 𝒞_2(p_2,q_2)\times \mathrm{}\times 𝒞_n(p_n,q_n).$$
The Abel transformation linearizes equations of motion on $`𝒞^{(n)}`$ by using first kind abelian differentials on the corresponding algebraic curves . The basis of first kind abelian differentials is uniquely related to the Stäckel matrix $`𝐒`$ .
Now let us turn to the superintegrable or degenerate systems in the classical mechanics. One of the main examples of the two-dimensional superintegrable systems is the isotropic harmonic oscillator, which has many common properties with the Drach degenerate systems. Recall, for the oscillator the Hamilton function and the second integral of motion look like
$$H=p_1^2+p_2^2+q_1^2+q_2^2,I=p_1^2+q_1^2p_2^2q_2^2.$$
Obviously, the angular momentum
$$K=q_1p_2p_1q_2=\frac{1}{2}\left(p_1\frac{dp_2}{dt}\frac{dp_1}{dt}p_2\right).$$
(2.5)
is the third integral of motion. Two pairs of quadratic integrals $`I_1=H,I_2=I`$ and $`\stackrel{~}{I}_1=H,\stackrel{~}{I}_2=K^2`$ are associated with the following Stäckel matrices
$$𝐒=\left(\begin{array}{cc}1& 1\\ 1& 1\end{array}\right),\text{and}\stackrel{~}{𝐒}=\left(\begin{array}{cc}1& 0\\ r^2& 1\end{array}\right),r^2=x^2+y^2$$
respectively. So, the corresponding equations of motion may be separated in the different curvilinear coordinate systems.
For this degenerate Stäckel system and for all other known superintegrable Stäckel systems with quadratic integrals of motion the number degree of freedom $`n>g`$ is always more then sum of genuses $`g_j`$ of the corresponding algebraic curves. In this case the number of independent first kind abelian differentials be insufficient for the inversion of the Abel-Jacobi map on $`𝒞^{(n)}`$.
To construct inversion of this map for the degenerate systems one has to complete a given basis of the differentials to the set of $`n`$ differentials. We have some freedom in a choice of complimentary differentials and, therefore, we can associate the different Stäckel matrices to one given Hamilton function . By using first kind abelian differentials one gets superintegrable Stäckel system with quadratic integrals only. Of course, we can try to add the second and third kind abelian differentials, but we do not know such examples.
Below we prove that for all the known superintegrable systems with a qubic integral $`K`$ the number degree of freedom $`n=2`$ is more then sum $`g=g_1+g_2=0`$ of genuses $`g_j=0`$ of the associated Riemann surfaces (2.4) too. The corresponding dynamics is splitting on two spheres
$$𝒞_{1,2}:\mu ^2=\alpha _{1,2}\lambda ^2+\beta _{1,2}\lambda +\gamma _{1,2},g_{1,2}=0,$$
(2.6)
where $`\alpha _j,\beta _j`$ and $`\gamma _j`$ be the constants of motion.
In variables $`\mu _{1,2}`$ (2.4) the additional cubic integral of motion for all the degenerate Drach systems looks like
$$K=\frac{det𝐒}{s_{21}s_{22}}\left(\mu _1\frac{d\mu _2}{dt}\frac{d\mu _1}{dt}\mu _2\right).$$
(2.7)
This generalized ”angular momentum” gives rise to the first order integrals (2.5) or the third order polynomials in momenta depending on the Stäckel matrices $`𝐒`$ and potentials $`U_j`$. In our case it will be qubic integral, which coincides with the initial Drach integral up to numerical factor.
To consider nonlinear algebra of integrals of motion for the Drach systems we shall introduce new generators $`\{N,a,a^{}\}`$ instead of the two quadratic integrals $`I_1=H,I_2=I`$, one qubic integrals $`K`$ and the constant of motion $`J`$ (1.3). Similar to oscillator these new generators have the following properties
$$\begin{array}{cc}\{N,a\}=a,\hfill & \{N,a^{}\}=a^{},\hfill \\ & \\ \{a,a^{}\}=\mathrm{\Phi }(I_1,I_2),\hfill & aa^{}=\mathrm{\Psi }(I_1,I_2).\hfill \end{array}$$
(2.8)
Here generator $`N(I_1,I_2)`$, functions $`\mathrm{\Psi }(I_1,I_2)`$ and $`\mathrm{\Phi }(I_1,I_2)`$ depend on the quadratic Stäckel integrals only. Two other generators $`a`$ and $`a^{}`$ be functions on the all three constants of motion $`I_1,I_2,K`$, such that
$$K=\rho (I_1,I_2)(aa^{}),J(I_1,I_2,K)=\frac{a+a^{}}{2}.$$
The relations (2.8) remind the deformed oscillator algebra, which is widely used for the superintegrable systems with quadratic integrals of motion . However, instead of the usual quadratic algebra of integrals we shall get more complicated algebras of integrals.
## 3 The Drach systems
Let us reproduce corrected Drach results in his notations
$`(a)U`$ $`=`$ $`{\displaystyle \frac{\alpha }{xy}}+\beta x^{r_1}y^{r_2}+\gamma x^{r_2}y^{r_1},\text{where}r_j^2+3r_j+3=0,`$ (3.1)
$`P`$ $`=`$ $`(xp_xp_yy)^3,w=x^2y^2/2,`$
$`(b)U`$ $`=`$ $`{\displaystyle \frac{\alpha }{\sqrt{xy}}}+{\displaystyle \frac{\beta }{(y\mu x)^2}}+{\displaystyle \frac{\gamma (y+\mu x)}{\sqrt{xy}(y\mu x)^2}},`$ (3.2)
$`P`$ $`=`$ $`3(xp_xp_yy)^2(p_x+\mu p_y),w=xy(y\mu x),`$
$`(c)U`$ $`=`$ $`\alpha xy+{\displaystyle \frac{\beta }{(yax)^2}}+{\displaystyle \frac{\gamma }{(y+ax)^2}},`$ (3.3)
$`P`$ $`=`$ $`3(xp_xp_yy)(p_x^2a^2p_y^2),w=(y^2a^2x^2)/2,`$
$`(d)U`$ $`=`$ $`{\displaystyle \frac{\alpha }{\sqrt{y(xa)}}}+{\displaystyle \frac{\beta }{\sqrt{y(x+a)}}}+{\displaystyle \frac{\gamma x}{\sqrt{x^2a^2}}},`$ (3.4)
$`P`$ $`=`$ $`3p_y\left[(xp_xp_yy)^2a^2p_x^2\right],w=y(x^2a^2),`$
$`(e)U`$ $`=`$ $`{\displaystyle \frac{\alpha }{\sqrt{xy}}}+{\displaystyle \frac{\beta }{\sqrt{x}}}+{\displaystyle \frac{\gamma }{\sqrt{y}}},`$ (3.5)
$`P`$ $`=`$ $`3p_yp_x(xp_xp_yy),w=2xy,`$
$`(f)U`$ $`=`$ $`\alpha xy+\beta y{\displaystyle \frac{2x^2+c}{\sqrt{x^2+c}}}+{\displaystyle \frac{\gamma x}{\sqrt{x^2+c}}},`$ (3.6)
$`P`$ $`=`$ $`3p_y^2(xp_xyp_y),w=(x^2+c)/2,`$
$`(g)U`$ $`=`$ $`{\displaystyle \frac{\alpha }{(y+3mx)^2}}+\beta (y3mx)+\gamma (ymx)(y9mx),`$ (3.7)
$`P`$ $`=`$ $`(p_x+3mp_y)^2(p_x3mp_y),w=m(y+3mx),`$
$`(h)U`$ $`=`$ $`(y+{\displaystyle \frac{mx}{3}})^{2/3}\left[\alpha +\beta (ymx/3)+\gamma (y^2{\displaystyle \frac{14mxy}{3}}+{\displaystyle \frac{m^2x^2}{9}})\right],`$ (3.8)
$`P`$ $`=`$ $`(p_x{\displaystyle \frac{mp_y}{3}})\left(p_x^2+{\displaystyle \frac{10mp_xp_y}{3}}+{\displaystyle \frac{m^2p_y^2}{9}}\right),w=m(y+{\displaystyle \frac{mx}{3}}),`$
$`(k)U`$ $`=`$ $`\alpha y^{1/2}+\beta xy^{1/2}+\gamma x,`$ (3.9)
$`P`$ $`=`$ $`3p_x^2p_y,w=y,`$
$`(l)U`$ $`=`$ $`\alpha \left(y{\displaystyle \frac{\rho x}{3}}\right)+\beta x^{1/2}+\gamma x^{1/2}(y\rho x),`$ (3.10)
$`P`$ $`=`$ $`3p_xp_y^2+\rho p_y^3,w=x.`$
Here $`\alpha `$, $`\beta `$, $`\gamma `$, $`\mu `$, $`\rho `$, $`a`$, $`c`$, and $`m`$ be arbitrary parameters. In compare with we corrected function $`w`$ in the case (g) (3.7) and revised potential $`U`$ in the case (k) (3.9). Namely this corrected Hamiltonian is in the involution with the initial Drach cubic integral $`K`$ (LABEL:drint).
With an exception of three cases (a) (3.1), (h) (3.8) and (l) (3.10), other Drach systems are degenerate or superintegrable Stäckel systems. The separation variables associated with the pair of quadratic integrals $`\{I_1=H,I_2\}`$ are the Stäckel variables. Equations of motion may be integrated in quadratures , but these quadratures depend on the value of quadratic integral $`I_2`$. Thus, instead of the solution of initial Drach problem related to integrals $`\{H,K\}`$ we shall solve the associated problem with quadratic integrals $`\{H,I_2\}`$.
In the case (h) (3.8) we also have the Stäckel systems . Only in this case (h) (3.8) dynamics is splitting on two tori and the number degrees of freedom is equal to the sum of genuses $`g=n=2`$, such that the corresponding system is non-degenerate.
In the case (l) (3.10) the Hamilton function coincides with the hamiltonian of the previous Stäckel system (3.9) at $`\rho =0`$. Here we shall not consider this generalized Stäckel system at $`\rho 0`$.
Below we shall consider the Drach integrals (LABEL:drint) up to linear transformations of the coordinates and a rescaling of these integrals. It allows us to remove some parameters in the Hamilton functions without loss of generality. To associate the degenerate Drach hamiltonians with the Stäckel matrices $`𝐒`$ we can join these systems into the four pairs of the systems with a common matrices $`𝐒`$.
### 3.1 Case (a)
In our previous paper , the first Drach system (3.1) has been related to the three-particle periodic Toda lattice in the center-of-mass frame. Namely, after canonical change of the time $`t=q_{n+1}`$ and the Hamiltonian $`H=p_{n+1}`$ at the extended phase space
$$d\stackrel{~}{t}=(xy)^1dt,\stackrel{~}{H}=xy(H+\delta ),$$
and after further canonical transformation of other variables
$$x=e^{\frac{q_1+iq_2}{2}},p_x=(p_1ip_2)e^{\frac{q_1+iq_2}{2}},y=e^{\frac{q_1iq_2}{2}},p_y=(p_1+ip_2)e^{\frac{q_1iq_2}{2}},$$
the Hamilton function (3.1) becomes
$$\stackrel{~}{H}=p_1^2+p_2^2+\beta e^{{\displaystyle \frac{1}{2}}q_1{\displaystyle \frac{\sqrt{3}}{2}}q_2}+\gamma e^{{\displaystyle \frac{1}{2}}q_1+{\displaystyle \frac{\sqrt{3}}{2}}q_2}+\delta e^{q_1}+\alpha .$$
It is the Hamiltonian of the tree-particle periodic Toda lattice in the center-of-mass frame. The separation variables survive at the change of the time. Thus, for the first Drach system we can separate variables and integrate equations of motions in quadratures repeating the calculations for the Toda chain .
Later in Thompson considered this system too. In fact, after point transformation
$$x=re^{i\varphi },p_x=\frac{e^{i\varphi }}{2}\left(p_rip_\varphi r^1\right),y=re^{i\varphi },p_y=\frac{e^{i\varphi }}{2}\left(p_r+ip_\varphi r^1\right)$$
the Drach hamiltonian $`H`$ (LABEL:drint) looks like
$$H=p_r^2+\frac{p_\varphi ^2}{r^2}+U(r,\varphi ),$$
up to numerical factor. Namely this Hamilton function was studied in and . The special substitution of the potential $`U(r,\varphi )`$ into the Drach equations leads to the following equation
$$U(r,\varphi )=\frac{f(\varphi )+f^{\prime \prime }(\varphi )}{r^3},f^{\prime \prime \prime }f^{\prime \prime }2f^{\prime \prime }f^{}3f^{}f=0,$$
introduced in . Of course, the same equation follows from the functional equation on the Toda potential .
### 3.2 Cases (b) and (e)
Put $`\mu =1`$ in (3.2). Let us introduce the Stäckel matrix
$$𝐒_{be}=\left(\begin{array}{cc}q_1^2& q_2^2\\ & \\ 1& 1\end{array}\right),$$
(3.11)
and take the following potentials
$$\begin{array}{ccc}(b)\hfill & U_1=2\alpha \frac{\beta 2\gamma }{q_1^2},\hfill & U_2=2\alpha \frac{\beta +2\gamma }{q_2^2},\hfill \\ & & \\ (e)\hfill & U_1=2\alpha +2(\beta +\gamma )q_1,\hfill & U_2=2\alpha 2(\beta \gamma )q_2.\hfill \end{array}$$
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ for the Drach systems (3.2) and (3.5), after the following canonical point transformation
$$x=\frac{(q_1q_2)^2}{4},p_x=\frac{p_1p_2}{q_1q_2},y=\frac{(q_1+q_2)^2}{4},p_y=\frac{p_1+p_2}{q_1+q_2}.$$
The second integrals of motion $`I_2`$ (2.1) are second order polynomials in momenta. The third independent integrals of motion $`K`$ are defined by (2.7), where
$$(b)\mu _1=q_1p_1,\mu _2=q_2p_2,(e)\mu _1=p_1,\mu _2=p_2.$$
From the above definitions we can introduce generators of the nonlinear algebra of integrals (2.8) and verify properties of this algebra
$`(b)`$ $`N={\displaystyle \frac{I_2}{4\sqrt{H}}},a=J+4\sqrt{H}K,a^{}=J4\sqrt{H}K,`$
$`aa^{}=16\left(4H(\beta +2\gamma )(2\alpha +I_2)^2\right)\left(4H(\beta 2\gamma )(2\alpha I_2)^2\right),`$
$`\{a,a^{}\}=256\sqrt{H}\left(I_2(I_22\alpha )(I_2+2\alpha )4H(\beta I_24\alpha \gamma )\right),`$
and
$`(e)`$ $`N={\displaystyle \frac{I_2}{2\sqrt{H}}},a=J+2\sqrt{H}K,a^{}=J2\sqrt{H}K,`$
$`aa^{}=16\left(H(2\alpha +I_2)(\beta \gamma )^2\right)\left(H(2\alpha I_2)+(\beta +\gamma )^2\right),`$
$`\{a,a^{}\}=64H^{3/2}(I_2H\beta ^2\gamma ^2).`$
### 3.3 Cases (c) and (g)
Put $`a=1`$ in (3.3) and $`m=1/3`$ in (3.7). Let us introduce the Stäckel matrix
$$𝐒_{cg}=\left(\begin{array}{cc}\frac{1}{2}& \frac{1}{2}\\ & \\ 1& 1\end{array}\right),$$
(3.12)
and take the following potentials
$$\begin{array}{ccc}(c)\hfill & U_1=\frac{\alpha q_1^2}{4}+\frac{\gamma }{q_1^2},& U_2=\frac{\alpha q_2}{4}\frac{\beta }{q_2^2},\\ & & \\ (g)\hfill & U_1=\frac{\gamma q_1^2}{3}+\frac{\alpha }{q_1^2},& U_2=\frac{4\gamma q_2^2}{3}\beta q_2.\end{array}$$
(3.13)
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ (3.3) and (3.7), after the following canonical point transformation
$$x=\frac{q_1q_2}{2},p_x=p_1p_2,y=\frac{q_1+q_2}{2},p_y=p_1+p_2.$$
The second integrals of motion $`I_2`$ (2.1) are the second order polynomials in momenta. The third independent integrals $`K`$ are defined by (2.7), where
$$(c)\mu _1=q_1p_1\mu _2=q_2p_2,(g)\mu _1=q_1p_1,\mu _2=p_2.$$
Generators and defining relations of the nonlinear algebra of integrals (2.8) look like
$`(c)`$ $`N={\displaystyle \frac{I_2}{2\sqrt{\alpha }}},a=J+2\sqrt{\alpha }K,a^{}=J2\sqrt{\alpha }K,`$
$`aa^{}=\left(H^2+4HI_2+4I_2^24\alpha \gamma \right)\left(H^24HI_2+4I_2^2+4\alpha \beta \right),`$
$`\{a,a^{}\}=32\sqrt{a}\left(I_2(H2I_2)(H+2I_2)\alpha \beta (H+2I_2)\alpha \gamma (H2I_2)\right),`$
and
$`(g)`$ $`N={\displaystyle \frac{I_2}{4}}\sqrt{{\displaystyle \frac{3}{\gamma }}},a=J+4\sqrt{{\displaystyle \frac{\gamma }{3}}}K,a^{}=J4\sqrt{{\displaystyle \frac{\gamma }{3}}}K,`$
$`aa^{}=1/9\left(8\gamma H16\gamma I_2+3\beta ^2\right)\left(3H^2+12HI_2+12I_2^2+16\alpha \gamma \right),`$
$`\{a,a^{}\}=64\left({\displaystyle \frac{\gamma }{3}}\right)^{3/2}\left({\displaystyle \frac{(2I_2+H)(4\gamma H24\gamma I_2+3b^2)}{4\gamma }}{\displaystyle \frac{16\alpha \gamma }{3}}\right).`$
### 3.4 Cases (d) and (f)
Put $`a=1`$ in (3.4) and $`c=1`$ in (3.6). Let us introduce two the Stäckel matrices
$$𝐒_d==\left(\begin{array}{cc}1& 1\\ & \\ \frac{1}{q_1^2}& \frac{1}{q_2^2}\end{array}\right),𝐒_f=\left(\begin{array}{cc}\frac{1}{q_1}& \frac{1}{q_2}\\ & \\ \frac{1}{q_1^2}& \frac{1}{q_2^2}\end{array}\right)$$
(3.14)
and takes the following potentials
$$\begin{array}{ccc}(d)\hfill & U_1=2\gamma +\frac{2\sqrt{2}(\alpha +\beta )}{q_1},& U_2=2\gamma \frac{2\sqrt{2}(\alpha \beta )}{q_2},\\ & & \\ (f)\hfill & U_1=\frac{\gamma }{2q_1}+\frac{(\alpha +2\beta )}{4},& U_2=\frac{\gamma }{2q_2}+\frac{(\alpha 2\beta )}{4}.\end{array}$$
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ (3.4) and (3.6) up to numerical factor, after the following explicit canonical transformations
$`(d)`$ $`x={\displaystyle \frac{q_1^2+q_2^2}{2q_1q_2}},p_x={\displaystyle \frac{(p_1q_1p_2q_2)q_1q_2}{q_1^2q_2^2}},`$ $`y=q_1q_2,p_y={\displaystyle \frac{p_1q_1+p_2q_2}{2q_1q_2}},`$
$`(f)`$ $`x={\displaystyle \frac{q_1q_2}{2\sqrt{q_1q_2}}},p_x={\displaystyle \frac{2(p_1q_1p_2q_2)\sqrt{q_1q_2}}{q_1+q_2}},`$ $`y=\sqrt{q_1q_2},p_y={\displaystyle \frac{p_1q_1+p_2q_2}{\sqrt{q_1q_2}}}.`$
The second integrals of motion $`I_2`$ (2.1) are the quadratic polynomials in momenta. The third independent integrals $`K`$ are defined by (2.7), where for the both systems one gets
$$\mu _1=q_1p_1,\mu _2=q_2p_2.$$
Generators of the nonlinear algebra of integrals (2.8) are given by
$$N=\sqrt{I_2},a=J+2\sqrt{I_2}K,a^{}=J2\sqrt{I_2}K,$$
which have the following properties
$`(d)`$ $`aa^{}=16\left(I_2(2\gamma H)+2(\alpha +\beta )^2\right)\left(I_2(2\gamma H)2(\alpha \beta )^2\right),`$
$`\{a,a^{}\}=64\sqrt{I_2}\left(I_2(2\gamma H)(2\gamma +H)+2(\alpha ^2+\beta ^2)H+8\alpha \beta \gamma \right),`$
$`(f)`$ $`aa^{}={\displaystyle \frac{1}{16}}\left(4I_2(\alpha +2\beta )+(\gamma 2H)^2\right)\left(4I_2(\alpha 2\beta )+(\gamma +2H)^2\right),`$
$`\{a,a^{}\}=4\sqrt{I_2}\left((2\beta \alpha )(2\beta +\alpha )I_2+\alpha H^2+2\beta \gamma H+{\displaystyle \frac{\alpha \gamma ^2}{4}}\right).`$
### 3.5 Cases (h) and (k)
Put $`m=3`$ in (3.8). Let us introduce two the Stäckel matrices
$$𝐒_h==\left(\begin{array}{cc}q_1& q_2\\ & \\ 1& 1\end{array}\right),𝐒_k=\left(\begin{array}{cc}q_1& q_2\\ & \\ 1& 1\end{array}\right)$$
(3.15)
and take the following potentials
$$\begin{array}{ccc}(h)\hfill & U_1=\frac{\alpha \gamma {\displaystyle \frac{\beta ^2}{16}}}{4\gamma }\frac{2\gamma q_1^3}{9},& U_2=\frac{\alpha \gamma {\displaystyle \frac{\beta ^2}{16}}}{4\gamma }\frac{2\gamma q_2^3}{9}\\ & & \\ (k)\hfill & U_1=\alpha +\beta q_1+\frac{\gamma q_1^2}{2},& U_2=\alpha +\beta q_2+\frac{\gamma q_2^2}{2}.\end{array}$$
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ (3.8) and (3.9) up to numerical factor, after the following explicit canonical transformations
$$\begin{array}{ccc}(h)\hfill & x=\frac{p_2p_1}{4\sqrt{\gamma }}+\frac{(3q_1+3q_2)^{3/2}}{54}+\frac{\beta }{16\gamma },& p_x=3\frac{p_1+p_2}{\sqrt{3q_1+3q_2}}+\sqrt{\gamma }(q_1q_2),\\ & & \\ & & \\ & y=\frac{p_2p_1}{4\sqrt{\gamma }}+\frac{(3q_1+3q_2)^{3/2}}{54}\frac{\beta }{16\gamma },& p_y=3\frac{p_1+p_2}{\sqrt{3q_1+3q_2}}\sqrt{\gamma }(q_1q_2),\\ & & \\ \text{and}\hfill & & \\ & & \\ (k)\hfill & x=\frac{q_1q_2}{2},p_x=p_1p_2,& y=\frac{(q_1+q_2)^2}{4},p_y=\frac{p_1+p_2}{q_1+q_2}.\end{array}$$
Note, in the case (h) (3.8) we used non-point canonical transformation in contrast with other Drach systems.
In the last case (k) (3.9) integral of motion $`I_2`$ (2.1) is the second order polynomial in momenta. The third independent integral $`K`$ is defined by (2.7), where
$$\mu _1=p_1\text{and}\mu _2=p_2.$$
Generators and defining relations of nonlinear algebra of integrals (2.8) look like
$`(k)`$ $`N={\displaystyle \frac{I_2}{\sqrt{2\gamma }}},a=J+\sqrt{2\gamma }K,a^{}=J\sqrt{2\gamma }K,`$
$`aa^{}=\left(2\gamma (I_2+\alpha )+(H+\beta )^2\right)\left(2\gamma (I_2\alpha )+(H\beta )^2\right),`$
$`\{a,a^{}\}=4\gamma \sqrt{2c}\left(H^2+2\gamma I_2+\beta ^2\right).`$
In the case (h) (3.8) the second integral of motion $`I_2`$ (2.1) is the second order polynomial in momenta $`\{p_1,p_2\}`$. However, after the non-point transformation of variables this integral $`I_2`$ becomes the qubic in momenta $`\{p_x,p_y\}`$ Drach integral $`K`$ (3.8). The corresponding dynamics is splitting on two tori and the third order polynomial (2.7) does not commute with the Hamilton function. Later this system has been rediscovered by Holt .
## 4 Other degenerate systems on the plane with a qubic integral of motion
In this section we consider the Stäckel systems on the plane with a qubic integral of motion defined by the following Hamilton function
$$H=\frac{1}{2}\left(p_x^2+p_y^2\right)+V(x,y).$$
As above, the corresponding qubic integral will be written at the Drach form (LABEL:drint).
On the plane we know four orthogonal systems of coordinates: elliptic, parabolic, polar and cartesian. Thus, we reproduce all the known results in correspondence with the type of the associated Stäckel matrix .
The systems whose Hamilton functions separable in cartesian coordinates:
$`(A)V`$ $`=`$ $`\alpha (4x^2+y^2)+\beta x+{\displaystyle \frac{\gamma }{y^2}},`$ (4.1)
$`P`$ $`=`$ $`p_xp_y^2,w={\displaystyle \frac{y}{6}},`$ (4.2)
$`(B)V`$ $`=`$ $`\alpha (x^2+y^2)+{\displaystyle \frac{\beta }{x^2}}+{\displaystyle \frac{\gamma }{y^2}},`$ (4.3)
$`P`$ $`=`$ $`(xp_yyp_x)p_xp_y,w={\displaystyle \frac{xy}{6}},`$
$`(C)V`$ $`=`$ $`\alpha (x^2+y^2)+\beta {\displaystyle \frac{xy}{(x^2y^2)^2}},`$ (4.4)
$`P`$ $`=`$ $`(xp_yyp_x)\left(p_x^2p_y^2\right),w={\displaystyle \frac{x^2y^2}{6}},`$
$`(D)V`$ $`=`$ $`\alpha (9x^2+y^2),`$ (4.5)
$`P`$ $`=`$ $`(xp_yp_xy)p_y^2,w={\displaystyle \frac{y^2}{18}},`$
The systems whose Hamilton functions separable in parabolic coordinates:
$`(F)V`$ $`=`$ $`(\alpha +{\displaystyle \frac{\beta }{r+x}}+{\displaystyle \frac{\gamma }{rx}})r^1,r=\sqrt{x^2+y^2},`$ (4.6)
$`P`$ $`=`$ $`(xp_yp_xy)^2p_x,w={\displaystyle \frac{yr^2}{12}},`$
$`(G)V`$ $`=`$ $`(\alpha +{\displaystyle \frac{\beta x}{y^2}})r^1,`$ (4.7)
$`P`$ $`=`$ $`(xp_yp_xy)^2p_x,w={\displaystyle \frac{yr^2}{12}},`$
$`(H)V`$ $`=`$ $`(\alpha +\beta \sqrt{r+x}+\gamma \sqrt{rx})r^1,`$ (4.8)
$`P`$ $`=`$ $`(xp_yp_xy)\left(2p_x^2+2p_y^2{\displaystyle \frac{\beta }{\sqrt{r+x}}}{\displaystyle \frac{\gamma }{\sqrt{rx}}}\right),w={\displaystyle \frac{r^2}{6}},`$
One system with the Hamilton function separable in polar coordinates:
$`(I)V`$ $`=`$ $`\alpha +{\displaystyle \frac{\beta }{\sqrt{x^2+y^2}}}+{\displaystyle \frac{\rho }{(\delta x+\gamma y)^2}}+{\displaystyle \frac{\gamma x\delta y}{\sqrt{x^2+y^2}(\delta x+\gamma y)^2}},`$ (4.9)
$`P`$ $`=`$ $`(p_xyp_yx)(\gamma p_x\delta p_y),w={\displaystyle \frac{(x^2+y^2)(\delta x+\gamma y)}{12}}.`$
It is new superintegrable system with a cubic integral of motion, which is a deformation of the degenerate kepler model. An application of the direct method or the Jacobi method does not allows us to obtain this system (4.9).
Three exceptional systems whose qubic integral of motion $`K`$ we can not rewrite in the ”generalized angular momentum” form (2.7):
$`(K)V`$ $`=`$ $`\alpha \sqrt{x}\pm \beta \sqrt{y},P={\displaystyle \frac{\beta }{\alpha }}p_x^3{\displaystyle \frac{\alpha }{\beta }}p_y^3,w=\sqrt{xy},`$ (4.10)
$`(L)V`$ $`=`$ $`\alpha \left(\sqrt{x}+\beta y\right),P=p_x^3,w={\displaystyle \frac{\sqrt{x}}{2\beta }},`$ (4.11)
$`(M)V`$ $`=`$ $`f^{}(\varphi )r^2,K=p_\varphi ^2\left(\mathrm{cos}\varphi p_r\mathrm{sin}\varphi r^1p_\varphi \right)+`$
$`+`$ $`\left(2f^{}(\varphi )\mathrm{cos}\varphi f(\varphi )\mathrm{sin}(\varphi )\right)p_r+\left(3f^{}(\varphi )\mathrm{sin}\varphi +f(\varphi )\mathrm{cos}\varphi \right)r^1p_\varphi .`$
At the case (M) (4) we used the standard polar coordinates $`\{r,p_r,\varphi ,p_\varphi \}`$ and the function $`f(\varphi )`$ has to satisfy the following equation
$$f^{\prime \prime }\left(3f^{}\mathrm{sin}\varphi +f\mathrm{cos}\varphi \right)+2f^{}\left(2f^{}\mathrm{cos}\varphi f\mathrm{sin}\varphi \right)=0.$$
At these exceptional cases the Hamilton functions (4.10,4.11,4) are separable at the cartesian and polar coordinates, respectively. However, the Stäckel potentials $`U_{1,2}`$ are not polynomials in variables of separation.
### 4.1 Cartesian coordinates, cases A-D
Let us introduce the Stäckel matrix
$$𝐒_{AD}=\left(\begin{array}{cc}\frac{1}{2}& \frac{1}{2}\\ & \\ 1& 1\end{array}\right),$$
(4.13)
and take the following potentials
$$\begin{array}{ccc}(A)\hfill & U_1=8\alpha q_1^2+2\beta q_1,& U_2=2\alpha q_2^2+\frac{2\gamma }{q_2^2},\\ & & \\ (B)\hfill & U_1=2\alpha q_1^2+\frac{2\beta }{q_1^2},& U_2=2\alpha q_2^2+\frac{2\gamma }{q_2^2},\\ & & \\ (C)\hfill & U_1=\frac{\alpha q_1^2}{2}\frac{\beta }{4q_1^2},& U_2=\frac{\alpha q_2^2}{2}+\frac{\beta }{4q_2^2},\\ & & \\ (D)\hfill & U_1=18\alpha q_1^2,& U_2=2\alpha q_2^2.\end{array}$$
(4.14)
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ (4.1,4.3,4.5) and (4.4) if
$$(AB,D)x=q_1,y=q_2$$
or after the following canonical transformation
$$(C)x=\frac{q_1q_2}{2},p_x=p_1p_2,y=\frac{q_1+q_2}{2},p_y=p_1+p_2.$$
The second integrals of motion $`I_2`$ (2.1) are the second order polynomials in momenta. The third independent integrals $`K`$ are calculated by (2.7), where variables
$$(A)\mu _1=p_1\mu _2=q_2p_2,(BC)\mu _1=q_1p_1,\mu _2=p_2.$$
determine the left hand side of the canonical algebraic curves (2.4). At the case (D) the variables
$$(D)\mu _1=p_2q_1\frac{p_1q_2}{3},\mu _2=p_2q_2$$
have not such natural algebro-geometric meaning.
Generators and defining relations of the nonlinear algebra of integrals (2.8) look like
$$(AB)N=\frac{I_2}{4\sqrt{2\alpha }},a=J+4\sqrt{2\alpha }K,a^{}=J4\sqrt{2\alpha }K,$$
such that
$`(A)aa^{}`$ $`=`$ $`4(4\alpha (2I_2+H)+\beta ^2)\left((2I_2H)^264\alpha \gamma \right),`$
$`\{a,a^{}\}`$ $`=`$ $`128\alpha \sqrt{2\alpha }(2I_2H)(6I_2+H+{\displaystyle \frac{\beta ^2}{2\alpha }})64\alpha \gamma ),`$
$`(B)aa^{}`$ $`=`$ $`\left((2I_2+H)^264\alpha \beta \right)\left((2I_2H)^264\alpha \gamma \right),`$
$`\{a,a^{}\}`$ $`=`$ $`16\sqrt{2\alpha }(((2I_2H)^264\alpha \gamma )(2I_2+H)`$
$`+((2I_2+H)^264\alpha \beta )(2I_2H)).`$
At the two last cases we have
$`(C)N`$ $`=`$ $`{\displaystyle \frac{I_2}{2\sqrt{2\alpha }}},a=J+2\sqrt{2\alpha }K,a^{}=J2\sqrt{2\alpha }K,`$
$`aa^{}`$ $`=`$ $`\left((2I_2H)^22\alpha \beta \right)\left((2I_2+H)^2+2\alpha \beta \right),`$
$`\{a,a^{}\}`$ $`=`$ $`8\sqrt{2\alpha }(((2I_2H)^22\alpha \beta )(2I_2+H)`$
$`+((2I_2+H)^2+2\alpha \beta )(2I_2H)),`$
and
$`(D)N`$ $`=`$ $`{\displaystyle \frac{I_2}{6\sqrt{2\alpha }}},a=J+6\sqrt{2\alpha }K,a^{}=J6\sqrt{2\alpha }K,`$
$`aa^{}`$ $`=`$ $`4(2I_2H)^3(2I_2+H),`$
$`\{a,a^{}\}`$ $`=`$ $`96\sqrt{2\alpha }(2I_2H)^2(4I_2+H).`$
At the case (D) (4.5) the quantum counterpart of this qubic deformed oscillator algebra has been used to study of the corresponding quantum superintegrable system .
### 4.2 Parabolic coordinates, cases F-H
Let us introduce two Stäckel matrices
$$𝐒_{F,G}=\left(\begin{array}{cc}1& 1\\ & \\ q_1^1& q_2^1\end{array}\right),𝐒_H=\left(\begin{array}{cc}q_1^2& q_2^2\\ & \\ 1& 1\end{array}\right),$$
(4.15)
and take the following potentials
$$\begin{array}{ccc}(F)\hfill & U_1=\frac{\alpha }{2q_1}\frac{\beta }{2q_1^2},& U_2=\frac{\alpha }{2q_2}\frac{\gamma }{2q_2^2},\\ & & \\ (G)\hfill & U_1=\frac{\alpha }{2q_1}\frac{\beta }{4q_1^2},& U_2=\frac{\alpha }{2q_2}+\frac{\beta }{4q_2^2},\\ & & \\ (H)\hfill & U_1=4\alpha 8\sqrt{2}\beta q_1,& U_2=4\alpha +8\sqrt{2}\gamma q_2.\end{array}$$
(4.16)
The corresponding Hamilton functions $`I_1`$ (2.1) coincide with the Hamilton functions $`H`$ (4.1,4.3) and (4.4) after the following canonical point transformations
$$(F,G)x=q_1+q_2,p_x=\frac{p_1q_2p_2q_2}{q_1q_2},y=2\sqrt{q_1q_2},p_y=\frac{(p_1p_2)\sqrt{q_1q_2}}{q_1q_2},$$
and
$$(H)x=q_1^2+q_2^2,p_x=\frac{1}{2}\frac{p_1q_1p_2q_2}{q_1^2q_2^2},y=2iq_1q_2,p_y=\frac{i}{2}\frac{p_1q_2p_2q_1}{q_1^2q_2^2}.$$
The second integrals of motion $`I_2`$ (2.1) are the second order polynomials in momenta. The third independent integrals $`K`$ are calculated by (2.7), where variables
$$(F,G)\mu _1=q_1p_1\mu _2=q_2p_2,(H)\mu _1=p_1,\mu _2=p_2,$$
define the left hand side of the canonical algebraic curves (2.4).
At all these cases (F,G) and (H) the generators and defining relations of the nonlinear algebra of integrals (2.8) look like
$$(FH)N=\frac{I_2}{2\sqrt{H}},a=J+2\sqrt{H}K,a^{}=J2\sqrt{H}K,$$
such that
$`(F)aa^{}`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left(8\beta H(2I_2+\alpha )^2\right)\left(8\gamma H(2I_2\alpha )^2\right),`$
$`\{a,a^{}\}`$ $`=`$ $`2\sqrt{H}\left(I_2(2I_2+\alpha )(2I_2\alpha )2H\left(\beta (2I_2\alpha )+\gamma (2I_2+\alpha )\right)\right),`$
$`(G)aa^{}`$ $`=`$ $`{\displaystyle \frac{1}{16}}\left(4\beta H+(2I_2+\alpha )^2\right)\left(4\beta H(2I_2\alpha )^2\right),`$
$`\{a,a^{}\}`$ $`=`$ $`2\sqrt{H}\left(I_2(2I_2+\alpha )(2I_2\alpha )2\alpha \beta H\right),`$
$`(H)aa^{}`$ $`=`$ $`16\left(H(I_24\alpha )+32\gamma ^2\right)\left(H(I_2+4\alpha )32\beta ^2\right),`$
$`\{a,a^{}\}`$ $`=`$ $`64H^{3/2}(HI_216\beta ^2+16\gamma ^2).`$
### 4.3 Polar coordinates, case I
Let us introduce Stäckel matrix
$$𝐒_I=\left(\begin{array}{cc}1& 0\\ & \\ q_1^2& 1\end{array}\right),$$
(4.17)
and take the following potentials
$$\begin{array}{ccc}(I)\hfill & U_1=\alpha +\frac{\beta }{q_1},& U_2=\frac{\gamma \mathrm{cos}(q_2)\delta \mathrm{sin}(q_2)+\rho }{\left(\delta \mathrm{cos}(q_2)+\gamma \mathrm{sin}(q_2)\right)^2},\end{array}$$
(4.18)
The corresponding Hamilton function $`I_1`$ (2.1) coincides with the Hamilton function $`H`$ (4.9) if $`q_1=r`$ and $`q_2=\varphi `$ are the standard polar coordinates on the plane. The second integrals of motion $`I_2`$ (2.1) are the second order polynomials in momenta. The third independent integrals $`K`$ are calculated by (2.7), where variables
$$(I)\mu _1=p_1,\mu _2=p_2\left(\delta \mathrm{cos}(q_2)+\gamma \mathrm{sin}(q_2)\right),$$
define the canonical algebraic curves (2.4).
At $`\delta =1`$ and $`\gamma =0`$ the generators and defining relations of the nonlinear algebra of integrals (2.8) look like
$`(I)N`$ $`=`$ $`\sqrt{I_2},a=J+2\sqrt{I_2}K,a^{}=J2\sqrt{I_2}K,`$
$`aa^{}`$ $`=`$ $`(4HI_24\alpha I_2+\beta ^2)(4I_2^24\rho I_2+1),`$
$`\{a,a^{}\}`$ $`=`$ $`8\sqrt{I_2}\left((2I_2\rho )\beta ^2+(12I_2^28\rho I_2+1)(H\alpha )\right).`$ (4.19)
This system does not contain in the list of the known integrable systems . At the cases (I) (4.9) and (M) (4) we have common leading part $`P`$ of the qubic integrals $`K`$. However, at the case (M) we can not rewrite the qubic integral in the ”generalised angular momentum” form.
## 5 The Lax representation
In we proposed some construction of the $`2\times 2`$ Lax matrices for the Stäckel systems with homogeneous Stäckel matrices and with uniform potentials $`U_j=U`$. The Drach-Stäckel systems fall out from this subset of the Stäckel systems. Nevertheless, we could construct the $`2\times 2`$ Lax matrices for these systems by using various covering of the initial spheres $`𝒞_{1,2}`$ (2.6).
Here we consider the $`4\times 4`$ Lax matrices for some Drach systems by using canonical transformations of the extended phase space, which induce transformations of the Lax matrices . Recall, if the Stäckel matrices $`𝐒`$ and $`\stackrel{~}{𝐒}`$ be distinguished the first row only, the corresponding Stäckel systems are related by canonical change of the time $`q_{n+1}=t`$ and conjugated momenta $`p_{n+1}=H`$
$$t\stackrel{~}{t},d\stackrel{~}{t}=\frac{det\stackrel{~}{𝐒}}{det𝐒}dt,H\stackrel{~}{H}=\frac{det𝐒}{det\stackrel{~}{𝐒}}H.$$
(5.1)
Thus, starting with the Stäckel systems related with matrix $`𝐒=𝐒_{cg}`$ (3.12) we can study systems associated with matrices $`\stackrel{~}{𝐒}=𝐒_{be}`$ (3.11) and $`\stackrel{~}{𝐒}=𝐒_k`$ (3.15). Here subscripts mean the type of the Stäckel matrices for the different Drach systems.
The Stäckel systems with the constant matrix $`𝐒_{cg}`$ possess the following $`4\times 4`$ Lax matrices
$$(\lambda )=\left(\begin{array}{cc}L_1(\lambda ,p_1,q_1)& 0\\ & \\ 0& L_2(\lambda ,p_2,q_2)\end{array}\right),$$
(5.2)
with independent $`2\times 2`$ non-trivial blocks $`L_j(\lambda )`$. For instance, two standard blocks may be chosen
$$L_j(\lambda )=\left(\begin{array}{cc}p_j& \lambda q_j\\ \left[\frac{\varphi _j}{\lambda q_j}\right]_{MN}& p_j\end{array}\right),L_j(\lambda )=\left(\begin{array}{cc}\frac{p_jq_j}{\lambda }& 1\frac{q_j^2}{\lambda }\\ \frac{p_j^2}{\lambda }\left[\frac{\varphi _j}{1\frac{q_j^2}{\lambda }}\right]_{MN}& \frac{p_jq_j}{\lambda }\end{array}\right).$$
Here $`\varphi (\lambda )`$ is a parametric function on spectral parameter $`\lambda `$ and $`[\xi ]_N`$ is the linear combinations of the Laurent projections .
According to , canonical transformations of the extended phase space induce shift of the Lax matrices depending on the Hamilton function. Thus, by using one known Lax matrix $`(\lambda )`$ (5.2) we can construct another Lax matrices. Namely, canonical transformations of the time (5.1) give rise the following shift of the corresponding Lax matrices
$$\stackrel{~}{}(\lambda )=(\lambda )\stackrel{~}{H}\left(\begin{array}{cccc}0& 0& 0& 0\\ a& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& b& 0\end{array}\right),a,b=\pm 1\text{or}\pm i,$$
(5.3)
where values of the constants $`a,b`$ depend on the chosen form of the blocks $`L_j(\lambda )`$.
Below we present some Lax matrices constructed by designated above scheme. In the case (c) the Lax matrix is given by
$$_c(\lambda )=\left(\begin{array}{cccc}p_1& \lambda q_1& 0& 0\\ (\lambda +q_1)\left(\frac{\alpha }{4}+\gamma q_1^2\lambda ^2\right)& p_1& 0& 0\\ 0& 0& ip_2& i(\lambda q_2)\\ 0& 0& i(\lambda +q_2)\left(\frac{\alpha }{4}\beta q_2^2\lambda ^2\right)& ip_2\end{array}\right),$$
so the spectral curve
$$\mathrm{\Gamma }(\lambda ,\mu ):det(_c(\lambda )\mu I)=0$$
is a product
$$\left(\mu ^2\frac{I_1}{2}I_2+\frac{\alpha \lambda ^2}{4}+\frac{\gamma }{\lambda ^2}\right)\left(\mu ^2\frac{I_1}{2}+I_2\frac{\alpha \lambda ^2}{4}+\frac{\beta }{\lambda ^2}\right)=0,$$
of the corresponding canonical Stäckel curves (2.6).
In the case (k) the Lax matrix is given by
$$\stackrel{~}{}_k=\left(\begin{array}{cccc}p_1& \lambda q_1& 0& 0\\ \frac{\gamma (\lambda +q_1)}{2}\beta & p_1& 0& 0\\ 0& 0& ip_2& i(\lambda +q_2)\\ 0& 0& \frac{i\gamma (\lambda q_2)}{2}+i\beta & ip_2\end{array}\right)+\stackrel{~}{H}\left(\begin{array}{cccc}0& 0& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& i& 0\end{array}\right),$$
where $`\stackrel{~}{H}=I_1`$ be the Hamilton function (3.9). As in the previous example the spectral curve
$$\left(\mu ^2\frac{\gamma \lambda ^2}{2}+(\beta +I_1)\lambda +\alpha +I_2\right)\left(\mu ^2+\frac{\gamma \lambda ^2}{2}+(\beta I_1)\lambda +\alpha I_2\right)=0,$$
is a product of the corresponding Stäckel curves (2.6).
In the case (b) the Lax matrix is given by
$$\stackrel{~}{}_b=\left(\begin{array}{cccc}\frac{p_1q_1}{\lambda }& 1\frac{q_1^2}{\lambda }& 0& 0\\ & & & \\ \frac{p_1^2(\beta 2\gamma )q_1^2}{\lambda }& \frac{p_1q_1}{\lambda }& 0& 0\\ 0& 0& \frac{p_2q_2}{\lambda }& 1+\frac{q_2^2}{\lambda }\\ & & & \\ 0& 0& \frac{p_2^2+(\beta +2\gamma )q_2^2}{\lambda }& \frac{p_2q_2}{\lambda }\end{array}\right)+\stackrel{~}{H}\left(\begin{array}{cccc}0& 0& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 1& 0\end{array}\right),$$
where $`\stackrel{~}{H}=I_1`$ be the Hamilton function (3.2). The corresponding spectral curve
$$\mathrm{\Gamma }(y,\mu ):det(\stackrel{~}{}_b(\lambda )yI)=0$$
is a product
$$\left(y^2I_1+\frac{2\alpha +I_2}{\lambda }\frac{\beta +2\gamma }{\lambda ^2}\right)\left(y^2I_1+\frac{2\alpha I_2}{\lambda }\frac{\beta 2\gamma }{\lambda ^2}\right)=0,$$
of the initial Stäckel curves (2.3), which could be rewritten in the canonical form (2.6).
All the spectral curves of these $`4\times 4`$ Lax matrices $`_𝒸`$, $`\stackrel{~}{}_k`$ and $`\stackrel{~}{}_b`$ give rise to the quadratic Stäckel integrals $`I_{1,2}`$ (2.1). The third integral $`K`$ (2.7) may be extracted from the same matrices by using multivariable universal enveloping algebras . In fact, this integral is a coefficient of the following multivariable polynomial
$$\text{tr}P_\pi (\lambda _1)(\lambda _2)(\lambda _3)(\lambda _4),$$
(5.4)
where $`P_\pi `$ be permutation operator of auxiliary spaces corresponding to a Young diagram $`\pi `$ . The nonlinear algebra of integrals (2.8) may be reproduced by using the Poisson bracket relations between the Lax matrices $`_j^k(\lambda _j)`$ . The formulae (5.4) for the $`256\times 256`$ matrices has been proved by using the computer algebra system Maple V.
## 6 Conclusion
Let us discuss the list of all the known integrable natural Hamiltonian systems in the plane with a qubic integral . We suppose that all these systems may be embedded into the family of the Stäckel systems , either into the subset of the generalized Stäckel systems or these systems may be related to the Toda lattices and the Calogero-Moser systems . As an example the last case (l) (3.10) of the Drach systems and the Fokas-Lagerstrom model belong to the generalized Stäckel systems . The complete classification will be presented in the forthcoming publication.
In this note we proved this proposition for all the Drach systems . Moreover, we rewrite the qubic integrals for the superintegrable Drach systems in common form (2.7). This generalized ”angular momentum” may be used to construct another $`n`$-dimensional superintegrable Stäckel systems with the cubic integrals of motion. For instance, let us consider the Hamilton function
$$H=\frac{1}{2}(p_x^2+p_y^2+p_z^2)+\frac{\gamma +\delta }{r}+\frac{1}{x^2+y^2}\left(\frac{\alpha (rz)}{r}+\frac{\beta (r+z)}{r}+U(y/x)\right),$$
(6.5)
where $`r=\sqrt{x^2+y^2+z^2}`$ and $`U(y/x)`$ be arbitrary function. The corresponding equations of motion are separable in the parabolic coordinates
$$q_1=r+z,q_2=rz,q_3=\mathrm{arctan}(y/x),$$
which related to the following Stäckel matrix
$$𝐒=\left(\begin{array}{ccc}\mathrm{𝟏}& \mathrm{𝟏}& \mathrm{𝟎}\\ 𝐪_\mathrm{𝟏}^\mathrm{𝟏}& 𝐪_\mathrm{𝟐}^\mathrm{𝟏}& \mathrm{𝟎}\\ 𝐪_\mathrm{𝟏}^\mathrm{𝟐}& 𝐪_\mathrm{𝟐}^\mathrm{𝟐}& \mathrm{𝟒}\end{array}\right).$$
The Hamilton function (6.5) coincides with the Stäckel integral $`I_1`$ (2.1) if
$$U_1=\frac{\alpha }{q_1^2}+\frac{\gamma }{q_1},U_2=\frac{\beta }{q_2^2}+\frac{\delta }{q_2},U_3=U(q_3).$$
Thus we have integrable Stäckel system with the independent integrals of motion $`I_{1,2}`$ and $`I_3`$, which are quadratic polynomials in momenta.
Canonical algebraic curves is defined in variables
$$\mu _1=p_1q_1,\mu _2=p_2q_2,\mu _3=p_3.$$
To substitute these variables in the ”generalized angular momentum” (2.7) one gets additional qubic in momenta integral of motion $`K`$. In the initial physical variables this integral $`K`$ looks like
$`K`$ $`=`$ $`(x^2+y^2)p_z^32zp_z^2(p_xx+p_yy)p_z(p_xx+p_yy)^2`$
$`+`$ $`r^2({\displaystyle \frac{H}{z}}(p_xx+p_yy)+p_z(p_x^2+p_y^2x{\displaystyle \frac{H}{x}}y{\displaystyle \frac{H}{y}}\left)\right).`$
It will be interesting to understand the algebro-geometric origin of this ”generalized angular momentum” (2.7). |
warning/0001/quant-ph0001108.html | ar5iv | text | # A modular functor which is universal for quantum computation
## 1 Introduction
The quantum computer was Feynman’s \[Fey\] last great idea. He understood that local $`\mathrm{`}\mathrm{`}`$quantum gates”, the basis of his model, can efficiently simulate the evolution of any finite dimensional quantum system and by extension any renormalizable system. The details of the argument are given in \[Ll\]. Topological quantum field theories (TQFTs), although possessing a finite dimensional Hilbert space, lack a Hamiltonian—the derivative of time evolution on which the Feynman-Lloyd argument is based. In \[FKW\], we provide a different argument for the poly-local nature of TQFTs showing that quantum computers efficiently simulate these as well. Here we give a converse to this simulation result. The Feynman-Lloyd argument is reversible, so we may summarize the situation as:
(1) finite dimensional quantum systems,
(2) quantum computers (meaning the quantum circuit model QCM \[D\]\[Y\]),
(3) certain topological modular functors (TMFs).
Each can efficiently simulate the others. We wrote TMF above instead of TQFT because we use only the conformal blocks and the action of the mapping class groups on these—not the general morphisms associated to 3-dimensional non-product bordisms.
We would like to thank Alexei Kitaev for conversations on our approach.
## 2 A universal quantum computer
The strictly 2-dimensional part of a TQFT is called a topological modular functor (TMF). The most interesting examples of TMFs are given by the SU(2) Witten-Chern-Simons theory at roots of unity \[Wi\]. These examples are mathematically constructed in \[RT\] using quantum groups (See also \[T\]\[Wa\]). A modular functor assigns to a compact surface $`\mathrm{\Sigma }`$ (with some additional structures detailed below) a complex vector space $`V(\mathrm{\Sigma })`$ and to a diffeomorphism of the surface (preserving structures) a linear map of $`V(\mathrm{\Sigma })`$. In the cases considered here $`V(\mathrm{\Sigma })`$ always has a positive definite Hermitian inner product $`<,>_h`$ and the induced linear maps preserve $`<,>_h`$, i.e. are unitary. The usual additional structures are fixed parameterizations of each boundary component, a labeling of each boundary component by an element of a finite label set $``$ with an involution $`\widehat{}:`$, and a Lagrangian subspace $`L`$ of $`H_1(\mathrm{\Sigma },)`$ (\[T\]\[Wa\]). Since our quantum computer is built from quantum-$`SU(2)`$-invariants of braiding, and the intersection pairing of a planar surface is $`0`$, $`L=H_1(\mathrm{\Sigma };)`$ and can be ignored. The parameterization of boundary components can be dropped. (The essential information which enhances the Kauffman bracket to the Jones polynomial is remembered by the $`\mathrm{`}\mathrm{`}`$blackboard framing” of the braid.) The involution $`\widehat{}`$ is simply the identity since the $`SU(2)`$-theory is self-dual. In fact, we can manage by only considering the $`SU(2)`$-Chern-Simons theory at $`q=e^{\frac{2\pi i}{r}},r=5`$ and so our label set will be the symbols $`\{0,1,2,3\}`$. Note that in our notation, $`0`$ labels the trivial representation, not $`1`$. Since we are suppressing boundary parameterizations, we may work in the disk with $`n`$ marked points-thought of crushed boundary components. Because we only need the $`\mathrm{`}\mathrm{`}`$uncolored theory” to make a universal model, each marked point is assigned the label $`1`$, and the boundary of the disk is assigned the label $`0`$. We consider the action of the braid group $`B(n)`$ which consists of diffeomorphisms of the disk which leave the $`n`$ marked points and the boundary set-wise invariant modulo those isotopic to the identity. The braid group has the well-known presentation:
$$\begin{array}{cc}B(n)=\{\sigma _1,\mathrm{},\sigma _{n1}|& \sigma _i\sigma _j\sigma _i^1\sigma _j^1=id\text{if}|ij|>1\\ & \sigma _i\sigma _j\sigma _i=\sigma _j\sigma _i\sigma _j\text{if}|ij|=1\},\end{array}$$
where $`\sigma _i`$ is the half right twist of the $`i`$-th marked point about the $`i+1`$-st marked point.
To describe our fault-tolerant computational model $`\mathrm{`}\mathrm{`}`$Chern-Simons5” CS5, we must deal with the usual error arising from decoherence as well as a novel $`\mathrm{`}\mathrm{`}`$qubit smearing error” resulting from imbedding the computational qubits within a modular functor super-space. To explain our approaches we initially ignore all errors; in particular formula (1) is a simplification valid only in the error-free context.
The state space $`S_k=(^2)^k`$ of our quantum computer consists of $`k`$ qubits, that is the disjoint union of $`k`$ spin=$`\frac{1}{2}`$ systems which can be described mathematically as the tensor product of $`k`$ copies of the state space $`^2`$ of the basic 2-level system, $`^2=\text{span}(|0>,|1>)`$. For each even integer $`k`$, we will choose an inclusion $`S_k\stackrel{i}{}V(D^2,\text{3k \hspace{0.33em} marked\hspace{0.33em} points})=V(D^2,3k)`$ and show how to use the action of the braid group $`B(3k)`$ on the modular functor to (approximately) induce the action of any poly-local unitary operator $`𝐔:S_kS_k`$. That is we will give an (in principle) efficient procedure for constructing a braid $`b=b(𝐔)`$ so that
$$i𝐔=V(b)i.$$
(1)
To see that this allows us to simulate the QCM, we need to explain: $`(i)`$ what we mean by the hypothesis $`\mathrm{`}\mathrm{`}`$poly-local” on $`𝐔`$, $`(ii)`$ what $`\mathrm{`}\mathrm{`}`$efficient” means, $`(iii)`$ what the effect of the two types of errors are on line (1), and $`(iv)`$ what measurement consists of within our model.
We begin by explaining how to map $`S_k`$ into $`V`$ and how to perform 1 and 2 qubit gates.
Let $`D`$ be the unit 2-dimensional disk and
$$\{\frac{11}{100k},\frac{12}{100k},\frac{13}{100k},\frac{21}{100k},\frac{22}{100k},\frac{23}{100k},\mathrm{},\frac{10k+1}{100k},\frac{10k+2}{100k},\frac{10k+3}{100k}\}$$
be a subset of $`3k`$ marked points on the $`x`$-axis. Without giving formulae the reader should picture $`k`$ disjoint sub-disks $`D_i,1ik`$, each containing one clump of $`3`$ marked points in its interior (these will serve as qubits) and further $`\left(\begin{array}{c}k\\ 2\end{array}\right)`$ disks $`D_{i,j},1i<jk`$, containing $`D_i`$ and $`D_j`$, but with $`D_{ij}D_l=\mathrm{},li\text{or}j`$ (which will allow 2-qubit gates). Strictly speaking, among the larger subdisks, we only need to consider $`D_{i,i+1},1i,i+1<k`$, and could choose a standard (linear) arrangement for these but there is no cost in the exposition to considering all $`D_{i,j}`$ above which will correspond in the model to letting any two qubits interact. Also, curiously, we will see that any of the numerous topologically distinct arrangements for the $`\{D_{i,j}\}`$ within $`D`$ may be selected without prejudice.
We define $`V_k^l`$ to be the $`SU(2)`$ Hilbert space of $`k`$ marked points in the interior with labels equal 1 and $`l`$ label on $`D`$. We need to understand the many ways in which $`V_m^0`$ arises via the $`\mathrm{`}\mathrm{`}`$gluing axiom” (\[Wa\]) from smaller pieces. The axiom provides an isomorphism:
$$V(X_\gamma Y)=_{\text{all consistent labelings}l}V(X,l)V(Y,l),$$
(2)
where the notation suppress all labels not on the 1-manifold $`\gamma `$ along which $`X`$ and $`Y`$ are glued. The sum is over all labelings of the components of $`\gamma `$ satisfying the conditions that matched components have equal labels. According to $`SU(2)`$-Chern-Simons theory \[KL\], for three-punctured spheres with boundary labels $`a,b,c`$, the Hilbert space $`V_{abc}`$ if
$$(i):a+b+c=\text{even},$$
$$(ii):ab+c,ba+b,ca+b\text{(triangle inequalities)}$$
(3)
$$(iii):a+b+c2(r2);$$
and $`V_{abc}=0`$ otherwise. The gluing axiom together with the above information allows an inductive calculation of $`V_k^l`$, where the superscript denotes the label on $`D`$. We easily calculate that
$$\text{dim}V_3^1=2,\text{dim}V_3^3=1,\text{dim}V_6^0=5,\text{dim}V_6^2=8.$$
(4)
Line (4) motivates taking $`V(D_i,\text{its 3 pts},\text{boundary label 1})`$ $`=:V_i^2`$ as our fundamental unit of computation, the qubit. We fix the choice of an arbitrary $`\mathrm{`}\mathrm{`}`$complementary vector” $`v`$ in the state space of $`D\backslash _{i=1}^kD_i`$
$`vV(D\backslash _{i=1}^kD_i`$, all boundary labels $`=1`$ except boundary of $`D=0`$) $`=:V_{\text{complement}}`$ (To keep this space nontrivial, we have taken k even.) Using $`v`$, the gluing axiom defines an injection:
$$i_v:(^2)^k_{i=1}^kV_i\stackrel{v}{}(_{i=1}^kV_i)V_{\text{complement}}\stackrel{\text{as summand}}{}V_{3k}^0$$
(5)
This composition $`i_v`$ determines what we will serve as our computational qubits within the modular functor $`V_{3k}^0`$. The reader familiar with \[FKW\] will notice that we use here a dual approach. In that paper, we imbedded the modular functor into a larger Hilbert space that is a tensor power; here we imbedded a tensor power into the modular functor.
The action of $`B(3)`$ on $`D_i`$ yields 1-qubit gates, whereas two qubit gates will be constructed using the action of $`B(6)`$ on $`D_{i,j}`$. Supposing our quantum computer $`S_k`$ is in state $`s`$, a given $`v`$ as above determines a state $`i_v(s)=svV_{3k}^0`$. Now suppose we wish to evolve $`s`$ by a 2-qubit gate $`g`$ acting unitarily on $`_i^2_j^2`$ and by $`id`$ on $`_l^2,li\text{or}j`$. Using gluing axiom (2) and the inclusion (5), we may write
$$s=\underset{h}{}t_hu_h,$$
(6)
where $`\{t_h\}`$ is a basis or partial basis for $`_i^2_j^2`$ and $`u_h_{li,j}_l^2`$, so $`sv=_h(t_hu_h)v`$. Decomposing along $`\gamma =D_{i,j}`$, we may write $`v=\alpha _0\beta _0+\alpha _2\beta _2`$, where $`\alpha _ϵV(D_{i,j}\backslash (D_iD_j),ϵ\text{on}\gamma )`$, $`ϵ=0`$ or $`2`$ and $`\beta _ϵV(D\backslash (_{li,j}D_lD_{ij}),ϵ\text{on}\gamma `$, and $`0`$ on $`D)`$. Thus
$$sv=\underset{h}{}t_hu_h\alpha _0\beta _0+\underset{h}{}t_hu_h\alpha _2\beta _2,$$
(7)
An element of $`B(6)`$ applied to the 6 marked points in $`D_iD_jD_{ij}`$ acts via a representation $`\rho ^0\rho ^2=:\rho `$ on $`V^0(D_{ij},\text{6 pts})V^2(D_{ij},\text{6 pts}),`$ where the superscript denotes the label appearing when the surface is cut along $`\gamma `$. In particular $`B(6)`$ acts on each factor $`t_h\alpha _0`$ and $`t_h\alpha _2`$ in (7). Note $`t_h\alpha _0`$ belongs to the summand of $`V^0(D_{ij},\text{6 pts})`$ corresponding to boundary labels on $`\left(D_{ij}\backslash (D_iD_j)\right)=0,1,1`$. There is an additional 1-dimensional summand corresponding to boundary labels 0,3,3-with 0,1,3 and 0,3,1 excluded by the triangle inequality $`(ii)`$ in (3) above. Similarly $`t_h\alpha _2`$ belongs to the summand of $`V^2(D_{ij},\text{6 pts})`$ with boundary labels=2,1,1. There are additional summands corresponding to (2,1,3), and (2,3,1) of dimensions 2 each.
Ideally we would find a braid $`b=b(g)B(6)`$ so that $`\rho ^0(b)(t_h\alpha _0)=gt_h\alpha _0`$ and $`\rho ^2(b)(t_h\alpha _2)=gt_h\alpha _2`$. Then referring to (7) we easily check that
$$\rho (b)(sv)=\underset{h}{}\left((gt_h)u_h\right)v,$$
(8)
i.e. $`\rho (b)`$ implements the gate $`g`$ on the state space $`S_k`$ of our quantum computer. In practice there are two issues: $`(i)`$ we cannot control the phase of the output of either $`\rho ^0`$ or $`\rho ^2`$, and $`(ii)`$ these outputs will be only approximations of the desired gate $`g`$. The phase issue $`(i)`$ leads to a change of the complimentary vector $`vv^{}`$ as follows as seen on line $`(\text{9})`$ below. This is harmless since ultimately we only measure the qubits.
$$sv=\underset{h}{}t_hu_h\alpha _0\beta _0+\underset{h}{}t_hu_h\alpha _2\beta _2$$
$``$ gate
$$sv=\omega _0\underset{h}{}gt_hu_h\alpha _0\beta _0+\omega _2\underset{h}{}gt_hu_h\alpha _2\beta _2$$
$$=\underset{h}{}\omega _0gt_hu_h\alpha _0\beta _0+\underset{h}{}\omega _2gt_hu_h\alpha _2\beta _2$$
$$=\underset{h}{}(gt_hu_h)(\omega _0\alpha _0\beta _0+\omega _2\alpha _2\beta _2)$$
$$=:\underset{h}{}(gt_hu_h)v^{}$$
(9)
The approximation issue is addressed by Theorem 2.1 below.
###### Theorem 2.1.
There is a constant $`C>0`$ so that for all unitary $`g:_i^2_j^2_i^2_j^2`$, there is a braid $`b_l`$ of length $`l`$ in the generators $`\sigma _i`$ and their inverses $`\sigma _i^1,1in1`$, so that:
$$\omega _0\rho ^0(b_l)gid_1+\omega _2\rho ^2(b_l)gid_4ϵ(>0)$$
(10)
for some unit complex numbers (phases) $`\omega _i,i=0,2`$ whenever $`ϵ`$ satisfies
$$lC\left(\frac{1}{ϵ}\right)^2.$$
(11)
We use $`||||`$ to denote the operator norms and the subscripts on $`id`$ indicate the dimension of the orthogonal component in which we are trying not to act.
The main work in proving Theorem 2.1 is to show that the closure of the image of the representation $`\rho :B(6)𝐔(5)\times 𝐔(8)`$ contains $`SU(5)\times SU(8)`$. Once this is accomplished the estimate (10) follows with some exponent $`2`$ from \[Ki\] and the refinement to exponent=2 which will appear in \[CN\] following a suggestion of the first author of the present paper. Also as explained in \[Ki\] there is a $`poly(\frac{1}{ϵ})`$ time classical algorithm which can be used to construct the approximating braid $`b_l`$ as a word in $`\{\sigma _i\}`$ and $`\{\sigma _i^1\}`$. The density theorem is the substance of Section 4.
The action $`\rho (b)`$ $`\mathrm{`}\mathrm{`}`$approximately” executes the gate $`g`$ on $`S_k`$ but not in the usual sense of approximation since the state space $`i_v(S_k)`$ itself is only approximately $`g`$ invariant. To convert this $`\mathrm{`}\mathrm{`}`$smearing of qubits” to errors of the type considered in the fault tolerant literature, after each $`g`$ is approximately executed by $`\rho (b)`$ we measure the labels around $`_{i=1}^kD_i`$ to project the new state $`\rho (b)(sv)`$ into the form $`s^{}v,s^{}S_k`$, with probability $`1𝒪(ϵ^2),|s^{}s|𝒪(ϵ)`$. With probability $`𝒪(ϵ^2)`$ the label measurement around $`D_i`$ does not yield one; in this case $`V^1(D_i;\text{3 pts})=:V_{1,1,1,1}^2`$ has collapsed to $`V_{1,1,1,3}`$ and it is as if a qubit has been $`\mathrm{`}\mathrm{`}`$traced out” of our state space. More specifically, if the label $`3`$ is measured on $`D_i`$, we replace $`V^3(D_i,\text{its 3 marked pts})`$ with a freshly cooled qubit $`V^1(D^{},\text{3 pts})`$ with a completely random initial state—an ancilli—which we have been saving for such an occasion. The reader may picture dragging $`D_i`$ off to the edge of the disk $`D`$ and dragging the ancilli $`D_i^{}`$ in as its replacement (and then renaming $`D^{}`$ by $`D_i`$.) The hypothesis that such ancilli are available is discussed below. The error model of \[AB\] is precisely suited to this situation; Aharanov and Ben-Or show in Chapter 8 that a calculation on the level of $`\mathrm{`}\mathrm{`}`$logical” qubits can be kept precisely on track with a probability $`\frac{2}{3}`$ provided the ubiquitous errors at the level of $`\mathrm{`}\mathrm{`}`$physical” qubits are of norm $`𝒪(ϵ)`$ (even if they are systematic and not random) and the large errors (in our case tracing a qubit) have probability also $`𝒪(ϵ)`$ for some threshold constant $`ϵ>0`$. For this, and all other fault tolerant models, entropy must be kept at bay by ensuring a $`\mathrm{`}\mathrm{`}`$cold” stream of ancilli $`|0>`$’s. In the context of our model we must now explain both the role of measurement and ancilla.
Given any essential simple closed curve $`\gamma `$ on a surface $`\mathrm{\Sigma }`$, the gluing formula reads:
$$V(\mathrm{\Sigma })=_lV(\mathrm{\Sigma }_{cut_\gamma },l)$$
(12)
so $`\mathrm{`}\mathrm{`}`$measuring a label” means that we posit for every $`\gamma `$ a Hermitian operator $`H_\gamma `$ with eigenvalues distinguishing the summands of the r.h.s. of (12) above. For a more comprehensive computational study, we would wish to posit that if $`\gamma `$ has length $`=L`$, then $`H_\gamma `$ can be computed in poly(L) time. For the present purpose we only need that $`H_\gamma ,\gamma =D_i\text{or}D_{i,j}`$ can be computed in constant time. Beyond measuring labels, we hypothesize that there is some way of probing the quantum state of the smallest nontrivial building blocks in the theory. For us these are the qubits $`=V_{1,1,1,1}^2`$. Fix a basis $`\{|0>,|1>\}`$ for $`V_{1,1,1,1}`$ and posit for each $`D_i,1ik`$, with label $`1`$ on its boundary, an observable Hermitian operator $`\sigma _z^i:V_{3k}^0V_{3k}^0`$ which acts as the Pauli matrix $`\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right)`$ in a fixed basis $`\{|0>,|1>\}`$ for that qubit. This is our repertoire of measurement: $`H_\gamma `$ is used to $`\mathrm{`}\mathrm{`}`$unsmear physical qubits” after each gate and the $`\sigma _z`$’s to read out the final state (according to von Neumann’s statistical postulate on measurement) after the computation is completed.
In fault tolerant models of computation it is essential to have available a stream of $`\mathrm{`}\mathrm{`}`$freshly cooled” ancilli qubits. If these are present from the start of the computation, even if untouched, they will decohere from errors in employing the identity operator. In the physical realization of a quantum computer unless stored zeros were extremely stable there would have to be some device (inherently not unitary!) for resetting ancilli to $`|0>`$, e.g. a polarizing magnetic field. As a theoretical matter unbounded computation requires such resetting. In a topological model such as $`V(\mathrm{\Sigma })`$ it is not unreasonable to postulate that $`|0>V_{1,1,1,1}=V^1(D_i,\text{3 pts})`$ is stable if not involved in any gates. An alternative hypothesis is that there is some mechanism outside the system analogous to the polarizing magnetic field above which can $`\mathrm{`}\mathrm{`}`$refrigerate” ancilli in the state $`|0>`$ until they are to be used. We refer below to either of these as the $`\mathrm{`}\mathrm{`}`$fresh ancilli” hypothesis. To correct the novel qubit smearing errors, we already encountered the need for ancilli in a random state $`\rho =\left(\begin{array}{cc}\frac{1}{2}& 0\\ 0& \frac{1}{2}\end{array}\right)`$. This state, of course, is easier to maintain.
Let us now return to line (1). Let $`𝐔`$ be the theoretical output of a quantum circuit $`𝒞`$ of (i.e. composition of ) gates to be executed on the physical qubit level so as to fault-tolerantly solve a problem instance of length $`n`$. We assume the problem is in $`BQP`$ and that the above composition has length $``$ $`poly(n)`$. Actually, due to error, $`𝒞`$ will output a completely positive trace preserving super-operator $`𝒪`$, called a physical operator. Now simulate $`𝒞`$ in the modular functor $`V`$ a gate at a time by a succession of braidings and $`H_\gamma `$-measurements. With regard to parallelism (necessary in all fault tolerant schemes), notice that disjoint 2 qubit gates can be performed simultaneously if $`D_{i,j}D_{i^{},j^{}}=\mathrm{}`$. For example this can always be arranged in the linear QCM for gates acting in $`D_{i,i+1}`$ and $`D_{j,j+1}`$ provided $`i+1j,j+1i`$, and $`ij`$, and even this model is shown to be fault tolerant \[AB\]. As noted above, the complementary vector $`vV_{\text{complement}}`$ evolves probabilistically as the simulation progresses . Different $`v`$’s will occur as a tensor factor in a growing number of probabilistically weighted terms. These new $`v`$-values are in the end unimportant; they simply label a computational state (to be observed with some probability) and are never read by the output measurements $`\sigma _z^i`$.
Now the two main theorems:
###### Theorem 2.2.
Let QCM denote the exact quantum circuit model. Suppose $`M`$ is a problem instance in BQP solved by a circuit $`𝒞`$ of length poly(L) where L is length(M). Let CS5 denote the model based on the SU(2)-Chern-Simons modular functor of braids at the fifth root of unity $`e^{\frac{2\pi i}{5}}`$ which we have described in this section: uncolored 3k-strand braids, $`H_\gamma `$ and $`\sigma _z^i`$ measurements, and $`\mathrm{`}\mathrm{`}`$fresh ancilli”. The braid group acts on the modular functor and within the functor one may identify k-qubits $`S_k`$. These actions together with label measurement $`H_\gamma `$’s define a probabilistic evolution of the initial (possibly mixed) state $`\alpha S_k`$. This evolution, defined gate-wise, evolves the mixed state $`\alpha vV_{\text{3k}}`$ of the modular functor to a new (probabilistic mixture of ) state(s) $`\beta `$. Performing $`\sigma _z^i`$-measurements on $`\beta `$ samples from the mixture drawing out a state $`\beta _l=\alpha _lv_l`$ and observing (according to von Neumann measurement) only the $`\alpha _l`$ factor. With probability $`\frac{3}{4}`$ the observations correctly solve the problem instance $`M`$. The number of marked points to be braided (=3k) and the length of the braiding exceed the size of the original circuit $`𝒞`$ by at most a multiplicative poly(log(L)) factor. Taken in triples, they represent the $`\mathrm{`}\mathrm{`}`$physical qubits” of the \[AB\] fault tolerant model, thus CS5 provides a model which efficiently and fault tolerantly simulates the computations of QCM. We note that the use of label measurements $`H_\gamma `$ introduces non-unitary steps in the middle of our simulation.
Proof: The structure of the proof relies heavily on Chapter 8 \[AB\] to reduce the QCM to a linear quantum circuit (with state space $`S_k`$) enjoying a very liberal error model (small systematic errors plus rare trace over qubit). In the final state $`\beta =p_l\beta _l`$, each $`\beta _l`$ admits a tensor decomposition according to the geometry: $`D=(_iD_i)(\text{complement})`$, but along the k boundary components $`_iD_i`$ all choices of labels 1 or 3 may appear. So if we write $`\beta _l=\alpha _lv_l`$ we must remember that associated to $`l`$ is an element $`[l]\{1,3\}^k`$ which defines the subspaces in which $`\alpha _l`$ and $`v_l`$ lie and that $`\beta _l`$ lies in the corresponding $`[l]`$ sector of the modular functor. All occurrences of the label 3 correspond to a $``$ tensor factor, $`V^3(D_i,\text{3 pts})V(D_i,\text{3 pts})`$ whereas the label 1 corresponds to a $`^2`$ factor. Thus in the \[AB\] framework each label 3 corresponds to a $`\mathrm{`}\mathrm{`}`$lost” or averaged qubit according to our replacement procedure $`D_iD^{}`$. Losing an occasional qubit from the computational space $`S_k`$ is the price we pay to $`\mathrm{`}\mathrm{`}`$unsmear” $`S_k`$ within the modular functor. Theorem 2.1 implies that for a braid length $`=𝒪(\frac{1}{ϵ^2})`$ a qubit will be lost with probability $`𝒪(ϵ^2)`$ and if no qubit is lost the gate will be performed with error $`𝒪(ϵ)`$ on pure states. Factoring a mixed state as a probabilistic combination of pure states and passing the error estimate across the probabilities we see that the $`𝒪(ϵ)`$ error bound holds on the super-operator trace norm as well. Thus for $`ϵ`$ sufficiently small (estimated $`<10^6`$ in Chapter 8 \[AB\]), observing (at random) $`\alpha _l`$ amounts to sampling from an error prone implementation of the quantum circuit $`𝒞`$. The error model is not entirely random in that the approximation procedure used to construct $`b_L`$ will have systematic biases. This implies that the $`𝒪(ϵ)`$ errors introduced in the functioning of each gate are not random and must be treated as $`\mathrm{`}\mathrm{`}`$malicious”. Fortunately the error model explained in Chapter 8 \[AB\] permits the small error to be arbitrary as long as the large error, e.g. qubit losses, occurs with a probability dominated by a small constant independent of the qubit and the computational history, as they do in our CS5 model. This completes the proof of Theorem 2.2 modulo the proof of the density Theorem 4.1.
We may define a variant of our model $`\mathrm{𝐂𝐒𝟓}`$, $`\mathrm{`}\mathrm{`}`$exact Chern-Simons at $`e^{\frac{2\pi i}{5}}`$”, $`\mathrm{𝐄𝐂𝐒𝟓}`$, in which we assume that all the braid groups act exactly (no error) on the modular functor $`V`$. Such a hypothesis is not outrageous since a physical implementation of a topological theory may itself confer fault tolerance, in that topological phenomena are inherently discrete. The only difference in the algorithm for modeling the QCM in $`\mathrm{𝐄𝐂𝐒𝟓}`$ is the simplification that $`H_\gamma `$ measurements are not performed in the middle of the simulation, but only at the very end, prior to reading out the qubits $`S_k`$ with $`\sigma _z^k`$ measurements.
###### Theorem 2.3.
There is an efficient and strictly unitary simulation of QCM by $`\mathrm{𝐄𝐂𝐒𝟓}`$. Thus given a problem instance $`M`$ of length $`L`$ in BQP, there is a classical poly(L) time algorithm for constructing a braid $`b`$ as a word of length poly(L) in the generators $`\sigma _i,1i\text{poly(L)=3k}`$. Applying $`b`$ to a standard initial state, $`\psi _{\text{initial}}V^0(D,\text{3k})`$, results in a state $`\psi _{\text{final}}V^0(D,\text{3k})`$, so that the results of $`H_\gamma `$ on $`D_i`$ followed by $`\sigma _z^i`$ measurements on $`\psi _{\text{final}}`$ correctly solve the problem instance $`M`$ with probability $`6`$.
Proof: In the quantum circuit model $`𝒞`$ for $`M`$ (implied by the problem lying in BQP) count the number $`n`$ of gates to be applied. Use line (11) to approximate each gate $`g`$ by a braid $`b`$ of length $`l`$ so that the operator norm error $`\rho (b)g`$ of the approximating gate will be less than $`\frac{1}{10n}`$. The composition of $`n`$ braids which gate-wise simulate the quantum circuit introduces an error on operator norm $`<0.1`$. It follows that our two measurement steps will return an answer (nearly) as reliable as the original quantum circuit $`𝒞`$: $`H_\gamma `$ projects to $`V^1(D,\text{3\hspace{0.33em} pts})`$ with (more than) $`90\%`$ probability and the subsequent probabilities of $`\sigma _z^1`$ measuring $`|0>`$ or $`|1>`$ differ from $`𝒞`$ by less than $`10\%`$.
Remark: Theorem 2.2 and 2.3 are complementary. One provided additional fault tolerance—fault tolerance beyond what might be inherent in a topological model—but at the cost of introducing intermediate non-unitary steps (i.e. measurements). The other eschews intermediate measurements by and so gives a strictly unitary simulation, but cannot confer additional fault tolerance. It is an interesting open problem whether fault tolerance and strictly unitary can be combined in a universal model of computation based on topological modular functors.
## 3 Jones’ representation of the braid groups
A TMF gives a family of representations of the braid groups and mapping class groups. In this section, we identify the representations of the braid groups from the SU(2) modular functor at primitive roots of unity with the irreducible sectors of the representation discovered by Jones whose weighted trace gives the Jones polynomial of the closure link of the braid \[J1\]\[J2\]. To prove universality of the modular functor for quantum computation, we only use this portion of the TMF. Therefore, we will focus on these representations.
First let us describe the Jones representation of the braid groups explicitly following \[We\]. To do so, we need first to describe the representation of the Temperley-Lieb-Jones algebras $`A_{\beta ,n}`$. Fix some integer $`r3`$ and $`q=e^{\frac{2\pi i}{r}}`$. Let $`[k]`$ be the quantum integer defined as $`[k]=\frac{q^{\frac{k}{2}}q^{\frac{k}{2}}}{q^{\frac{1}{2}}q^{\frac{1}{2}}}`$. Note that $`[k]=[k]`$, and $`[2]=q^{\frac{1}{2}}+q^{\frac{1}{2}}`$. Then $`\beta :=[2]^2=q+\overline{q}+2=4cos^2(\frac{\pi }{r})`$. The algebras $`A_{\beta ,n}`$ are the finite dimensional $`C^{}`$algebras generated by $`1`$ and projectors $`e_1,\mathrm{},e_{n1}`$ such that
1. $`e_i^2=e_i`$, and $`e_i^{}=e_i`$,
2. $`e_ie_{i\pm 1}e_i=\beta ^1e_i`$,
3. $`e_ie_j=e_je_i`$ if $`|ij|2`$,
and there exists a positive trace $`tr:_{n=1}^{\mathrm{}}A_{\beta ,n}`$ such that $`tr(xe_n)=\beta ^1tr(x)`$ for all $`xA_{\beta ,n}`$.
The Jones representation of $`A_{\beta ,n}`$ is the representation corresponding to the G.N.S construction with respect to the above trace. An important feature of the Jones representation is that it splits as a direct sum of irreducible representations indexed by some 2-row Young diagrams, which we will refer to as sectors. A Young diagram $`\lambda =[\lambda _1,\mathrm{}\lambda _s],\lambda _1\lambda _2\mathrm{}\lambda _s`$ is called a $`(2,r)`$ diagram if $`s2`$ (at most two rows) and $`\lambda _1\lambda _2r2`$. Let $`_n^{(2,r)}`$ denote all $`(2,r)`$ diagrams with $`n`$ nodes. Given $`\lambda _n^{(2,r)}`$, let $`T_\lambda ^{(2,r)}`$ be all standard tableaus $`\{t\}`$ with shape $`\lambda `$ satisfying the inductive condition which is the analogue of $`(iii)`$ in (3): when $`n,n1,\mathrm{},2,1`$ are deleted from $`t`$ one at a time, each tableau appeared is a tableau for some $`(2,r)`$ Young diagram. The representation of $`A_{\beta ,n}`$ is a direct sum of irreducible representations $`\pi _\lambda ^{(2,r)}`$ over all $`(2,r)`$ Young diagrams $`\lambda `$. The representation $`\pi _\lambda ^{(2,r)}`$ for a fixed $`(2,r)`$ Young diagram $`\lambda `$ is given as follows: let $`V_\lambda ^{(2,r)}`$ be the complex vector space with basis $`\{\stackrel{}{v}_t,tT_\lambda ^{(2,r)}\}`$. Given a generator $`e_i`$ in the Temperley-Lieb-Jones algebra and a standard tableau $`tV_\lambda ^{(2,r)}`$. Suppose $`i`$ appears in $`t`$ in row $`r_1`$ and column $`c_1`$, $`i+1`$ in row $`r_2`$ and column $`c_2`$. Denote by $`d_{t,i}=c_1c_2(r_1r_2)`$, $`\alpha _{t,i}=\frac{[d_{t,i}+1]}{[2][d_{t,i}]}`$, and $`\beta _{t,i}=\sqrt{\alpha _{t,i}(1\alpha _{t,i})}`$. They are both non-negative real numbers and satisfy the equation $`\alpha _{t,i}=\alpha _{t,i}^2+\beta _{t,i}^2`$. Then we define
$$\pi _\lambda ^{(2,r)}(e_i)(\stackrel{}{v}_t)=\alpha _{t,i}\stackrel{}{v}_t+\beta _{t,i}\stackrel{}{v}_{g_i(t)},$$
(13)
where $`g_i(t)`$ is the tableau obtained from $`t`$ by switching $`i`$ and $`i+1`$ if $`g_i(t)`$ is in $`T_\lambda ^{(2,r)}`$. If $`g_i(t)`$ is not in $`T_\lambda ^{(2,r)}`$, then $`\alpha _{t,i}`$ is $`0`$ or $`1`$ given by its defining formula. This can occur in several cases. It follows that $`\pi _\lambda ^{(2,r)}`$ with respect to the basis $`\{\stackrel{}{v}_t\}`$ is a matrix consisting of only of $`2\times 2`$ and $`1\times 1`$ blocks. Furthermore, the $`1\times 1`$ blocks are either $`0`$ or $`1`$, and the $`2\times 2`$ blocks are
$$\left(\begin{array}{cc}\alpha _{t,i}& \beta _{t,i}\\ \beta _{t,i}& 1\alpha _{t,i}\end{array}\right).$$
(14)
The identity $`\alpha _{t,i}=\alpha _{t,i}^2+\beta _{t,i}^2`$ implies that (14) is a projector. So all eigenvalues of $`e_i`$ are either $`0`$ or $`1`$.
The Jones representation of the braid groups is defined by
$$\rho _{\beta ,n}(\sigma _i)=q(1+q)e_i.$$
(15)
Combining (15) with the above representation of the Temperley-Lieb-Jones algebra, we get Jones’ representation of the braid groups, denoted still by $`\rho _{\beta ,n}`$:
$$\rho _{\beta ,n}:B_nA_{\beta ,n}𝐔(N_{\beta ,n}),$$
where the dimension $`N_{\beta ,n}=_{\lambda _n^{(2,r)}}\text{dim}V_\lambda ^{(2,r)}`$ grows asymptotically as $`\beta ^n`$.
When $`|q|=1`$, as we have seen already, Jones’ representation $`\rho _{\beta ,n}`$ is unitary. To verify that $`\rho (\sigma _i)\rho ^{}(\sigma _i)=1`$, note $`\rho ^{}(\sigma _i)=\overline{q}(1+\overline{q})e_i^{}`$. So we have $`\rho (\sigma _i)\rho ^{}(\sigma _i)=q\overline{q}+(1+q)(1+\overline{q})e_ie_i^{}(1+q)e_i(1+\overline{q})e_i^{}=1`$. We use the fact $`e_i^{}=e_i`$ and $`e_i^2=e_i`$ to cancel out the last 3 terms.
From the definition, $`\rho _{\beta ,n}`$ also splits as a direct sum of representations over $`(2,r)`$-Young diagrams. A sector corresponding to a particular Young diagram $`\lambda `$ will be denoted by $`\rho _{\lambda ,\beta ,n}`$.
Now we collect some properties about the Jones representation of the braid groups into the following:
###### Theorem 3.1.
(i) For each $`(2,r)`$-Young diagram $`\lambda `$, the representation $`\rho _{\lambda ,\beta ,n}`$ is irreducible.
(ii) The matrices $`\rho _{\lambda ,\beta ,n}(\sigma _i)`$ for $`i=1,2`$ generate an infinite subgroup of $`𝐔(2)`$ modulo center for $`r3,4,6,10`$.
(iii) Each matrix $`\rho _{\lambda ,\beta ,n}(\sigma _i),1in1,`$ has exactly two distinct eigenvalues $`1,q`$.
(iv) For the (2,5)-Young diagram $`\lambda =[4,2]`$, $`n=6`$, the two eigenvalues $`1,q`$ of every $`\rho _{\lambda ,\beta ,6}(\sigma _i)`$ have multiplicity of 3 and 5, respectively.
The proofs of (i) and (ii) are in \[J2\]. For (iii), first note that the matrix $`\rho _{\lambda ,\beta ,n}(\sigma _1)`$ is a diagonal matrix with respect to the basis $`\{\stackrel{}{v}_t\}`$ with only two distinct eigenvalues $`1,q`$. Now (iii) follows from the fact that all braid generators $`\sigma _i`$ are conjugate to each other. For (iv), simply check the explicit matrix for $`\rho _{\lambda ,\beta ,6}(\sigma _1)`$ at the end of this section.
Now we identify the sectors of the Jones representation with the representations of the braid groups coming from the $`SU(2)`$ Chern-Simons modular functor. The $`SU(2)`$ Chern-Simons modular functor $`\mathrm{𝐂𝐒𝐫}`$ of level $`r`$ has been constructed several times in the literature (for example, \[RT\]\[T\]\[Wa\]\[G\]). Our construction of the modular functor $`\mathrm{𝐂𝐒𝐫}`$ is based on skein theory \[KL\]. The key ingredient is the substitute of Jones-Wenzl idempotents for the intertwiners of the irreducible representations of quantum groups \[RT\]\[T\]\[Wa\]. This is the same $`SU(2)`$ modular functor as constructed using quantum groups in \[RT\] (see \[T\]) which is regarded as a mathematical realization of the Witten-Chern-Simons theory. All formulae we need for skein theory are summarized in Chapter 9 of \[KL\] with appropriate admissible conditions. Fix an integer $`r3`$. Let $`A=\sqrt{1}e^{\frac{2\pi i}{4r}}`$, and $`s=A^2`$, and $`q=A^4`$. (Note the confusion caused by notations. The $`q`$ in \[KL\] is $`A^2`$ which is our $`s`$ here. But in Jones’ representation of the braid groups \[J2\], $`q`$ is $`A^4`$. In all formulae in \[KL\], $`q`$ should be interpreted as $`s`$ in our notation.) The label set $``$ of the modular functor $`\mathrm{𝐂𝐒𝐫}`$ will be $`\{0,1,\mathrm{},r2\}`$ and the involution is the identity. We are interested in a unitary modular functor and the one in \[G\] is not unitary. We claim that if we follow the same construction of \[G\] using our choice of $`A`$ and endow all state spaces of the modular functor with the following Hermitian inner product, the resulting modular functor $`\mathrm{𝐂𝐒𝐫}`$ is unitary.
Given a surface $`\mathrm{\Sigma }`$, a pants decomposition of $`\mathrm{\Sigma }`$ determines a basis of $`V(\mathrm{\Sigma })`$: each basis element is a tensor product of the basis elements of the constituent pants. The desired inner products are determined by axiom $`(2.14)`$ \[Wa\] if we specify an inner product on each space $`V_{abc}`$. Our choice of $`A`$ makes all constants $`S(a)`$ appearing in the axiom $`(2.14)`$ \[Wa\] positive. Consequently, positive definite Hermitian inner products on all spaces $`V_{abc}`$ determine a positive definite Hermitian inner product on $`V(\mathrm{\Sigma })`$. The vector space $`V_{abc}`$ of the three punctured sphere $`P_{abc}`$ is defined to be the skein space of the disk $`D_{abc}`$ enclosed by the seams of the punctured sphere $`P_{abc}`$. The numbering of the three punctures induces a numbering of the three boundary $`\mathrm{`}\mathrm{`}`$points” of the disk $`D_{abc}`$ labeled by $`\{a,b,c\}`$. Suppose $`t`$ is a tangle on $`D_{abc}`$ in the skein space of $`D_{abc}`$, and let $`\overline{t}`$ be the tangle on $`D_{abc}`$ obtained by reflecting the disk $`D_{abc}`$ through the first boundary point and the origin. Then the inner product $`<,>_h:V_{abc}\times V_{abc}`$ is as follows: given two tangles $`s`$ and $`t`$ on $`D_{abc}`$, their product $`<s,t>_h`$ is the Kauffman bracket evaluation of the resulting diagram on $`S^2`$ obtained by gluing the two disks with $`s`$ and $`\overline{t}`$ on them respectively, along their common boundaries with matching numberings. Extending $`<,>_h`$ on the skein space of $`D_{abc}`$ linearly in the first coordinate and conjugate linearly in the second coordinate, we obtain a positive definite Hermitian inner product on $`V_{abc}`$. It is also true that the mapping class groupoid actions in the basic data respect this Hermitian product, and the fusion and scattering matrices $`F`$ and $`S`$ also preserve this product. So $`\mathrm{𝐂𝐒𝐫}`$ is indeed a unitary modular functor.
This modular functor $`\mathrm{𝐂𝐒𝐫}`$ defines representations of the central extension of the mapping class groups of labeled extended surfaces, in particular for $`n`$-punctured disks $`D_n^m`$ with all interior punctures labeled 1 and boundary labeled $`m`$. If $`m1`$, then the mapping class group is the braid group $`B_n`$. If $`m=1`$, then the mapping class group is the spherical braid group $`𝒮B_{n+1}=(0,n+1)`$. Recall that we suppress the issues of framing and central extension as they are inessential in our discussion. Also the representation of the mapping class groups coming from $`\mathrm{𝐂𝐒𝐫}`$ will be denoted simply by $`\rho _r`$.
###### Theorem 3.2.
Let $`D_n^m`$ be as above.
(1): If $`m+n`$ is even, and $`m1`$, then $`\rho _r`$ is equivalent to the irreducible sector of the Jones representation $`\rho _{\lambda ,\beta ,n}`$ for the Young diagram $`\lambda =[\frac{m+n}{2},\frac{mn}{2}]`$ up to phase.
(2): If $`n`$ is odd, and $`m=1`$, then the composition of $`\rho _r`$ with the natural map $`\iota :B_n𝒮B_{n+1}`$ is equivalent to the irreducible sector of the Jones representation $`\rho _{\lambda ,\beta ,n}`$ for the Young diagram $`\lambda =[\frac{n+1}{2},\frac{n1}{2}]`$ up to phase.
The equivalence of these two representations was first established in a non-unitary version \[Fu\]. A computational proof of this theorem can be obtained following \[Fu\]. So we will be content with giving some examples for $`r=5`$.
For the $`(2,5)`$ Young diagram $`\lambda =[2,1]`$, $`n=3`$ with an appropriate ordering of the basis:
$`\rho _{[2,1],\beta ,3}(\sigma _1)=\left(\begin{array}{cc}1& 0\\ 0& q\end{array}\right),`$
$`\rho _{[2,1],\beta ,3}(\sigma _2)=\left(\begin{array}{cc}\frac{q^2}{q+1}& \frac{q\sqrt{[3]}}{q+1}\\ \frac{q\sqrt{[3]}}{q+1}& \frac{1}{q+1}\end{array}\right),`$ where quantum $`[3]=q+\overline{q}+1.`$
For the $`(2,5)`$ Young diagram $`\lambda =[3,3]`$, $`n=6`$, the representation is 5-dimensional. With an appropriate ordering of the basis, we have:
$`\rho _{[3,3],\beta ,6}(\sigma _1)=\left(\begin{array}{ccccc}1& & & & \\ & q& & & \\ & & 1& & \\ & & & q& \\ & & & & q\end{array}\right),`$
$`\rho _{[3,3],\beta ,6}(\sigma _2)=\left(\begin{array}{ccccc}\frac{q^2}{q+1}& \frac{q\sqrt{[3]}}{q+1}& & & \\ \frac{q\sqrt{[3]}}{q+1}& \frac{1}{q+1}& & & \\ & & \frac{q^2}{q+1}& \frac{q\sqrt{[3]}}{q+1}& \\ & & \frac{q\sqrt{[3]}}{q+1}& \frac{1}{q+1}& \\ & & & & q\end{array}\right).`$
For the $`(2,5)`$ Young diagram $`\lambda =[4,2]`$, $`n=6`$, the representation is 8-dimensional. Here the inductive condition on basis elements make one standard tableau illegal, so the representation is not 9-dimensional as it would be if $`r>5`$. This is the restriction analogous to $`(iii)`$ in (3) for the modular functor. With an appropriate ordering of the basis:
$`\rho _{[4,2],\beta ,6}(\sigma _1)=\left(\begin{array}{cccccccc}1& & & & & & & \\ & q& & & & & & \\ & & 1& & & & & \\ & & & q& & & & \\ & & & & 1& & & \\ & & & & & q& & \\ & & & & & & q& \\ & & & & & & & q\end{array}\right).`$
## 4 A Density theorem
In this section, we prove the density theorem.
###### Theorem 4.1.
Let $`\rho :=\rho _{[3,3]}\rho _{[4,2]}:B_6𝐔(5)\times 𝐔(8)`$ be the Jones representation of $`B_6`$ at the $`5`$-th root of unity $`q=e^{\frac{2\pi i}{5}}`$. Then the closure of the image of $`\rho (B_6)`$ in $`𝐔(5)\times 𝐔(8)`$ contains $`SU(5)\times SU(8)`$.
By Theorem 3.2, this is the same representation $`\rho :=\rho ^0\rho ^2:B_6𝐔(5)\times 𝐔(8)`$ in the $`SU(2)`$ Chern-Simons modular functor at the $`5`$-th root of unity used in Section 2 to build a universal quantum computer. In the following, a key fact used is that the image matrix of each braid generator under the Jones representation has exactly two eigenvalues $`\{1,q\}`$ whose ratio is not $`\pm 1`$. This strong restriction allows us to identify both the closed image and its representation.
Proof: First it suffices to show that the images of $`\rho _{[3,3]}`$ and $`\rho _{[4,2]}`$ contain $`SU(5)`$ and $`SU(8)`$, respectively. Supposing so, if $`K=\overline{\rho (B_6)}(SU(5)\times SU(8))`$, then the two projections $`p_1:KSU(5)`$ and $`p_2:KSU(8)`$ are both surjective. Let $`N_2`$ (respectively $`N_1`$) be the kernel of $`p_1`$ (respectively $`p_2`$). Then $`N_1`$ (respectively $`N_2`$) can be identified as a normal subgroup of $`SU(5)`$ (respectively $`SU(8)`$). By Goursat’s Lemma (page 54, \[La\]), the image of $`K`$ in $`SU(5)/N_1\times SU(8)/N_2`$ is the graph of some isomorphism $`SU(5)/N_1SU(8)/N_2`$. As the only nontrivial normal subgroups of $`SU(n)`$ are finite groups, this is possible only if $`N_1=SU(5)`$ and $`N_2=SU(8)`$. Therefore, $`K=SU(5)\times SU(8)`$.
The proofs of the density for $`\rho _{[3,3]}`$ and $`\rho _{[4,2]}`$ are similar. So we prove both cases at the same time and give separate argument for the more complicated case $`\rho _{[4,2]}`$ when necessary.
Let $`G`$ be the closure of the image of $`\rho _{[3,3]}`$ (or $`\rho _{[4,2]}`$) in $`𝐔(5)`$ (or $`𝐔(8)`$) which we will try to identify. By Theorem 3.1, G is a compact subgroup of $`𝐔(m)(m=5\text{or}\mathrm{\hspace{0.33em}8})`$ of positive dimension. Denote by $`V`$ the induced $`m`$-dimensional faithful, irreducible complex representation of $`G`$. The representation $`V`$ is faithful since $`G`$ is a subgroup of $`𝐔(m)`$. Let $`H`$ be the identity component of $`G`$. What we actually show is that the derived group of $`H`$, $`Der(H)=[H,H]`$, is actually $`SU(m)`$. We will divide the proof into several steps.
Claim 1: The restriction of $`V`$ to $`H`$ is an isotypic representation, i.e. a direct sum of several copies of a single irreducible representation of $`H`$.
Proof: As $`G`$ is compact, $`V=_PV_P`$, where $`P`$ runs through some irreducible representations of $`H`$, and $`V_P`$ is the direct sum of all the copies of $`P`$ contained in $`V`$. Since $`H`$ is a normal subgroup, and the braid generators $`\sigma _i`$ topologically generate $`G`$, the $`\sigma _i`$’s permute transitively the isotypic components $`V_P`$ \[CR, Section 49\]. If there is more than 1 such component, then some $`\sigma _i`$ acts nontrivially, so it must permute these blocks.
Now we need a linear algebra lemma:
###### Lemma 4.2.
Suppose $`W`$ is a vector space with a direct sum decomposition $`W=_{i=1}^nW_i`$, and there is a linear automorphism $`T`$ such that $`T:W_iW_{i+1}`$ $`1in`$ cyclically. Then the product of any eigenvalue of $`T`$ with any $`n`$-th root of unity is still an eigenvalue of $`T`$.
Proof: Choose a basis of $`W`$ consisting of bases of $`W_i,i=1,2,\mathrm{},n`$. If $`k`$ is not a multiple of $`n`$, then $`\text{tr}T^k=0`$, as all diagonal entries are 0 with respect to the above basis. Let $`\{\lambda _i\}`$ be all eigenvalues of $`T`$. ( They may repeat.) Consider all values of $`\text{tr}T^m=\lambda _{i}^{}{}_{}{}^{m}(m=1,2,\mathrm{})`$ which are sums of $`m`$-th powers of all eigenvalues of $`T`$. These sums of $`m`$-th powers of $`\{\lambda _i\}`$ are invariant if we simultaneously multiply all the eigenvalues $`\{\lambda _i\}`$ by an $`n`$-th root of unity $`\omega `$: $`(\omega \lambda _i)^m=\omega ^m\lambda _{i}^{}{}_{}{}^{m}=\omega ^m\lambda _{i}^{}{}_{}{}^{m}`$ which is equal to $`\text{tr}T^m=\lambda _{i}^{}{}_{}{}^{m}`$ because when $`m`$ is not a multiple of $`n`$, they are both $`0`$, and when $`m`$ is, $`\omega ^m=1`$. These values $`\text{tr}T^m`$ uniquely determine the eigenvalues of $`T`$, and therefore the set of the eigenvalues of $`T`$ is invariant under multiplication by any $`n`$-th root of unity.
Back to claim 1, if there is more than one isotypic component, then some $`\sigma _i`$ will have an orbit of length at least 2. It is impossible to have an orbit of length 3 or more by the above lemma as this will lead to at least 3 eigenvalues. If the orbit is of length 2 and as $`\rho (\sigma _i)`$ has only two eigenvalues $`\{a,b\}`$, by the lemma, $`\{a,b\}`$ are also eigenvalues. It follows that $`a=b`$ which is impossible when $`q1`$.
Claim 2: The restriction of $`V`$ to $`H`$ is an irreducible representation.
Proof: By claim 1, $`V|_H`$ has only one isotypic component. If $`V|_H`$ is reducible, then the isotypic component is a tensor product $`V_1V_2`$, where $`V_1`$ is the irreducible representation of $`H`$ in the isotypic component and $`V_2`$ is a trivial representation of $`H`$ with $`\text{dim}V_22`$. If $`V_1`$ is 1-dimensional, then $`\rho (\sigma _i),i=1,2`$ generate a finite subgroup of $`𝐔(m)`$ modulo center which is excluded by Theorem 3.1. So we have $`\text{dim}V_12`$. Now we recall a fact in representation theory: a representation of a group $`\rho :GGL(V)`$ is irreducible if and only if the image $`\rho (G)`$ of $`G`$ generates the full matrix algebra $`\text{End}(V)`$. As $`V_1`$ is an irreducible representation of $`H`$, the image $`\rho (H)`$ generates $`\text{End}(V_1)\text{id}_2`$, where the subscript of $`id`$ indicate the tensor factor. As the elements $`\sigma _i`$ normalize $`H`$, they also normalize the subalgebra $`\text{End}(V_1)\text{id}_2`$ in $`\text{End}(V_1V_2)`$. Consequently they act as automorphisms of the full matrix algebra $`\text{End}(V_1)`$. Any automorphism of a full matrix algebra is a conjugation by a matrix, so the braid generators $`\sigma _i`$ act via conjugation (up to a scalar multiple) as invertible matrices in $`\text{End}(V_1)\text{id}_2`$ modulo its centralizer. It is not hard to see the centralizer of $`\text{End}(V_1)\text{id}_2`$ in $`\text{End}(V_1V_2)`$ is $`\text{id}_1\text{End}(V_2)`$. Therefore, the braid generators $`\sigma _i`$ act via conjugation as invertible matrices in $`\text{End}(V_1)\text{End}(V_2)`$, i.e. they preserve the tensor decomposition. This is impossible by the following eigenvalue analysis. Consider a braid generator $`\sigma _i`$, its image $`\rho (\sigma _i)`$ is a tensor product of two matrices each of sizes at least 2. Since $`\rho (\sigma _i)`$ has only two eigenvalues, neither factor matrix can have 3 or more eigenvalues. If both factor matrices have two eigenvalues, the fact that $`\rho (\sigma _i)`$ has 2 eigenvalues in all implies that the ratio of these two eigenvalues is $`\pm 1`$ which is forbidden. If one factor matrix is trivial, then $`\rho (\sigma _i)`$ acts trivially on this factor. As all braid generators are conjugate to each other, so the whole group $`G`$ will act trivially on this factor which implies that $`V`$ is a reducible representation of $`G`$. This case cannot happen either, as $`V`$ is an irreducible representation of $`G`$.
Claim 3: The derived group, $`Der(H)=[H,H]`$, of $`H`$ is a semi-simple Lie group, and the further restriction of $`V`$ to $`Der(H)`$ is still irreducible.
Proof: By claim 2, $`V|_H`$ is a faithful, irreducible representation, so $`H`$ is a reductive Lie group \[V, Theorem 3.16.3\]. It follows that the derived group of $`H`$ is semi-simple. It also follows that the derived group and the center of $`H`$ generate $`H`$. By Schur’s lemma, the center act by scalars. So $`V|_{Der(H)}`$ is still irreducible.
Claim 4: Every outer automorphism of $`Der(H)`$ has order 1, 2, or 3.
First we recall a simple fact in representation theory. If $`V`$ is an irreducible representation of a product group $`G_1\times G_2`$, then $`V`$ splits as an outer tensor product of irreducible representations of $`G_i,i=1,2`$. The restriction of $`V`$ to $`G_1`$ has only one isotypic component, and the restriction of $`V`$ to $`G_2`$ lies in the centralizer of the image of $`G_1`$. So the representation splits.
Proof: It suffices to prove the same statement for the universal covering $`Der^{uc}(H)`$ of $`Der(H)`$, as the automorphism group of $`Der(H)`$ is a subgroup of the automorphism group of $`Der^{uc}(H)`$.
For the 5-dimensional case: as 5 is a prime, $`Der^{uc}(H)`$ is a simple group. It is well-known that any outer automorphism of a simple Lie group is of order 1, 2, or 3.
For the 8-dimensional case, if $`Der^{uc}(H)`$ is a simple group, it can be handled as above, so we need only to consider the split cases. If $`Der^{uc}(H)`$ splits into two simple factors, then one factor must be $`SU(2)`$: of all simply connected simple Lie groups, only $`SU(2)`$ has a 2-dimensional irreducible representation. So the outer automorphism group is either $`Z_2`$ when both factors are $`SU(2)`$, or the same as the outer automorphism group of the other simple factor. Our claim holds. If there are three simple factors, they must all be $`SU(2)`$. The outer automorphism group is the permutation group on three letters $`S_3`$. Again our claim is true.
Claim 5: For each braid generator $`\sigma _i`$, we can choose a corresponding element $`\stackrel{~}{\sigma _i}`$ lying in the derived group $`Der(H)`$ which also has exactly two eigenvalues, whose ratio is not $`\pm 1`$. The multiplicity of each eigenvalue of $`\stackrel{~}{\sigma _i}`$ is the same as that of $`\sigma _i`$. (The choice of $`\stackrel{~}{\sigma _i}`$ is not unique, but its two eigenvalues have ratio $`q`$.)
Proof: Since $`Der(H)`$ is still a normal subgroup of $`G`$, and the braid generators $`\sigma _i`$ normalize $`Der(H)`$, so they determine outer-automorphisms of $`Der(H)`$. By claim 4, an outer-automorphism of $`Der(H)`$ is of order 1, 2, or 3. Hence $`\sigma _i^6`$ acts as an inner automorphism of $`Der(H)`$. By Schur’s lemma, each $`\sigma _i^6`$ is the product of an element in $`Der(H)`$ with a scalar, though the decomposition is not unique. Fix a choice for an element $`\stackrel{~}{\sigma _i}`$ in $`Der(H)`$. Then it has exactly two desired eigenvalues.
To complete the proof of Theorem 4.1, we summarize our situation: we have a nontrivial semi-simple group $`Der^{uc}(H)`$ with an irreducible unitary representation. Furthermore, it has a special element $`x`$ whose image under the representation has exactly two distinct eigenvalues whose ratio is not $`\pm 1`$.
For the 5-dimensional case, $`Der^{uc}(H)`$ is a simple Lie group. Going through the list \[MP\] of pairs $`(G,\varpi )`$, where $`G`$ is a simply connected Lie group and $`\varpi `$ a dominant weight. The only possible 5-dimensional irreducible representations are as follows: rank=1, $`(SU(2),4\varpi _1)`$, rank=2, $`(Sp(4),\varpi _2)`$, and rank=4, $`(SU(5),\varpi _i),i=1,4`$. By examining the possible eigenvalues, we can exclude the first two cases as follows: for the first case, suppose $`\alpha ,\beta `$ are the two eigenvalues of the above element $`x`$ in $`SU(2)`$, then under the representation $`4\varpi _1`$ the eigenvalues of the image of $`x`$ are $`\alpha ^i\beta ^j,i+j=4`$, where $`i`$ and $`j`$ both are non-negative integers. The only possibility is two eigenvalues whose ratio is $`\pm 1`$. For the second case, since $`5`$ is an odd number, any element in the image has a real eigenvalue. Other eigenvalues come in mutually reciprocal pairs. Again the only possibility is two eigenvalues whose ratio is $`\pm 1`$. Therefore, the only possible pair is the third case which gives $`Der^{uc}(H)=SU(5)`$. As $`V`$ is a faithful representation of $`Der(H)`$, the image of $`Der(H)`$ is the same as that of $`Der^{uc}(H)`$ which is $`SU(5)`$.
The 8-dimensional case for $`\rho _{[4,2]}`$ is similar. By \[MP\], we see the possible pairs for simply connected simple groups are $`(SU(2),7\varpi _1)`$, $`(SU(3),\varpi _1+\varpi _2)`$, $`(Spin(7),\varpi _3)`$, $`(Sp(8),\varpi _1)`$, $`(Spin(8),\varpi _i),i=1,3,4`$ and $`(SU(8),\varpi _i),i=1,7`$, where $`\varpi _i`$ is the fundamental weight. The same eigenvalue analysis will exclude all but the $`(SU(8),\varpi _i)`$ case. The proof follows the same pattern as above with the following novelties. Case 2 is the adjoint representation of $`SU(3)`$, if the special element $`xSU(3)`$ has eigenvalues $`\{\alpha ,\beta ,\gamma \}`$, the image matrix of $`x`$ will have eigenvalue $`1`$ with multiplicity 2 and all six pair-wise ratios of $`\{\alpha ,\beta ,\gamma \}`$, so they are $`\pm 1`$. For case 4, recall that if $`\lambda `$ is an eigenvalue of a symplectic matrix, so is $`\lambda ^1`$ with the same multiplicity, thus there are candidates for the special element $`x`$, but all such elements have the property that the multiplicity for both eigenvalues is 4. Notice by Theorem 3.1 (iv), the multiplicity of the two distinct eigenvalue in $`\stackrel{~}{\rho }(\sigma _i)`$ is 3 and 5, respectively. Case 5 is done just as case 4. This excludes all the unwanted simple groups. We have to consider also the product cases. For product of two or three simple factors, the same analysis of eigenvalues as at the end of the proof of claim 2 excludes them. Actually, there are only four cases here: $`SU(2)\times SU(2)`$, $`SU(2)\times SU(4)`$, $`SU(2)\times Sp(4)`$ and $`SU(2)\times SU(2)\times SU(2)`$. This completes the proof of our density theorem. |
warning/0001/hep-lat0001011.html | ar5iv | text | # Determinant of a new fermionic action on a lattice - (I)
## I Introduction
As the Nielsen-Ninomiya theorem states, we necessarily meet the difficulty of so-called fermion doubling problem when we formulate fermion fields on a lattice. In practical calculations, Wilson fermions have been widely used, where an additional term which vanishes in the naive continuum limit is introduced at the expense of the chiral symmetry. An alternative scheme was proposed by Kogut and Susskind . In this scheme the chiral symmetry is maintained as discrete one and doubler fermions are regarded as fermions in other species. In $`(1+D)`$ dimensions, the Kogut-Susskind (KS) formalism describes a theory with $`2^{\frac{1+D}{2}}`$ degenerate quark flavors ($`2^{1+D}`$ components).
Recently it has been shown that lattice fermionic actions satisfying the Ginsparg-Wilson relation may provide a solution of the chirality problem . In these attempts the modified chiral symmetry operator is used in stead of $`\gamma _5`$. The actions are local in the sense that the fermionic matrix are bounded by $`Ce^{\gamma x}`$. However, it has been proved that actions with the Ginsparg-Wilson relation cannot be ”ultralocal” . From the practical point of view, the ultralocality (the couplings drop to zero beyond a finite number of lattice spacings) is also important. Thus ultralocal fermionic actions with better features than, for example, the KS action is awaited, though not a final solution of the chirality problem.
In the recent papers , we proposed a new type of fermionic action on a $`(1+D)`$-dimensional lattice. The action is ultralocal and constructed so that fermion fields satisfy the bosonic type of dispersion relation. In this sense there are no extra poles in the propagator. We found that the minimal number of fermion components, dimensions of the spinor space, is $`2^{D1}`$ in the Minkowski case and $`2^D`$ in the Euclidean case, which should be compared with $`2^{1+D}`$ of the KS fermion. Furthermore our action has the discrete chiral symmetry as well.
It is much of interest to investigate the numerical feasibility of our new fermionic action. When dynamical fermions are included, the property of fermion determinants is crucial in numerical calculations. For example, some methods proposed to treat fermionic freedoms rely on reality or positivity of the fermion determinants .
In this paper we report the analytical and numerical results on the fermion determinants of our new action in $`(1+1)`$ dimensions. Our main concern is on U(1) gauge group, but some results on SU(N) gauge group are also presented. In the case of U(1) gauge group, we will see that calculations with specific conditions for temporal link variables are stable and satisfactory though results without the conditions are unstable. The reason why we need those conditions will be discussed in detail.
In Sec.2, we recapitulate our formalism for later convenience. The analytical and numerical results on fermion determinants will be presented in Sec.3, which is followed by the summary.
## II New fermionic action
In the previous paper , we proposed a new fermionic action on the Euclidean lattice. Though the action respects the discrete chiral symmetry like one in the KS action, fermion fields in this action have $`2^D`$ components in $`(1+D)`$ dimensions, which should be compared with $`2^{1+D}`$ in the case of KS fermions. In this section we briefly sketch our formalism for later convenience.
The action can be written with the fermion matrix $`\mathrm{\Lambda }`$ as
$`S_f={\displaystyle \underset{m,n}{}}\psi _m^{}\mathrm{\Lambda }_{m,n}\psi _n,`$ (1)
where the summation is over lattice points and spinor indices, and our fermion matrix is defined by
$`\mathrm{\Lambda }=1S_0^{}U_E.`$ (2)
Here $`U_E`$ is the Euclidean time evolution operator and $`S_\mu `$ is the unit shift operator defined as
$`S_\mu \psi (x^0,x^1,\mathrm{},x^\mu ,\mathrm{},x^D)=\psi (x^0,x^1,\mathrm{},x^\mu +1,\mathrm{},x^D)(\mu =0,1,\mathrm{},D).`$ (3)
We required that the propagator has no extra poles and found that $`U_E`$ has the form
$`U_E=1{\displaystyle \underset{i=1}{\overset{D}{}}}{\displaystyle \frac{r_E}{2}}\left\{iX_i\left(S_iS_i^{}\right)+\left(1Y_i\right)\left(S_i2+S_i^{}\right)\right\},`$ (4)
where $`r_E`$ is the ratio of the temporal lattice constant to the spatial one and $`X`$’s and $`Y`$’s, which are matrices with respect to spinor indices, should satisfy the following algebra:
$`\{\begin{array}{ccc}\{X_i,X_j\}& =& {\displaystyle \frac{2}{r_E}}\delta _{ij},\hfill \\ \{X_i,Y_j\}& =& 0,\hfill \\ \{Y_i,Y_j\}& =& 2\left({\displaystyle \frac{1}{r_E}}\delta _{ij}+1\right),\hfill \end{array}`$ (8)
where $`i`$ and $`j`$ run from $`1`$ to $`D`$. The matrix $`2(\delta _{ij}/r_E+1)`$ is positive definite for any positive $`r_E`$, therefore $`X`$’s and $`Y`$’s can be assumed hermitian,
$`X_i^{}=X_i,Y_i^{}=Y_i.`$ (9)
The matrices $`X`$’s and $`Y`$’s are written in terms of the Clifford algebra $`\mathrm{\Gamma }_1,\mathrm{\Gamma }_2,\mathrm{},\mathrm{\Gamma }_{2D}`$ as
$`X_i={\displaystyle \frac{\mathrm{\Gamma }_i}{\sqrt{r_E}}},Y_i={\displaystyle \underset{j=1}{\overset{D}{}}}\alpha _{ij}\mathrm{\Gamma }_{D+j},`$ (10)
where
$`\mathrm{\Gamma }_i^{}=\mathrm{\Gamma }_i,\{\mathrm{\Gamma }_i,\mathrm{\Gamma }_j\}=2\delta _{ij}(i,j=1,\mathrm{},2D),`$ (11)
and $`\alpha _{ij}`$ are certain real constants. The dimension of the irreducible representation for $`\mathrm{\Gamma }`$’s is $`2^D`$ and accordingly $`\psi `$ has $`2^D`$ components.
Now we give some useful properties of the time evolution operator $`U_E`$. We can see immediately that
$`U_E^{}=U_E,`$ (12)
as $`X`$’s and $`Y`$’s are hermitian. Since the inverse of $`U_E`$ is given by
$`U_E^1=1{\displaystyle \underset{i=1}{\overset{D}{}}}{\displaystyle \frac{r_E}{2}}\left\{iX_i\left(S_iS_i^{}\right)+\left(1+Y_i\right)\left(S_i2+S_i^{}\right)\right\},`$ (13)
we find that $`U_E`$ is related to its inverse as:
$`\mathrm{\Gamma }_{2D+1}U_E\mathrm{\Gamma }_{2D+1}=U_E^1,`$ (14)
where
$`\mathrm{\Gamma }_{2D+1}`$ $``$ $`i\mathrm{\Gamma }_1,\mathrm{\Gamma }_2\mathrm{}\mathrm{\Gamma }_{2D},`$ (15)
$`\mathrm{\Gamma }_{2D+1}^{}=\mathrm{\Gamma }_{2D+1},\mathrm{\Gamma }_{2D+1}^2=1,`$ $`\{\mathrm{\Gamma }_i,\mathrm{\Gamma }_{2D+1}\}=0(i=1,\mathrm{},2D).`$ (16)
The interaction of the fermion with gauge fields is introduced by replacing the unit shift operators by covariant ones:
$`S_\mu S_\mu (x)U_{x,x+\widehat{\mu }}S_\mu ,`$ (17)
where $`\widehat{\mu }`$ is the unit vector along the $`\mu `$’th direction, and $`U_{x,y}`$ is a link variable connecting sites $`x`$ and $`y`$.
The fermion matrix Eq.(2) and the time evolution operator Eq.(4) become
$`\mathrm{\Lambda }(x)=1S_0^{}(x)U_E(x),`$ (18)
and
$`U_E(x)=1{\displaystyle \underset{i=1}{\overset{D}{}}}{\displaystyle \frac{r_E}{2}}\left\{iX_i\left(S_i(x)S_i^{}(x)\right)+\left(1Y_i\right)\left(S_i(x)2+S_i^{}(x)\right)\right\}.`$ (19)
With gauge fields coupled to the fermions, however, Eqs.(13) and (14) do not hold any more in arbitrary dimensions except for $`(1+1)`$ dimensions. In $`(1+1)`$ dimensions these relations are
$`U_E^1(x)=1{\displaystyle \frac{r_E}{2}}\left\{iX_1\left(S_1(x)S_1^{}(x)\right)+\left(1+Y_1\right)\left(S_1(x)2+S_1^{}(x)\right)\right\},`$ (20)
and
$`\mathrm{\Gamma }_3U_E(x)\mathrm{\Gamma }_3=U_E^1(x),`$ (21)
which are used in the next section.
## III Analytical and numerical results of our fermion determinant in $`(1+1)`$ dimensions
### A Reality of the determinant
In this section we show the analytical and numerical results of our fermion determinant in the $`(1+1)`$-dimensional case.
First, we investigate the determinant of the unit shift operator $`S_0(x)`$ in $`(1+D)`$ dimensions. Let us denote the link variables by $`U_{x,x+\widehat{\mu }}`$ which live on the links connecting two neighboring lattice sites $`x`$ and $`x+\widehat{\mu }`$. Now we consider the case of the Abelian gauge group U(1), and they can be written in the form
$`U_{x,x+\widehat{\mu }}=e^{i\theta _\mu (x)},`$ (22)
where $`\theta _\mu (x)`$ is restricted to the compact domain \[0, 2$`\pi `$).
Then, $`S_0(x)`$ can be constructed as a $`N_TN_𝒙N_s\times N_TN_𝒙N_s`$ matrix where $`N_T`$ and $`N_𝒙`$ is the size of the lattice in the Euclidean time and the spatial directions, respectively, and $`N_s`$ is the number of spinor components. $`S_0(x)`$ is diagonal in the spatial and spinor space and consists of block $`N_T\times N_T`$ matrices $`Q(𝒙)`$ belonging to the Euclidean time space. Thus we can write the determinant of $`S_0(x)`$ as follows:
$`detS_0(x)=\left|\begin{array}{cccccc}Q(1)& & & & & \\ & \mathrm{}& & & & \\ & & Q(2)& & & \\ & & & \mathrm{}& & \\ & & & & Q(N_𝒙)& \\ & & & & & \mathrm{}\end{array}\right|,`$ (29)
where $`Q(𝒙)`$ for each $`𝒙`$ appear $`N_s`$ times as $`S_0(x)`$ is unity in the spinor space.
The determinant of $`Q(𝒙)`$ is easily calculated by noticing its matrix form:
$`detQ(𝒙)`$ $`=`$ $`\left|\begin{array}{ccccc}0& U_{(1,𝒙),(2,𝒙)}& & & \\ & 0& U_{(2,𝒙),(3,𝒙)}& & \\ & & \mathrm{}& \mathrm{}& \\ & & & 0& U_{(N_{T1},𝒙),(N_T,𝒙)}\\ U_{(N_T,𝒙),(1,𝒙)}& & & & 0\end{array}\right|`$ (35)
$`=`$ $`()^{N_T1}U_{(1,𝒙),(2,𝒙)}U_{(2,𝒙),(3,𝒙)}\mathrm{}\left(U_{(N_T,𝒙),(1,𝒙)}\right)`$ (36)
$`=`$ $`()^{N_T}\mathrm{exp}\left\{i{\displaystyle \underset{t=1}{\overset{N_T}{}}}\theta _0(t,𝒙)\right\}.`$ (37)
We used the anti-periodic boundary condition for the time direction, which is represented by a negative sign in the lower-left corner of the matrix $`Q`$.
As a result, in $`(1+D)`$ dimensions we find
$`detS_0(x)`$ $`=`$ $`{\displaystyle \underset{𝒙}{}}\left[()^{N_T}\mathrm{exp}\left\{i{\displaystyle \underset{t=1}{\overset{N_T}{}}}\theta _0(t,𝒙)\right\}\right]^{N_s}`$ (38)
$`=`$ $`\mathrm{exp}\left\{iN_s{\displaystyle \underset{t}{}}{\displaystyle \underset{𝒙}{}}\theta _0(t,𝒙)\right\}=e^{i\alpha },`$ (39)
where the angular variable $`\alpha `$ is defined by
$`\alpha =N_s{\displaystyle \underset{t}{}}{\displaystyle \underset{𝒙}{}}\theta _0(t,𝒙).`$ (40)
From now we assume $`D=1`$ and show that the determinant of the Euclidean time evolution operator $`U_E(x)`$ is unity. By Eq.(21) we easily find $`detU_E(x)=\pm 1`$. Since $`X`$’s and $`Y`$’s obey Eq.(10) and in the limit $`r_E0`$
$`r_EX_i0,r_EY_i0,`$ (41)
$`detU_E(x)`$ is equal to $`1`$ at $`r_E=0`$. From the continuity of $`detU_E(x)`$ with respect to $`r_E0`$, we conclude
$`detU_E(x)=1,`$ (42)
in the $`(1+1)`$-dimensional case.
Collecting the above results, the relations (21) and (42), and the hermiticity of $`U_E(x)`$, we can obtain a relation on the phase of our fermion determinant in $`(1+1)`$ dimensions as follows:
$`\left(det\mathrm{\Lambda }(x)\right)^{}`$ $`=`$ $`det\left(1S_0^{}(x)U_E(x)\right)^{}=det\left(1U_E(x)S_0(x)\right)`$ (43)
$`=`$ $`detU_E(x)detS_0(x)det\left(S_0^{}(x)U_E(x)^11\right)`$ (44)
$`=`$ $`()^{\text{dim. of }\mathrm{\Lambda }}detS_0(x)det\mathrm{\Lambda }(x).`$ (45)
Since the number of components of our fermion is $`N_s=2`$ and consequently the fermion matrix $`\mathrm{\Lambda }(x)`$ always has an even number of dimensions, we get
$`\left(det\mathrm{\Lambda }(x)\right)^{}=e^{i\alpha }det\mathrm{\Lambda }(x).`$ (46)
Accordingly in the case of $`(1+1)`$ dimensions with U(1) gauge fields, our fermion determinant is real under the condition
$`\alpha =2\pi n\text{(}n\text{: integer)},`$ (47)
where $`\alpha `$ is defined in Eq.(40). The condition specified by Eq.(47) will be called ”GT-condition” in this paper, where the temporal link variables are globally constrained as $`_x\theta _0(x)=\pi n`$ and on the infinite lattice the condition is achieved by a gauge transformation. We also define ”T-condition” as $`\theta _0(x)=0`$, which corresponds to the temporal gauge condition on the infinite lattice.
At a glance, it seems strange that the determinant has $`\alpha `$-dependence, because $`\alpha `$ can be gauged away. However, this is not true for a finite lattice, since all gauge transformations on the infinite lattice are not allowed on a finite lattice with periodic boundary conditions. Thus, the dependence comes from finiteness of the lattice. For example, $`_t\theta _0(t,𝒙)`$ is invariant under any gauge transformation on a finite lattice with periodic boundary conditions. The temporal gauge condition is certainly a consistent gauge in an infinite system, but not in a finite system. However, this restriction is just an artifact arising when we try to approximate an infinite system by a finite one. We could think in the following way. Given an infinite system, we impose the temporal gauge condition, which is legitimate. Then we take an finite sub-system to approximate the original total system. The difference is due to size and boundary effects and would disappear in the infinite volume limit. The same is also true for the GT-condition. It should be noted that the fermion determinant is invariant under gauge transformations on a finite lattice, but the determinant does depend on our ”gauge conditions”: the T- or the GT-condition.
Now, we will give the numerical results in the $`(1+1)`$-dimensional U(1) gauge theory with our fermion action. In the following numerical simulations link variables are updated by the Metropolis method and determinants are calculated by the LU decomposition. So there are no systematic errors in the determinants. The way to generate a sequence of configurations in the Monte Carlo computation is as follows. In the T-condition, we fix $`U_0(x)=1`$ everywhere. After this, all temporal link variables are not changed and we only update the link variables of spatial direction using the Metropolis method. In the GT-condition, we prepare a configuration satisfying $`_x\theta _0(x)=\pi n`$. Then a temporal link variable $`e^{i\theta _0(x)}`$ is replaced by $`e^{i\theta _0(x)+i\chi }`$, where $`\chi `$ is a random number between $`\pi `$ and $`\pi `$. At the same time another temporal link variable, which is chosen randomly, is multiplied by $`e^{i\chi }`$. And we use the above configuration as the trial one in a Metropolis acceptance test. It can be shown that this procedure satisfies the micro-reversibility requirement, and therefore also the detailed balance condition.
In Fig.1(a) - 1(c) we show the distribution of the fermion determinants in the complex plane for some $`\beta `$. The statistics is $`2000`$ thermalizations and $`2000`$ measurements. In Fig.1(a), no conditions are imposed on temporal link variables. Fig.1(b) is the result in the T-condition. In Fig.1(c), the GT-condition with $`n=0`$ is imposed, namely the condition $`\alpha =0`$ in Eq.(40) is always kept in Monte Carlo updates. The distribution in Fig.1(a) has a doughnut-like structure with a center around the origin, so that we can expect that the convergence for any observation is very poor. On the other hand, Fig.1(b) and 1(c) show that the determinants in the T- and the GT-condition are real as expected and also positive. The distribution in the imaginary direction is due to numerical errors (note the difference in the scales of real and imaginary parts). The positivity of determinants is important from the numerical point of view and will be discussed later.
To test the convergence in the above three types of conditions, we measured the expectation value of a plaquette value as a function of $`\beta `$ using the following formula:
$`W_p={\displaystyle \frac{{\displaystyle 𝑑U𝑑\psi 𝑑\psi ^{}W_pe^{S_gS_f}}}{{\displaystyle 𝑑U𝑑\psi 𝑑\psi ^{}e^{S_gS_f}}}}={\displaystyle \frac{{\displaystyle 𝑑UW_pdet\mathrm{\Lambda }(x)e^{S_g}}}{{\displaystyle 𝑑Udet\mathrm{\Lambda }(x)e^{S_g}}}}={\displaystyle \frac{W_pdet\mathrm{\Lambda }(x)_0}{det\mathrm{\Lambda }(x)_0}}.`$ (48)
Here $`S_g`$ is the usual action for link fields, and $`d\psi `$, $`d\psi ^{}`$ and $`dU`$ stand for $`_{x,\alpha }d\psi _\alpha (x)`$, $`_{x,\alpha }d\psi _\alpha ^{}(x)`$ and $`_{x,\mu }dU_{x,x+\widehat{\mu }}`$, respectively.
In Fig.2 we can see that the expectation values in the T- and the GT-condition display gentle curves, while one without any conditions is spiky especially for small $`\beta `$. Here, the plotted points are the average over $`10`$ results, each of which is obtained by $`2000`$ Monte Carlo iterations at each $`\beta `$. And the error bars are evaluated by using the standard deviation of the $`10`$ results. If the error bars are not displayed, they are not visible at this scale.
The poorness of the convergence without any conditions can be traced back to the behavior of the denominator in Eq.(48). The expectation value $`det\mathrm{\Lambda }(x)_0`$ not only takes complex values but also suddenly increases or decreases during sampling. Fig.3 shows the absolute value of the averaged fermion determinant as a function of the update iterations for each condition. We find that the convergence in the case of no conditions is ill while in the cases of other two conditions very fine. Generally the more iterations are expected to improve the convergence. Without any conditions, however, this is not the case since the expectation value of the determinant must be in the empty hole of a doughnut-like structure in Fig.1(a). Thus it is difficult to improve the convergence for any expectation value within a reasonable number of iterations.
Now what is the difference between the T- and the GT-condition? We expect that the GT-condition is superior to the T-condition since the former condition is much weaker condition than the latter. To see this, we show in Figs.4(a) and 4(b) the averaged plaquette value in the pure U(1) gauge theory as a function of $`\beta `$ for three types of conditions. We find the line without any conditions and one with the GT-condition are very close to each other, but the line calculated in the T-condition is slightly upper than those. As expected, this difference comes from the size effects in each condition and obviously tends to decrease as the lattice size becomes larger. Thus we conclude that the GT-condition has better feature on a finite lattice: the fermion determinant is positive and the finite-size effects are smaller.
### B Spectrum of the fermion matrix
Before discussing the positivity of the fermion determinant, we study the spectrum of $`\mathrm{\Lambda }(x)`$. First, we introduce a discrete rotational symmetry in the complex plane for eigenvalues of $`S_0^{}(x)U_E(x)`$. Defining $`(T)_{t,t^{}}=\mathrm{exp}(i2\pi t/N_T)\delta _{t,t^{}}`$, we find
$`TS_0^{}(x)U_E(x)T^1=\mathrm{exp}\left(i{\displaystyle \frac{2\pi }{N_T}}\right)S_0^{}(x)U_E(x).`$ (49)
This relation implies that if $`\lambda `$ is some eigenvalue of $`S_0^{}(x)U_E(x)`$, then $`\mathrm{exp}\left(i2\pi /N_T\right)\lambda `$ is also its eigenvalue. Next using the following relation
$`det\left(U_E(x)S_0(x)\lambda ^{}\text{1}\right)`$ $`=`$ $`det\left(U_E(x)S_0(x)\right)det\left(\text{1}\lambda ^{}S_0^{}(x)U_E(x)^1\right)`$ (50)
$`=`$ $`det\left(U_E(x)S_0(x)\right)det\left(\text{1}\lambda ^{}S_0^{}(x)U_E(x)\right),`$ (51)
we rewrite the eigenvalue equation
$`det\left(S_0^{}(x)U_E(x)\lambda \text{1}\right)=0`$ (52)
into the form:
$`det\left(S_0^{}(x)U_E(x)\lambda _{}^{}{}_{}{}^{1}\text{1}\right)=0.`$ (53)
Accordingly, from Eqs.(52) and (53) if $`\lambda `$ is some eigenvalue of $`S_0^{}(x)U_E(x)`$, then $`\lambda _{}^{}{}_{}{}^{1}`$ is also its eigenvalue.
Let us look at the numerical results of the spectrum of our fermion matrix $`\mathrm{\Lambda }(x)=1S_0^{}(x)U_E(x)`$ in the $`(1+1)`$-dimensional U(1) theory. Figs.5(a) - 5(c) display the spectrum in the case of no conditions, the T- and GT-condition, respectively. From the figures we confirm the two properties of the spectrum discussed above for each condition. The spectrum in the GT-condition is so similar to the one with no conditions that one couldn’t distinguish them at a glance. It is interesting that the determinant is real only in the former case.
Now, using above two properties for $`S_0^{}(x)U_E(x)`$, we discuss the positivity of our fermion determinant. We will give a plausible reason for the positivity not a complete proof.
First we write the determinant as follows:
$`det\left(1S_0^{}(x)U_E(x)\right)`$ (54)
$`=\left(1\lambda _1\right)\left(1{\displaystyle \frac{1}{\lambda _{1}^{}{}_{}{}^{}}}\right)\left(1\lambda _2\right)\left(1{\displaystyle \frac{1}{\lambda _{2}^{}{}_{}{}^{}}}\right)\mathrm{}`$ (55)
$`={\displaystyle \frac{1}{()^{N_TN_X}}}{\displaystyle \frac{\left(1\lambda _1\right)\left(1\lambda _{1}^{}{}_{}{}^{}\right)}{\lambda _{1}^{}{}_{}{}^{}}}{\displaystyle \frac{\left(1\lambda _2\right)\left(1\lambda _{2}^{}{}_{}{}^{}\right)}{\lambda _{2}^{}{}_{}{}^{}}}\mathrm{}`$ (56)
In Eq.(56), the denominator must be real, since the numerator is positive and $`det\mathrm{\Lambda }(x)`$ is real under the T- or the GT-condition as shown before. The denominator is a continuous function of the background configuration and furthermore it is never vanished, because $`S_0^{}(x)U_E(x)`$ has always an inverse. Once the denominator is positive for some background configuration, it keeps on having a positive value for any configuration, if one configuration can be transformed continuously to the other within the T- or the GT-condition.
Next, we show that our fermion determinant is positive when all link variables are set unity. In this case $`S_0^{}U_E`$ can be easily diagonalized in the momentum space. The eigenvalues are expressed by the eigenvalues $`\lambda _{U}^{}{}_{m}{}^{}`$ of $`U_E`$ as
$`\lambda _{U}^{}{}_{m}{}^{}\mathrm{exp}\left(i{\displaystyle \frac{2n+1}{N_T}}\pi \right)(n=0,\mathrm{},N_T1).`$ (57)
Therefore
$`det\mathrm{\Lambda }={\displaystyle \underset{m,n}{}}\left(1\lambda _{U}^{}{}_{m}{}^{}\mathrm{exp}\left(i{\displaystyle \frac{2n+1}{N_T}}\pi \right)\right)`$ (58)
is clearly positive, as $`\lambda _{U}^{}{}_{m}{}^{}`$ can be shown positive.
The configuration with all link variables unity clearly satisfies the T- and the GT-condition. This seems to complete our proof for the positivity of the determinant. It should be noticed that we might be faced with configurations where some $`\lambda `$ with $`\left|\lambda \right|=1`$ is not degenerated in spite of the symmetry $`\lambda `$ and $`\lambda _{}^{}{}_{}{}^{1}`$. In this case the relation Eq.(56) does not hold and our proof fails. However, such configurations are very hard to happen.
### C SU(N) gauge group
When the gauge group is SU(N), we can also make similar discussion to the U(1) gauge group case and find the positivity of our fermion determinant. In this case, from Eq.(35) we immediately see
$`detQ(𝒙)`$ $`=`$ $`\pm \left|U_{(1,𝒙),(2,𝒙)}\right|\left|U_{(2,𝒙),(3,𝒙)}\right|\mathrm{}\left|U_{(N_T1,𝒙),(0,𝒙)}\right|`$ (59)
$`=`$ $`\pm 1,`$ (60)
so we find $`detS_0(x)=1`$ without any conditions. Consequently, the determinant of our fermion matrix is always real in the case of SU(N) gauge fields. And we find that our fermion determinant is positive in the $`(1+1)`$-dimensional SU(N), since the two symmetries of the spectrum of $`\mathrm{\Lambda }(x)`$ discussed above are also satisfied. In fact we can confirm the symmetries from the spectra shown in Figs.6(a) and 6(b) for SU(2) and SU(3), respectively. And, in Figs.7(a) and 7(b), we display the numerical results for the distribution of our fermion determinant in the $`(1+1)`$-dimensional SU(2) and SU(3) gauge theories. The results show that fermion determinants are real and positive in both cases.
We conclude that our fermion determinant $`det\mathrm{\Lambda }(x)`$ in (1+1)-dimensions is real and positive for U(1) gauge group under the T- or the GT-condition and for SU(N) gauge group without any conditions, though we have only a plausible reason for the positivity.
## IV Discussion and summary
When we applied our new action without any conditions to U(1) gauge theory on a lattice, we were faced with the problem of convergence in Monte Carlo simulation. In this note, we showed that we could avoid this problem by imposing the T- or the GT-condition. In order to make this situation clear, as an example, we study the propagator of the fermi field,
$`\psi (y)\psi ^{}(y^{})={\displaystyle \frac{{\displaystyle 𝑑\psi 𝑑\psi ^{}𝑑U\psi (y)\psi ^{}(y^{})e^{S_gS_f}}}{{\displaystyle 𝑑\psi 𝑑\psi ^{}𝑑Ue^{S_gS_f}}}}.`$ (61)
Integrating in Eq.(61) with respect to $`\psi (x)`$ and $`\psi ^{}(x)`$ we have
$`\psi (y)\psi ^{}(y^{})={\displaystyle \frac{{\displaystyle 𝑑Udet\left(1S_0^{}(x)U_E(x)\right)\left(1S_0^{}(x)U_E(x)\right)_{y,y^{}}^1e^{S_g}}}{{\displaystyle 𝑑Udet\left(1S_0^{}(x)U_E(x)\right)e^{S_g}}}}.`$ (62)
We introduce new variables $`\mathrm{\Theta }`$ and $`\stackrel{~}{\theta }_0(x)`$ instead of link variables $`U_{x,x+\widehat{0}}=e^{i\theta _0(x)}`$,
$`\theta _0(x)=\mathrm{\Theta }+\stackrel{~}{\theta }_0(x),`$ (63)
and using the following equation
$`{\displaystyle \underset{x}{}d\stackrel{~}{\theta }_0(x)d\mathrm{\Theta }\delta \left(\underset{x}{}\stackrel{~}{\theta }_0(x)\right)}`$
$`={\displaystyle \underset{x}{}d\theta _0(x)d\mathrm{\Theta }\delta \left(\underset{x}{}\theta _0(x)N_TN_X\mathrm{\Theta }\right)}`$
$`={\displaystyle \frac{1}{N_TN_X}}{\displaystyle \underset{x}{}d\theta _0(x)},`$
we get
$`\psi (y)\psi ^{}(y^{})`$ (64)
$`={\displaystyle \frac{{\displaystyle 𝑑\stackrel{~}{U}𝑑\mathrm{\Theta }\delta \left(\underset{x}{}\stackrel{~}{\theta }_0(x)\right)det\left(1e^{i\mathrm{\Theta }}\stackrel{~}{S}_0(x)^{}U_E(x)\right)\left(1e^{i\mathrm{\Theta }}\stackrel{~}{S}_0(x)^{}U_E(x)\right)_{y,y^{}}^1e^{\stackrel{~}{S}_g}}}{{\displaystyle 𝑑\stackrel{~}{U}𝑑\mathrm{\Theta }\delta \left(\underset{x}{}\stackrel{~}{\theta }_0(x)\right)det\left(1e^{i\mathrm{\Theta }}\stackrel{~}{S}_0(x)^{}U_E(x)\right)e^{\stackrel{~}{S}_g}}}},`$ (65)
where the tilde represents replacing $`\theta _0(x)`$ by $`\stackrel{~}{\theta }_0(x)`$.
The phase of determinant $`det(1e^{i\mathrm{\Theta }}\stackrel{~}{S}_0(x)^{}U_E(x))`$ takes any value between $`0`$ and $`2\pi `$ as is shown in Fig.1(a). The summation of $`det(1e^{i\mathrm{\Theta }_n}\stackrel{~}{S}_0(x)^{}U_E(x))`$ over a sequence $`\mathrm{\Theta }_n(n=1,2,\mathrm{})`$ which is chosen at random are canceled out accidentally, then the denominator of Eq.(65) becomes very small. This is the origin of unstable behavior in Monte Carlo simulation. Indeed in Sec.3, under the T- or the GT-condition which ensures that the variable $`\mathrm{\Theta }`$ is fixed zero, we proved the determinant $`det(1S_0^{}(x)U_E(x))`$ in the $`(1+1)`$-dimensional lattice to be real for all configurations of link variables and to be positive for most ones.
We have another reason for fixing $`\mathrm{\Theta }`$, namely, imposing some condition like the T- or the GT-condition. Since the integrand of the numerator in Eq.(65) is equal to the cofactor of matrix $`(1e^{i\mathrm{\Theta }}\stackrel{~}{S}_0(x)^{}U_E(x))`$ whose elements are linear functions of $`e^{i\mathrm{\Theta }}`$, it is a polynomial of $`e^{i\mathrm{\Theta }}`$. After we integrate it with respect to $`\mathrm{\Theta }`$, all terms of this polynomial vanish besides constant terms. For the denominator of Eq.(65) we can say the same thing as the above, thus we have
$`\psi (y)\psi ^{}(y^{})=\text{1}_{y,y^{}},`$ (66)
which is an undesired result. Contrarily, if we choose some condition, like the T- or the GT-condition, we may avoid this trouble.
When we try applying this fermion to the SU(N) lattice gauge theory, we can expect that there is no necessity for imposing some condition because the element like $`e^{i\mathrm{\Theta }}`$ does not belong to the group SU(N). In fact, it is shown that the fermion determinant is real for all configurations and that it is positive for most configurations. And the numerical simulation for SU(2) and SU(3) gauge theory shows that the fermion determinants are real and positive. We will discuss the application of our fermions to SU(N) lattice gauge theory in higher dimensions elsewhere.
## Acknowledgments
We would like to thank Prof. H. Yamamoto for very fruitful discussions in the early stages of this work. |
warning/0001/gr-qc0001031.html | ar5iv | text | # 4d neutral dilatonic black holes and (4+𝑝) dimensional nondilatonic black 𝑝-branes
## Abstract
It is shown that, in contrast to the case of extreme 4d dilatonic black holes, 4d neutral dilatonic black holes with horizon singularities can not be interpreted as nonsingular nondilatonic black $`p`$-branes in ($`4+p`$) dimensions, regardless of the number of extra dimensions $`p`$. That is, extra dimensions do not remove naked singularities of 4d neutral dilatonic black holes.
PACS: 04.50.+h, 04.20.Jb
It has been pointed out by Gibbons, Horowitz, and Townsend that for certain values of the dilaton coupling $`\alpha =\sqrt{p/(p+2)}`$, where $`p`$ is an odd integer, the horizon singularities of extreme 4d dilaton black holes can be removed by extra dimensions. That is, for certain numbers of extra dimensions ($`p=odd`$) the extreme 4d dilaton black hole can be viewed from a $`(4+p)`$ dimensional perspective as a singularity free nondilatonic black $`p`$-brane. It can be inferred that this result may not hold for the case of nonextreme dilaton black holes. In fact, Wesson and de Leon have investigated the neutral spherically symmetric 4d black hole solutions for 5d Kaluza-Klein theory and have shown that the horizon singularities of the neutral nonrotating dilatonic solutions are not removed by the extra dimension, i.e. the 5d nondilatonic neutral black 1-brane is not singularity free, except for the Schwarzschild case. This was done by computing the 5d Kretschmann invariant. The purpose of this paper is to extend this result to an arbitrary number $`p`$ of extra dimensions. This can be accomplished by considering the neutral, spherically symmetric 4d dilaton solutions (in nonisotropic coordinates) of Brans-Dicke (BD) theory that arise from a reduction of the ($`4+p`$) dimensional action for pure gravity. In nonisotropic coordinates the Schwarzschild case is easily identified, and, without appealing to scalar no hair theorems, it can be seen that all other cases arise from singular black $`p`$-branes, regardless of the value of $`p`$.
In ($`4+p`$) dimensions the metric is given by
$$ds^2=\stackrel{~}{g}_{MN}dX^MdX^N=g_{\mu \nu }dx^\mu dx^\nu \delta _{mn}B^2(x)dy^mdy^n,$$
(1)
where $`X^M=(x^\mu ,y^m)`$, with $`M`$,$`N=0,\mathrm{},p+3`$, $`\mu `$,$`\nu =0,\mathrm{},3`$, $`m`$,$`n=1,\mathrm{},p`$, and the metric has a signature diag $`g_{MN}=(+,,\mathrm{},)`$. The action for pure gravity is
$$S_{(4+p)}=\frac{1}{16\pi \overline{G}}d^4xd^py\sqrt{\stackrel{~}{g}}\stackrel{~}{R},$$
(2)
where $`\overline{G}`$ is the gravitation constant for the ($`4+p`$) dimensional theory and $`\stackrel{~}{g}=|det(\stackrel{~}{g}_{MN})|`$. We define the scalar field $`\phi =\frac{1}{G}\left(\frac{B}{B_0}\right)^p`$, where $`B_0`$ is a constant, $`g=detg_{\mu \nu }`$, $`R=g^{\mu \nu }R_{\mu \nu }`$ is the 4d Ricci scalar, and the usual Newtonian gravitation constant $`G`$ is related to $`\overline{G}`$ by $`\overline{G}=GV_p`$, where the “volume” of the compactified internal space $`V_p=B_0^pd^py`$ is assumed to be finite. The action then takes the form
$$S=\frac{1}{16\pi }d^4x\sqrt{g}\left\{\phi R+\frac{\omega }{\phi }g^{\mu \nu }_\mu \phi _\nu \phi \right\},\omega =1+\frac{1}{p}.$$
(3)
This is just the action for Brans-Dicke theory in vacuum with a BD scalar field $`\phi `$ and a fixed BD parameter $`\omega =1+\frac{1}{p}`$. Therefore, vacuum solutions for the Brans-Dicke theory corresponding to the particular value of BD parameter $`\omega =1+\frac{1}{p}`$ are also vacuum solutions for ($`4+p`$) dimensional Einstein gravity, and these solutions can be interpreted as neutral nondilatonic black $`p`$-branes in ($`4+p`$) dimensions.
The static spherically symmetric BD vacuum solutions , have been studied previously by Campanelli and Lousto , Saa Rama , and Kim . Campanelli and Lousto (CL) have represented the solutions in the form
$$\begin{array}{cc}ds^2=A^{m+1}dt^2A^{n1}dr^2r^2A^nd\mathrm{\Omega }^2,& \\ \phi =\phi _0A^{\frac{1}{2}(m+n)},A=\left(1\frac{2r_0}{r}\right),& \end{array}$$
(4)
where $`d\mathrm{\Omega }^2=d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2`$ and $`\phi _0`$, $`r_0`$, $`m`$, and $`n`$ are constants. In terms of these solutions parametrized by $`m`$ and $`n`$, the BD parameter is given by
$$\omega =2\frac{(m^2+n^2+mn+mn)}{(m+n)^2}.$$
(5)
The Schwarzschild solution is obtained for $`m=n=0`$, in which case $`\phi `$, and hence the extra dimensional scale factor $`B`$, is a constant. By examining the 4d Kretschmann curvature invariant $`I=R_{\mu \nu \alpha \beta }R^{\mu \nu \alpha \beta }`$ Campanelli and Lousto have concluded that no physical singularity exists at the surface $`r=2r_0`$ for solutions where either $`n1`$ or $`m=n=0`$. For nonzero values of $`n`$ where $`n>1`$ the surface represents a naked singularity ( in the BD frame defined by the metric and scalar field appearing in the BD action).
However, the $`n1`$ solutions of BD theory are not solutions of the ($`4+p`$) dimensional pure gravity theory. To see this, we define $`s=m+n`$ and use the CL value of the BD parameter in (5), which can then be rewritten in the form
$$s=\frac{(n1)\pm \sqrt{(n1)^22n(n2)(2+\omega )}}{(2+\omega )}.$$
(6)
For $`s`$ to be real valued the quantity $`f=(n1)^22n(n2)(2+\omega )`$ must be nonnegative. For $`1\omega 0`$ we find that the constraint $`f0`$ excludes all values of $`n1`$. Therefore the ($`4+p`$) dimensional Einstein theory does not give rise to any of the 4d BD solutions for which $`n1`$. However, the Schwarzschild solution is allowed. We will make use of these facts shortly.
To determine whether the singularities of the 4d BD solutions at $`r=2r_0`$ can be removed by the extra dimensions, we compute the ($`4+p`$) dimensional Kretschmann scalar $`\stackrel{~}{I}=\stackrel{~}{R}_{ABMN}\stackrel{~}{R}^{ABMN}`$ for the nondilatonic neutral black $`p`$-brane. Using the symmetry properties of the Riemann tensor, $`\stackrel{~}{R}_{ABMN}=\stackrel{~}{R}_{BAMN}=\stackrel{~}{R}_{ABNM}=\stackrel{~}{R}_{MNAB}`$, we can write $`\stackrel{~}{I}=I+4\stackrel{~}{R}^{\alpha b\mu n}\stackrel{~}{R}_{\alpha b\mu n}+\stackrel{~}{R}^{abmn}\stackrel{~}{R}_{abmn}`$, where the four dimensional Kretschmann scalar $`I=R^{\alpha \beta \mu \nu }R_{\alpha \beta \mu \nu }`$ has been computed by CL and can be written in the form
$$I=\frac{4r_0^2}{r^6}A^{(2n+2)}\left\{\left(\frac{r_0}{r}\right)^2J_1(m,n)+4\left(\frac{r_0}{r}\right)J_2(m,n)+6J_3(m,n)\right\},$$
(7)
where the $`J_i(m,n)`$ are polynomial functions of $`m`$ and $`n`$ and are given in . A calculation for $`\stackrel{~}{I}`$ then leads to the result
$$\stackrel{~}{I}=\frac{4r_0^2}{r^6}A^{(2n+2)}(r;m,n),$$
(8)
where
$$(r;m,n)=\left\{\left(\frac{2r_0}{r}\right)^2K_1(m,n)+\left(\frac{2r_0}{r}\right)K_2(m,n)+K_3(m,n)\right\},$$
(9)
with
$$\begin{array}{cc}K_1(m,n)\hfill & =4pq^2\{4n+4q+4\hfill \\ & +\{(q+n)^2(q+n)(n1)+\left(\frac{p1}{2}\right)q^2\hfill \\ & +\frac{1}{4}[(m+1)^2+(n1)^2+2n^2]\}\}+\frac{1}{4}J_1(m,n),\hfill \\ K_2(m,n)\hfill & =4pq^2(4n+4q10)+2J_2(m,n),\hfill \\ K_3(m,n)\hfill & =24pq^2+6J_3(m,n),\hfill \end{array}$$
(10)
and we have defined $`q=\left(\frac{m+n}{2p}\right)`$.
In order for $`\stackrel{~}{I}`$ to be finite on the surface $`r=2r_0`$, we must have that either (1) $`n1`$ or (2) $`n>1`$ and $`(2r_0;m,n)=0`$. As pointed out previously, the solutions satisfying the first condition are inaccessible to the ($`4+p`$) dimensional Einstein theory, so that we must focus on solutions which satisfy the second condition, which is satisfied for $`(r;m,n)A^2`$ with $`n=0`$. This is equivalent to the conditions
$$K_2(m,0)=2K_1(m,0),K_3(m,0)=K_1(m,0).$$
(11)
Using the $`K_i`$ in (10), we find that the only solution satisfying these conditions is the one for which $`m=n=0`$, i.e. the Schwarzschild solution. Therefore, except for the Schwarzschild case, the neutral nonrotating solutions of the 4d dilatonic BD theory cannot be viewed as nonsingular nondilatonic $`p`$-branes in ($`4+p`$) dimensions. Furthermore, the only nondilatonic black $`p`$-brane solutions for which the extra dimensions become visible near $`r=2r_0`$ are those possessing a singularity at $`r=2r_0`$. |
warning/0001/cond-mat0001001.html | ar5iv | text | # Statistical thermodynamics of membrane bending mediated protein-protein attractions
## 1 Introduction
Membrane proteins interact directly via screened electrostatic, van der Waal’s, and hydrophobic forces. These are short ranged, operating typically over distances of less than a nanometer. Proteins can also interact indirectly via the bilayer in which they are dissolved. In particular, a protein that is “geometrically mismatched” to the bilayer will induce deformations that affect neighboring proteins. These “solvent induced forces” (the membrane lipids being the solvent) are generated by bending deformations of the bilayer and typically act over many protein diameters.
Membrane associated proteins can aggregate due to bilayer bending mediated interactions. For example, aquaporin AQP1 and CD59 aggregate to tips of pipette-drawn tubules \[Cho et al., 1999, Discher & Mohandas, 1996\]. Previous studies using a continuum approximation for the intervening bilayer membrane, have treated protein-protein interactions and found an $`r^4`$ repulsion between two identical inclusions \[Goulian, Bruinsma & Pincus, 1993, Kim, et al., 1998, Park & Lubensky, 1996, Dommersnes, Fournier & Galatola, 1998\]. Goulian et al.\[Goulian, Bruinsma & Pincus, 1993\] also find a weak attractive ($`k_BT/r^4`$) interaction arising from Casimir forces resulting from suppressed thermodynamic fluctuations of the intervening membrane. Here, we study in detail a direct mechanical origin for protein-protein attractive interactions. Although bending induced forces between multiple inclusions are not pairwise additive, \[Kim, et al., 1998, Kim, et al., 1999\], in this paper, we restrict ourselves to low protein densities where only pairwise interactions are relevant. We find that the interplay between protein noncircularity \[Kim, et al., 1999\] and background Gaussian curvature curvature dramatically affect protein-protein attractions and thermodynamics.
Many membrane proteins are noncircular in the plane of the membrane, including adsorbed polypeptides such as MARCKS \[Myat et al. 1997\], and bacteriorhodopsin \[Luecke et al. 1999\], which consists of seven transmembrane helices arranged in an elliptical configuration. Small domains, dimers, or droplets of molecules such as cholesterol or specific lipids can themselves behave effectively as membrane inclusions. Droplets need not be rigid to induce attractive forces among themselves; fluctuations in the droplet shape itself may lead to an effective attraction.
In Section 2 we briefly review the mechanical theory of inclusion-induced bilayer bending \[Helfrich, 1973, Kim, et al., 1998, Netz & Pincus, 1995\]. The lipid membrane is approximated by a thin plate that resists out-of-plane bending. Inclusions such as integral membrane proteins, or surface adsorbed molecules, impose boundary conditions along the contact line between the membrane and the protein. Using elastic plate theory to describe the membrane deformations, we derive the energy for two identical inclusions as a function of their relative position within the membrane surface.
In Section 3, we show that the rotational and translational time scales can be separated, so that we can thermally average out the fast rotational degrees of freedom. The resulting effective potential between two proteins is attractive provided the inclusions are sufficiently non-circular. We use the effective potential to compute the second virial coefficient and show how the attractive interactions affect the two-dimensional protein osmotic pressure. Finally, we discuss biological processes where membrane induced long-ranged protein-protein attractions may play an intermediate role, and propose possible measurements.
## 2 Membrane inclusions and height deformation
Small membrane deformations (on the scale of the lipid or protein molecules) can be accurately modeled using standard plate theory \[Landau & Lifshitz, 1985, Helfrich, 1973\]
$$\stackrel{~}{E}[H(𝐒),K(𝐒)]=2b𝑑𝐒H^2(𝐒)+b_g𝑑𝐒K(𝐒),$$
(1)
where $`H(𝐒)`$ and $`K(𝐒)`$ are the local mean and Gaussian curvatures, and $`b`$ and $`b_g`$ are their associated elastic moduli. We have assumed a symmetric bilayer and a vanishing spontaneous mean curvature in the absence of the membrane-deforming proteins of interest. For uniform $`b_g`$, the Gaussian contribution (the second integral in Eq. 1), when integrated over the entire surface yields a constant that is independent of the relative configurations of the embedded proteins \[Kim, et al., 1998, Struik 1994\]. Thus, the Gaussian energy term can be ignored when considering protein-protein interaction energies.
In an expansion of the free energy about that of a flat interface, $`H(𝐒)(1/2)^2h(x,y)`$, where $`^2`$ is the two-dimensional, in-plane Laplacian, and $`h(x,y)`$ is a small, slowly varying height deformation from the flat state (cf. Fig 1). Minimizing $`\stackrel{~}{E}[h(𝐒)]`$ with respect to $`h(x,y𝐒)`$ yields the biharmonic equation
$$^4h(x,y)=2^2H(𝐫)=0.$$
(2)
First consider membrane deformations about an isolated, circularly symmetric inclusion of radius $`a`$. If the bilayer midplane contacts<sup>2</sup><sup>2</sup>2The contact angle $`\gamma `$ incorporates the details of the molecular interactions between the included/adsorbed protein with the lipid molecules. A molecular dynamics simulation would in principle provide the basis for a quantitative estimate of $`\gamma `$, but is beyond the scope of this paper. We will simply assume that $`\gamma `$ is a phenomenological parameter determined by the local chemistry, in complete analogy with the standard liquid-gas-solid contact angle. the protein perimeter $`𝐂`$ (see Fig. 1) at an angle $`\gamma `$, the appropriate solution to Eq. 2 is $`h(r)=\gamma \mathrm{ln}(r/a)`$ for $`r>a`$. We have excluded terms in $`h(r)`$ of the form $`r^2\mathrm{ln}r,r^2,\text{const.}`$ because they are unbounded in energy (Eq. 1), or do not satisfy the contact angle boundary condition at $`r=a`$. Since $`^2\mathrm{ln}(r/a)=2H(r)=0`$ for $`r>a`$, there is no mean curvature bending energy (proportional to $`b`$) residing in the bilayer.
Now consider cases where more than one inclusion are present, or where the contact angles, heights of contact, or the shapes of the membrane associated proteins are noncircular. Three types of noncircularity can arise. The inclusion itself may be noncircular (e.g. elliptical), the height of contact of the bilayer midplane to the inclusion may vary along the perimeter $`𝐂`$ of the protein, and the contact angle itself may vary along $`𝐂`$. These noncircular boundary effects arise from the detailed microscopic nature of the protein and its interaction with the lipid molecules. When more than one protein is present, the deformations surrounding each protein are not circularly symmetric. A nonvanishing mean curvature, $`H(𝐫)`$, that gives bounded bending energies can be represented by a multipole expansion,
$$H(r,\theta )=\underset{n=2}{\overset{\mathrm{}}{}}r^n\left(a_n\mathrm{cos}n\theta +b_n\mathrm{sin}n\theta \right),$$
(3)
where $`(r,\theta )`$ is the radial and angular coordinate about an arbitrary origin. Upon substitution of Eq. 3 into Eq. 1, we find the bending energy $`\stackrel{~}{E}b_{n=2}^{\mathrm{}}(a_n^2+b_n^2)`$. To determine $`a_n,b_n`$, we solve $`_{}^2h(r,\theta )=H(r,\theta )`$ and impose boundary conditions (see Appendix A) on $`h(r,\theta )`$ at $`𝐂`$. In the limit of small noncircularity or low protein concentrations, the largest nondivergent terms are associated with $`n=2`$. Wiggly inclusion cross-sections or highly oscillating boundary conditions only weakly affect membrane bending-mediated protein-protein interactions via $`n>2`$ terms. We derive the full multibody, interaction energy in Appendix A. The two-body interaction energy measured in units of $`k_BT`$ is
$$E(R,\theta _1,\theta _2;\mathrm{\Delta },K_b,\mathrm{\Omega })=\left|\frac{e^{2i\mathrm{\Omega }}}{R^2}+K_b\frac{\mathrm{\Delta }}{2}e^{2i\theta _1}\right|^2+\left|\frac{e^{2i\mathrm{\Omega }}}{R^2}+K_b\frac{\mathrm{\Delta }}{2}e^{2i\theta _2}\right|^2.$$
(4)
The dimensionless separation distance $`R`$, protein ellipticity $`\mathrm{\Delta }`$, and background curvature $`K_b`$ are given by
$$\begin{array}{c}R\frac{r}{R_0},R_0a\sqrt{\gamma }B^{1/4}\hfill \\ \mathrm{\Delta }\overline{\epsilon }\sqrt{B}\hfill \\ K_ba\sqrt{B}\left(\frac{^2h_b(x_1,x_2)}{x_1^2}\right),\hfill \end{array}$$
(5)
where $`B\pi b/k_BT`$ is the dimensionless bending stiffness, and $`\overline{\epsilon }O(\epsilon )`$ quantifies the noncircular nature of each inclusion (see Appendix A). The angle $`\mathrm{\Omega }`$ is measured between the line joining the protein centers and the principle axis of curvature defined by the background Gaussian curvature (see Fig. 2). The angles $`\theta _1,\theta _2`$ are measured between the principle axes of proteins 1 and 2 and the same principle axis. The quantity $`K_b`$ measures the local, externally induced (via other distant proteins or external bending forces) background curvature in this principle axis direction. We show in Appendix A that the dominant effect of distant proteins is to induce mean curvature deformations that decay as $`1/r^2`$, but constant negative Gaussian curvatures. The local curvature $`K_b`$ arises only from deformations that are of zero mean curvature. In what follows, our statistical thermodynamic analyses will be applied to the pair interaction energy given by Eq. 4 with the convention $`\mathrm{\Delta },K_b0`$.
## 3 Rotationally averaged interactions
Proteins that are not attached to the cytoskeleton are free to rotate and diffuse within in the membrane. The interaction potential between two membrane-deforming inclusions is a complicated, nonseparable function of their relative major axis angles and separation distance (cf. Eq. 4). Although the energy is a function of the specific separations and angles between two membrane associated proteins, rotational degrees of freedom are sampled faster than the translational of freedom. This can be shown by the following argument.
A small solvent molecule in solution has a rotational correlation time of the order $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}1`$ ns, while its translational diffusion constant is $`D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}10^6`$cm<sup>2</sup>/s. Therefore, in the time $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}`$ it takes for a small solvent molecule to lose rotational correlation, it would have translated
$$\delta r\sqrt{\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}}0.1\text{nm}.$$
(6)
For membrane constituents, such as bilayer lipid molecules, $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}15`$ns, and $`D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}10^7`$cm<sup>2</sup>/s, where $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}`$ corresponds to rotation about the molecular axis parallel to the normal vector of the membrane \[Marsh, 1990\]. As with small molecules in bulk solution, membrane-bound lipid molecules also move $`\delta R0.1\text{nm}`$ during a rotational correlation time. For larger membrane inclusions such as proteins, both rotation and translational diffusion are slower. If $`\mathrm{\Lambda }`$ is the relative size of the protein radius with respect to the effective lipid radius, protein rotational correlation times increase by $`\mathrm{\Lambda }^3`$ while $`D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}`$ decreases with $`a`$. Membrane proteins that are not too large, $`\mathrm{\Lambda }10`$ say, diffuse $`\delta r1`$nm during the time over which it has lost rotational correlation. Therefore, in the time it takes for a typical inclusion to rotate about its axis, it has diffused less than its own size, typically $``$1nm. This estimate is consistent with fluorescence measurements that find $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}0.11`$ms \[Yamada et al. 1999\]. Rotational time scales for larger proteins may not be much faster than translational motions, therefore, our subsequent model is most appropriate for small, unhindered membrane proteins. We must eventually verify that the protein-protein separation $`r`$ of interest is greater than the typical diffusion distance $`\delta r`$.
The time scale separation can be implemented by statistically averaging over the principle axis angles of the two inclusions while keeping the distance $`R`$ and angle $`\mathrm{\Omega }`$ between them fixed. Weighting the exact two particle energy over its own Boltzmann weight,
$$E_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })=Z^1_0^{2\pi }E(R;\mathrm{\Delta },K_b\theta _1,\theta _2)e^{E(R,\theta _1,\theta _2;\mathrm{\Delta },K_b,\mathrm{\Omega })}𝑑\theta _1𝑑\theta _2,$$
(7)
where the rotational partition function
$$Z_0^{2\pi }e^{E(R,\theta _1,\theta _2;\mathrm{\Delta },K_b,\mathrm{\Omega })}𝑑\theta _1𝑑\theta _2.$$
(8)
Upon substitution of Eq. 4 into Eqs. 7 and 8, and performing the integration (see Appendix B),
$$E_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;K_b,\mathrm{\Omega })=\frac{2\xi ^2}{\mathrm{\Delta }^2}+\frac{\mathrm{\Delta }^2}{2}2\xi \frac{I_1(\xi )}{I_0(\xi )}$$
(9)
where
$$\begin{array}{c}\xi =\mathrm{\Delta }\sqrt{\frac{1}{R^4}+\frac{2K_b}{R^2}\mathrm{cos}2\mathrm{\Omega }+K_b^2}\hfill \end{array}$$
(10)
The effective interaction potential between two inclusions is defined by the difference between the membrane bending energies of two inclusions separated at distance $`R`$ and at infinite separation,
$$\begin{array}{cc}\hfill U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })& =E_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })E_{\mathrm{e}\mathrm{f}\mathrm{f}}(\mathrm{})\hfill \\ & \frac{2\xi ^2}{\mathrm{\Delta }^2}\frac{2\xi I_1(\xi )}{I_0(\xi )}\left[2K_b^22\mathrm{\Delta }K_b\frac{I_1(\mathrm{\Delta }K_b)}{I_0(\mathrm{\Delta }K_b)}\right],\hfill \end{array}$$
(11)
For fixed ellipticity $`\mathrm{\Delta }`$, the set of parameters $`K_b,\mathrm{\Omega }`$ and $`R`$ that gives rise to a minimum at $`R^{}<\mathrm{}`$, if it exists, is implicitly determined by
$$\left(\frac{U_{\mathrm{e}\mathrm{f}\mathrm{f}}}{R}\right)_R^{}=0.$$
(12)
For sufficiently small $`R`$, $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}2/R^4`$, as in the circular protein case.
### 3.1 Zero background curvature
First consider the case of two isolated proteins embedded in a flat membrane. In the absence of external mechanical forces that impose background membrane deformations, and with other inclusions sufficiently far away, $`H_b=K_b=0`$, and $`\xi =|\mathrm{\Delta }|/R^2`$. The effective potential (Eq. 11) becomes
$$U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b=0)=\frac{2}{R^4}\left(\frac{2\mathrm{\Delta }}{R^2}\right)\frac{I_1(\mathrm{\Delta }/R^2)}{I_0(\mathrm{\Delta }/R^2)}.$$
(13)
Without background curvature $`(K_b=0)`$, there are no defining principle axes, and $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ is independent of how the angle of the segment joining the inclusion centers is aligned. Clearly, an effective attractive interaction can arise for $`\mathrm{\Delta }/R^21`$, when $`I_1(\mathrm{\Delta }/R^2)/I_0(\mathrm{\Delta }/R^2)1`$, and $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b=0)1/R^4\mathrm{\Delta }/R^2`$. Although the interaction (Eq. 4) yields both repulsive as well as attractive forces, the Boltzmann thermal average in Eq. 7 favors the lower energy configurations of $`\theta _1,\theta _2`$. Hence the pair of inclusions spends more time in attractive configurations, resulting in a residual attraction in $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)`$. In the $`K_b=0`$ limit, the large $`R`$ behavior of Eq. 13 is
$$U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)=\frac{2\mathrm{\Delta }^2}{R^4}+O(R^6).$$
(14)
Since the potential becomes repulsive at short distances, an effective ellipticity $`\mathrm{\Delta }>\mathrm{\Delta }^{}\sqrt{2}`$ is necessary for the existence of a minimum in $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)`$.
Figure 3a shows the $`\theta `$independent effective interaction potential as a function of $`R`$ for various $`\mathrm{\Delta }`$. As $`\mathrm{\Delta }`$ is increased from $`\mathrm{\Delta }^{}=\sqrt{2}`$, the minimum radius $`R^{}`$ determined by Eq. 12, decreases rapidly from $`R^{}\mathrm{}`$. The $`\mathrm{\Delta }>\mathrm{\Delta }^{}`$ dependence of $`R^{}`$ is plotted in Figure 3b. Also shown are the corresponding magnitudes of the global minima of $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b=0)`$ as a function of $`\mathrm{\Delta }`$.
Figures 3a,b show that appreciable attractive wells can persist at distances $`R1`$. For example, the minimum determined by the set of parameters $`\mathrm{\Delta }2,R^{}0.9`$ and $`|U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R^{})|2(k_BT)`$ can arise for $`\gamma 1,\epsilon /a0.25`$ and $`b30k_BT`$. The separation corresponding to the minimum energy in this case is $`r^{}/a9`$, or nine times the inclusion radius. Our initial assumptions using measured and estimated rotational/translational diffusion constants for typical membrane proteins are validated since $`\delta r/a12r^{}/a`$. We conclude that thermally averaged noncircular membrane deformations can induce long-ranged attractive interactions of at least $`2k_BT`$ up to distances $`10a`$.
### 3.2 Effect of local Gaussian curvature, $`H_b=0,K_b0`$
A local background curvature may arise due to a nonuniform distribution of distant membrane proteins, or an externally imposed deformation. For example, in the experiments of \[Cho et al., 1999, Discher & Mohandas, 1996\], a cell is manipulated by a micropipette. A lipid neck is drawn into the pipette and curvature is externally imposed. Regions near the base of the neck will have a large negative Gaussian curvature. Similarly, membrane fusion and fission processes in endo/exocytosis involves intermediate shapes with constricted necks containing Gaussian curvature. These regions may be “externally” imposed by proteins involved in vesiculation (e.g. dynamin or motor proteins). The Gaussian curvature in this case may also result from lipid structural or composition changes \[Schmidt et al., 1999\]. Therefore, curvature can couple to membrane protein or lipid shapes.
The Gaussian curvature of the membrane between the two proteins establishes local axes of principle curvature such that $`a_{x_1}^2h(x_1,x_2)=a_{x_2}^2h(x_1,x_2)\eta _bK_b>0`$. Since we assume $`H_b=0`$, the background deformation between the two proteins will resemble a saddle with principle curvatures of equal magnitudes (cf. Fig. 2). The rotationally averaged (over $`\theta _1,\theta _2`$) effective interaction $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })`$ will generate attractions at specific orientation angles $`\mathrm{\Omega }`$ even if $`\mathrm{\Delta }<\mathrm{\Delta }^{}`$. This can be most easily seen by expanding equation 11 (the rotationally averaged interaction $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })`$) in powers of $`1/R`$ for large $`R`$:
$$U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R\mathrm{};\mathrm{\Delta },K_b,\mathrm{\Omega })\frac{A_2}{R^2}\mathrm{cos}2\mathrm{\Omega }+\frac{A_4}{R^4}+O\left(R^6\right),$$
(15)
where explicit forms for $`A_2,A_4`$ are given in Appendix A. The appearance of $`A_20`$ when $`K_b>0`$ immediately generates a minimum, however small. Even when ellipticity vanishes ($`\mathrm{\Delta }=0`$), $`A_2K_b\mathrm{cos}2\mathrm{\Omega }<0`$ for appropriate $`\mathrm{\Omega }`$.
The physical origin of attractions in the presence of background curvature can be readily seen by considering Figure 2. With our convention $`\gamma >0`$, circular proteins situated at low regions of the saddle ($`\mathrm{\Omega }\pi /2`$) develop attractive interactions, while those with $`\mathrm{\Omega }0`$ always repel. Recall from previous studies that two circular protein repel with a $`R^4`$ potential \[Goulian, Bruinsma & Pincus, 1993, Kim, et al., 1998, Park & Lubensky, 1996, Dommersnes, Fournier & Galatola, 1998\]. This is a direct consequence of placing a second protein in the Gaussian curvature of the first one. When the background curvature of the membrane in the region between two proteins augments the individual Gaussian curvatures around the first protein (near $`\mathrm{\Omega }=0`$), the $`R^4`$ repulsion is also enhanced. Conversely, if the background curvature mitigates the saddle induced by an individual inclusion (near $`\mathrm{\Omega }=\pi /2`$), the other inclusion sees not only a diminished repulsion, but a mutual attraction at large enough distances. This is because the individual Gaussian curvature around a protein (arising from $`h(r)\gamma \mathrm{ln}(r/a)`$) decays as $`1/r^4`$ and eventually becomes smaller than the imposed constant background Gaussian curvature associated with $`K_b`$. Attractive effects of the background curvature eventually manifest themselves when $`\mathrm{\Omega }\pi /2`$.
Figure 4a shows the effects of a small amount of local background curvature on the effective interaction potential. For small ellipticity $`\mathrm{\Delta }\mathrm{\Delta }^{}`$, minima appear for large enough angles $`\mathrm{\Omega }`$ (approximately for $`\mathrm{\Omega }>\pi /4`$). For similar background curvatures, but much larger ellipticities, the potential develops a repulsive barrier before becoming attractive for certain $`\mathrm{\Omega }`$. This signals that $`A_4<0`$ for large enough $`\mathrm{\Delta }`$ and is depicted in Fig. 4b for $`\mathrm{\Delta }=2.5`$. In the limit of small $`K_b`$, $`A_4<0`$ when
$$\mathrm{\Delta }>\mathrm{\Delta }^{}+\frac{K_b^2}{8}\left(3+\frac{\sqrt{2}}{2}(3+\mathrm{sin}^22\mathrm{\Omega })\right)+O\left(K_b^4\right).$$
(16)
There is yet an additional, qualitatively different feature of $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })`$ when both $`\mathrm{\Delta }`$ and $`K_b`$ are large. Although typical values of $`K_b`$ (see Eq. 5) in biological settings is $`K_b1`$, we find that large values of $`K_b`$ and $`\mathrm{\Delta }`$ give rise to double minima in the interaction potential, especially near $`\mathrm{\Omega }\pi /2`$. Figure 4c shows double minima for $`\mathrm{\Omega }=7\pi /16,\pi /2`$. Additional higher order coefficients such as $`A_6/R^6`$, etc. are required to quantitatively describe multiple minima. The two minima are a consequence of the two independent physical effects that prefer energy minima; local Gaussian curvature associated with $`K_b`$ and effective ellipticity $`\mathrm{\Delta }`$. Typically, the weaker, longer-ranged minimum is predominantly the signature of a large $`K_b`$, while the deeper, shorter-ranged minimum (such as that shown in Fig. 3a and 4b) is a feature of ellipticity $`\mathrm{\Delta }>\mathrm{\Delta }^{}`$. Saddles of order $`K_b>1`$ correspond to principle radii of curvature on the order of $`10`$ times the protein size $`a`$, and are thus regions of extreme Gaussian deformations. Regions of such warp may would only exist as transient, small systems such as fusion necks. Henceforth, we will restrict ourselves to $`K_b`$ small enough to only induce one minimum.
Angles $`\mathrm{\Omega }`$ which yield attractive interactions can be estimated by considering $`A_2,A_4`$. Assuming $`A_4>0`$, values of $`A_2<0`$ give attractive interactions when $`\pi /4<\mathrm{\Omega }<\pi /4`$. When $`A_2>0`$, proteins with orientation $`\pi /4<|\mathrm{\Omega }|<3\pi /4`$ will experience attractive forces. However, these conditions are modified if $`A_4<0`$, when some angles within $`\pi /4<\mathrm{\Omega }<\pi /4`$ can yield attraction even if $`A_2>0`$. This case corresponds to Fig. 4b where a repulsive barrier at $`R>R^{}`$ arises. A minimum can still arise even at angles where $`A_2\mathrm{cos}2\mathrm{\Omega }>0`$ due to the $`R^4`$ behavior. The matching to repulsive behavior at smaller $`R`$ requires consideration of $`+R^6`$ terms.
The top panel of Figure 5 shows the radius corresponding to the only minimum of the effective potential $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ as a function of $`K_b`$, for $`\mathrm{\Delta }=0.5,2`$ and $`4`$. Both east-west and north-south configurations are shown, with intermediate angles $`\mathrm{\Omega }`$ interpolating between the curves. For small ellipticity, the local principle curvature $`K_b`$ is the predominant source of attraction at larger distances, shown by the thick dashed curve. Increasing $`K_b`$ destabilizes the effective energy minima near $`\mathrm{\Omega }=0`$. Above a certain background Gaussian curvature intensity, the effective potential minimum evaporates to $`R^{}\mathrm{}`$ for proteins situated at $`\mathrm{\Omega }=0`$ (solid curves), and the attraction is washed out. For small $`K_b`$, the two effects, ellipticity and background Gaussian curvature, complement each other near $`\mathrm{\Omega }=\pi /2`$ in reinforcing an energy minimum. Consistent with Fig. 3a for $`\mathrm{\Delta }>\sqrt{2}`$, $`R^{}`$ in Fig. 5 (thick curves) is smaller for larger $`\mathrm{\Delta }`$. The bottom panel plots the corresponding minimum energies.
The $`\mathrm{\Omega }`$-dependence of $`R^{}`$ and the minimum energy is shown in Figure 6. As expected, or large $`\mathrm{\Delta }\sqrt{2}`$, both $`R^{}`$ and $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R^{},\mathrm{\Omega })`$ are fairly insensitive to $`\mathrm{\Omega }`$. When $`\mathrm{\Delta }`$ is small, the energy minima and their associated radii $`R^{}`$, caused predominantly by $`K_b`$, are very sensitive to orientation $`\mathrm{\Omega }`$. These behaviors are consistent with the energy profiles shown in Fig. 4b. In fact, for small enough $`\mathrm{\Delta }`$, the minima near $`\mathrm{\Omega }0`$ are annihilated, independent of $`K_b`$. Thus, we see a qualitative difference between attractive potentials generated by intrinsic ellipticity and background Gaussian curvature.
## 4 The second virial coefficient
We now consider the influence of the effective protein-protein attractions on a low density ensemble of inclusions. By analogy with the molecular origins of the osmotic second virial coefficients of proteins in solution \[Neal, Asthagiri & Lenhoff, 1998\], we will consider the bending energy contributions to the second virial coefficient for a two-dimensional protein equation of state. The membrane mediated interactions however, are much longer-ranged than those in solution \[Neal, Asthagiri & Lenhoff, 1998\]. Consider the thermodynamic limit and times long enough such that
$$\tau \frac{\mathrm{}^2}{D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}}\tau _{\mathrm{r}\mathrm{o}\mathrm{t}},$$
(17)
where $`D_{\mathrm{t}\mathrm{r}\mathrm{a}\mathrm{n}\mathrm{s}}`$ is the protein translational diffusion constant. On the time scale $`\tau `$, the inclusions are relatively free to diffuse about the bilayer. They interact among themselves via the rotationally averaged potential $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ that manifests itself on time scales $`\tau _{\mathrm{r}\mathrm{o}\mathrm{t}}`$. For very low protein densities (large protein separation $`\mathrm{}`$), the two dimensional protein osmotic pressure will be nearly that of an ideal gas, analogous to a low density gaseous phase surfactant monolayer at the air water interface. Finite protein size $`a`$, and longer-ranged elastically-coupled interactions will give nonideal gas properties. The first correction to ideality in the equation of state is given by the second virial coefficient \[McQuarrie, 1976\]:
$$\frac{\mathrm{\Pi }}{k_BT}=\mathrm{\Gamma }+B_2\mathrm{\Gamma }^2+O\left(B_3\mathrm{\Gamma }^3\right),$$
(18)
where $`\mathrm{\Gamma }`$ is the surface concentration of protein and $`B_2`$ is computed using the formula
$$\begin{array}{cc}\hfill B_2(\mathrm{\Delta },K_b)& \frac{1}{2Z}_0^{2\pi }_0^{\mathrm{}}\left(e^{E(R,\theta _1,\theta _2;\mathrm{\Delta },K_b,\mathrm{\Omega })+E_{\mathrm{e}\mathrm{f}\mathrm{f}}(\mathrm{})}1\right)R𝑑R𝑑\mathrm{\Omega }𝑑\theta _1𝑑\theta _2\hfill \\ & =\frac{1}{2}_0^{2\pi }_0^{\mathrm{}}\left(e^{U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })+U_{\mathrm{}}(\mathrm{\Delta },K_b)}1\right)R𝑑R𝑑\mathrm{\Omega },\hfill \end{array}$$
(19)
The second virial $`B_2`$ represents the small fraction of pairwise interacting proteins. Equations 18 and 19 are nondimensionalized such that the surface density $`\mathrm{\Gamma }1`$ is measured by the number of proteins in area $`R_0^2`$ (see Eq. 5) and the protein osmotic pressure $`\mathrm{\Pi }`$ is measured in units of $`k_BT/R_0^2`$. Equation 19 is exact and does not require the separation of rotational and translational diffusion times needed for the derivation of $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R;\mathrm{\Delta },K_b,\mathrm{\Omega })`$. Here, we do not consider how integrating out the rotational degrees of freedom affect the fixed translational degree of freedom. Instead, we are considering times long enough for equilibration of both degrees of freedom, and their combined contribution to the equation of state via $`B_2`$.
The physical origin and value of $`K_b`$ used in Eq. 19 is as follows. The local curvature felt by the interacting pair represents an interaction between this pair and some other distant proteins. However, the virial equation of state (Eq. 18) is a systematic expansion in surface density expanded about an ideal, noninteracting ensemble. Since membrane bending mediated interactions are not pairwise additive \[Kim, et al., 1998\], one might be tempted to assume that the presence of other proteins would modify the interaction energy $`E`$ used in the expression for $`B_2`$. However, these more complicated interactions would depend upon the concentration of the other background proteins, and would generate terms higher order in $`\mathrm{\Gamma }`$. In other words, we start at densities so low that the protein ensemble is completely noninteracting. As the density is slightly increased, a pair of protein molecules occasionally interact and perhaps form dimers, with each pair ignorant of any other protein. At this still rather low density, the probability three or more proteins approach each other is negligible. When the density is further increased, one needs to consider the higher order virial terms. Therefore, to second order in $`\mathrm{\Gamma }`$, the deviation of the equation of state from ideality is completely determined by the two-body interaction $`E(R,\theta _1,\theta _2;\mathrm{\Delta },\overline{K}_b,\mathrm{\Omega })`$ and is independent of nonpairwise effects \[McQuarrie, 1976\]. Note however, that the two-body interaction will depend on the $`\overline{K}_b`$ associated with externally forced, zero mean curvature membrane deformations. Therefore, for the expansion Eq. 18 to be consistent, the value of $`K_b=\overline{K}_b`$ to be used in Eq. 19 is that owing solely to external force-generated Gaussian curvatures, independent of the protein density.
Figure 7a shows the numerically computed second virial coefficient as a function of inclusion ellipticity for various $`\overline{K}_b`$. As expected for small $`\overline{K}_b`$, the virial coefficient becomes increasingly negative as the ellipticity increases. The value for circular inclusions $`B_2(\mathrm{\Delta }=0,\overline{K}_b=0)=\pi ^{3/2}/\sqrt{2}`$ corresponds to purely repulsive disks with mutual interaction $`U(R)=2/R^4`$. At ellipticity $`\mathrm{\Delta }1.69`$, $`B_2(1.69,\overline{K}_b=0)0`$ corresponding to a protein solution that is ideal to second order in surface density. Although when $`\mathrm{\Delta }1.69>\mathrm{\Delta }^{}=\sqrt{2}`$, $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ has an attractive minimum, its effects are negated by the repulsive $`R^4`$ part of the interaction such that the overall, effective contribution to $`B_2`$ vanishes. For $`\mathrm{\Delta }>1.69`$, the effective attraction between membrane proteins begins to manifest itself. The second virial is modified by externally imposed Gaussian curvature. Recall that when $`\overline{K}_b0`$, certain angles $`\mathrm{\Omega }`$ lead to attractive interactions, even for small $`\mathrm{\Delta }<\mathrm{\Delta }^{}`$. Since we are now thermodynamically averaging over protein positions and $`\mathrm{\Omega }`$ in addition to $`\theta _1,\theta _2`$, the inclusions will spend more time at attractive, lower energy angles $`\mathrm{\Omega }`$, hence lowering $`B_2`$. Consistent with Fig. 4, larger values of $`\overline{K}_b`$ for $`\mathrm{\Delta }>\mathrm{\Delta }^{}`$ lead to stronger repulsions at small $`\mathrm{\Omega }`$, which average into $`B_2`$, making it less negative.
The dependence of $`B_2`$ on $`\overline{K}_b`$ is indicated in Figure 7b. In the absence of ellipticity, $`B_2`$ is given by the integral
$$B_2(\mathrm{\Delta }=0;\overline{K}_b)=\pi _0^{\mathrm{}}\left[e^{2/R^4}I_0(4\overline{K}_b/R^2)1\right]R𝑑R.$$
(20)
For $`\mathrm{\Delta }>0`$, $`B_2`$, found from numerical integration of the full expression Eq. 19, are also shown in Fig. 7b. For $`\overline{K}_b=0`$, increasing ellipticity decreases inclusion repulsions and $`B_2`$. As in Fig. 7a, large $`\overline{K}_b`$ and $`\mathrm{\Delta }`$ tend to increase $`B_2`$.
Equation 4 was used in Eq. 19 to compute the curves shown in Figures 7; thus, the protein-protein interaction was assumed to consist of contributions only from membrane bending. The hard core, excluded area of each protein, $`\pi a^2`$, can be included by modifying $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)`$ by setting $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(Ra/R_0)=\mathrm{}`$. Although we expect this additional repulsive term to further reduce the effective sampling area of the inclusions, and increase the second virial coefficient, we find that for all reasonable values of $`R_0`$, $`B_2`$ does not change noticeably from those shown in Figs. 7. The hard core part of the potential, due to e.g. close-ranged van der Waals repulsion, is not statistically sampled by the inclusions since the membrane bending induced interactions ($`1/R^4`$) already keeps them far apart.
Since nonpairwise interactions manifest themselves only at third and higher order in $`\mathrm{\Gamma }`$, we can estimate their importance by comparing $`B_2\mathrm{\Gamma }^2`$ with $`B_3\mathrm{\Gamma }^3`$. For nonpairwise interactions to be thermodynamically relevant it is necessary but not sufficient that the surface density
$$\mathrm{\Gamma }\left|\frac{B_2}{B_3}\right|.$$
(21)
Although multibody interactions may be important microscopically, their effects on the low density equation of state, cannot be resolved. Even if the density is high enough for $`B_3\mathrm{\Gamma }^3`$ to be measurable, the value of $`B_3`$ is found via a four-dimensional integral over configurations of three membrane proteins. All orientations and distances will be averaged and all components of their interactions, repulsive, attractive, pairwise, and nonpairwise will be included. In other words, one cannot uniquely determine the potential $`U`$ from a measurement of $`B_n`$.
## 5 Discussion and Conclusions
Proteins beyond the range of screened electrostatic, or van der Waals molecular forces can exert forces on one another by virtue of the deformation they impose on the lipid bilayer. These interactions can be attractive if the proteins have a noncircular cross sectional shape or if the local membrane deformation is saddle shaped (negative Gaussian curvature). For bending rigidities $`b30k_BT`$, and protein shape ellipticities $`\epsilon /a0.25`$, we find attractive interactions of a few $`k_BT`$ acting at a range of $`5`$ protein diameters. Thus proteins of radii $`5`$nm can interact at distances of $`50`$ nm, much further than any interaction between similar molecules in solution. On a flat membrane ($`H_b=K_b=0`$), an effective ellipticity $`\overline{\epsilon }>(2k_BT/\pi b)^{1/2}`$ is necessary for a potential minimum to emerge between a pair of proteins.
Although we have presented the model in terms of integral membrane proteins, noncircular peripheral proteins, or lipid/cholesterol/proteins aggregates can also induce local membrane bending. Elastically-coupled interactions between peripheral membrane proteins can mediate dissociation and binding. Dimerization of noncircular peripheral proteins lowers their absolute energy below that of separated ones, and so they are less likely to dissociate from the membrane. Similarly, in one dimension, adsorbed proteins can bend DNA and affect the binding of a second nearby protein. The effective interactions in that case also depends upon the protein orientation \[Diamant & Andelman 1999, Rudnick & Bruinsma, 1999\]. Moreover, proteins and protein aggregates need not be rigid, as we have assumed here. Noncircular distortions of a lipid or protein domains can fluctuate in such a way as to yield domain-domain interactions of exactly the same form considered here.
Elastically mediated attractions can also manifest themselves in the aggregation of circular proteins. Once circular proteins overcome short-ranged repulsions and dimerize, barriers to further aggregation of these elliptical dimers are reduced by dimer-dimer attractions described by Eq. (11). If the inclusions are themselves dimers or higher aggregates that persist on the time scale of rotation, bending mediated attraction would enhance further aggregation.
We have only considered the mechanical energies of the intervening lipid bilayer. However, an ensemble of membrane-bound proteins or a mixture of lipids can manifest itself through other forces. For example, the presence of charged membrane components can induce bending \[Chou, Jarić & Siggia, 1997\] and initiate endo/exocytosis and organelle trafficking. Protein-protein interactions arising from screened electrostatic forces operate at much shorter distances than those of the bending induced forces. Therefore, by spatially organizing charged membrane components, elastic interactions may also play an indirect role in large scale electrostatically induced membrane deformation.
We also considered an ensemble of surface proteins elastically coupled by membrane deformation and computed the deviation of its equation of state from that of an ideal solute. Although membrane-mediated protein-protein interactions are nonpairwise additive \[Kim, et al., 1998\], only the two-particle interaction is relevant for sparsely distributed proteins. On a flat membrane the second virial coefficient $`B_2<0`$ when $`\overline{\epsilon }0.95\sqrt{k_BT/b}`$ ($`\mathrm{\Delta }1.69`$). At this ellipticity, the elastically induced $`1/r^4`$ repulsive interactions just compensate for the rotationally averaged attractions. This dependence on $`b/k_BT`$ suggests that the cell can regulate protein-protein interactions by varying the lipid composition and hence the bending rigidity of the bilayer. Thus, the formation of cholesterol rafts and lipid domains may have an indirect role in mediating long-ranged surface protein aggregation and activity.
Finally, we propose possible experiments in artificial membrane systems where the surface density can be made small enough for a virial expansion to be valid. Although the two-dimensional osmotic pressure would be difficult to measure accurately, measurements of the association time between dimerized proteins are feasible. Measurements have been made of the lifetimes of gramicidin A channels composed of dimers of barrels in opposite bilayer leaflets as a function of bilayer thickness \[Kolb & Bamberg 1977, Elliot et al. 1983\]. Measurements of dimer lifetimes as a function of lipid tail length may reveal the dependence of the attractive interactions outlined in this paper. In fact, since the second virial coefficient measures the time-averaged fraction of proteins in dimers at low density, their lifetimes are proportional to $`B_2`$ for attracting proteins. Moreover, these association lifetimes can be measured in the presence of externally imposed Gaussian deformations. Even though an imposed Gaussian curvature increases the interaction well depth at $`\mathrm{\Omega }\pi /2`$, and destroys the attractions for proteins near $`\mathrm{\Omega }0`$, the overall statistical effect, is to enhance binding, as is evident from Figures 7. Therefore, we expect that dimer lifetimes can be enhanced for proteins residing in regions of large magnitudes of Gaussian curvature such as the base of extruded tubules. This may be instrumental in recruiting fusagens to the correct location for membrane budding. Dimer lifetimes potentially can be measured by fluorescence transfer of specifically designed hydrophilic moieties attached to membrane proteins. Nonpairwise interactions can only be probed directly by measuring lifetimes and aggregation rates of trimers. This would require statistical analyses of chemical or fluorescence activity among two differently tagged membrane proteins, or single molecule diffusion studies.
We thank J. B. Keller and J. C. Neu for helpful comments and many enlightening discussions. T. C. acknowledges support from the National Science Foundation via grant DMS-9804370. K. K. is supported by a grant from the Wellcome Trust, and G. O. is supported by NSF grant DMS-9220719.
## Appendix A Interaction energy among noncircular inclusions
We consider the boundary conditions that the height, $`h(r,\theta )`$, must satisfy and the effects of noncircular proteins on the interaction energies \[Kim, et al., 1999\]. Consider proteins with chemistry that changes the cross-sectional protein shape from circularity by an amount $`\epsilon `$. The concomitant changes in lipid contact height and angle are also assumed to be modified by $`O(\epsilon )`$. As shown in Fig. 1, the protein perimeter, measured from the protein center is, to order $`O(\epsilon )`$,
$$𝐂(a+\epsilon \mathrm{cos}2(\theta \theta _i))𝐧,$$
(22)
where $`𝐧`$ is the unit normal vector to the curve $`𝐂`$ projected onto the bilayer midplane, and $`\epsilon \mathrm{cos}2(\theta \theta _i)`$ is a small, angle-dependent perturbation measuring the deviation from circularity of protein $`i`$. Upon expanding the general boundary conditions $`h(𝐂)=\delta h(\theta )`$ and $`𝐧h(𝐂)=\gamma \delta \gamma (\theta )`$ to lowest order in $`\epsilon `$, we arrive at effective boundary conditions:
$$\begin{array}{c}h(a)\delta h(\theta \theta _i)+\gamma \epsilon \mathrm{cos}2(\theta \theta _i)+O(\epsilon ^2)\hfill \\ _rh(a)\gamma \left(1+\frac{\epsilon }{a}\mathrm{cos}2(\theta \theta _i)\right)\delta \gamma (\theta \theta _i)+O(\epsilon ^2)\hfill \end{array}$$
(23)
where we have for simplicity also assumed the variations $`\delta h(𝐂)`$ and $`\delta \gamma (𝐂)`$ to be also of order $`\epsilon `$.
In the limit of small noncircularity or low protein concentrations, the dominant nondivergent contribution of $`H(𝐫)`$ to the energy $`\stackrel{~}{E}`$ is $`a_2^2+b_2^2`$. The deformation $`h(r,\theta )`$ that satisfies $`^2h(r,\theta )=2H(r,\theta )`$ and Eqs. 23 can be written in the form
$$h(r,\theta )\gamma \mathrm{ln}\left(\frac{r}{a}\right)+\underset{n=2}{\overset{\mathrm{}}{}}\left(f_n(r)\mathrm{cos}n\theta +g_n(r)\mathrm{sin}n\theta \right),$$
(24)
and determine $`a_2,b_2`$. When the proteins have intrinsic noncircularity ($`\epsilon 0`$), $`a_2^2+b_2^2`$ turns out to be the magnitude of the local Gaussian curvature (since $`H_b=0`$), modified by additional $`\theta _i`$dependent terms \[Kim, et al., 1999\]. The local Gaussian curvature due to the other $`j`$ far field proteins, in either case, is calculated using the leading order term $`h(\stackrel{}{r})\gamma \mathrm{ln}|\stackrel{}{r}\stackrel{}{r}_j|`$, which is simply a superposition of the longest-ranged $`\mathrm{ln}r`$ terms about each inclusion. The total bending energy $`\stackrel{~}{E}[H(r,\theta )]`$ for an ensemble of $`N`$ inclusions can be written in the complex form \[Kim, et al., 1999\],
$$\stackrel{~}{E}=\pi b\gamma ^2\underset{j}{}\left|\underset{ij}{}\frac{a^2}{(z_iz_j)^2}\frac{\overline{\epsilon }}{2}e^{2i\theta _j}\right|^2.$$
(25)
where $`z_i=x_i+iy_i`$ is the position of the $`i^{th}`$ protein in the complex plane, and
$$\overline{\epsilon }\left(\frac{\epsilon }{a}\right)\left(\gamma +2\frac{\delta h}{a}\delta \gamma \right)$$
(26)
measures the effective ellipticity of the identical proteins. Now consider two relatively isolated, identical proteins $`i,j=1,2`$. The effects of proteins far away are felt via a local Gaussian curvature emanating from these background proteins. Upon explicitly separating these contributions, the pair interaction energy becomes
$$\stackrel{~}{E}(r,\theta _1,\theta _2;\eta _b,\mathrm{\Omega })=\pi b\left[\left|\frac{a^2\gamma e^{2i\mathrm{\Omega }}}{r^2}+\eta _b\frac{\overline{\epsilon }}{2}e^{2i\theta _1}\right|^2+\left|\frac{a^2\gamma e^{2i\mathrm{\Omega }}}{r^2}+\eta _b\frac{\overline{\epsilon }}{2}e^{2i\theta _2}\right|^2\right]$$
(27)
where
$$\eta _ba\frac{^2h_b(𝐒)}{x_1^2}=a\frac{^2h_b(𝐒)}{x_2^2}$$
(28)
is the curvature in the $`𝐱_1`$ principle direction due to far-field background inclusions or externally induced deformations $`h_b\gamma \mathrm{ln}|zz_j|,j3`$. The mean curvature expanded about a noncircular protein (Eq. 3) results in a deformation $`h(r,\theta )`$ with terms proportional to $`r^2\mathrm{cos}2\theta ,r^2\mathrm{sin}2\theta `$ \[Kim, et al., 1998\]. These term carry zero mean curvature, but constant negative Gaussian curvature. From the expansion Eq. 3, the only mean curvature contribution can be seen to decay as $`r^2`$, which we neglect. A further contribution to the local saddle curvature, $`\eta _b^2`$, felt by the two proteins can arise from externally applied mechanical forces that deform the bilayer in an appropriate way. The angles $`\theta _1,\theta _2`$ are the angles of the principle axes of the inclusion shape (or the height or contact angle functions $`\delta h,\delta \gamma `$) measured from the principle background curvature axis $`𝐱_1`$. The angle $`\mathrm{\Omega }`$ measures the angle between the principle background curvature axis and the segment joining the centers of the two inclusions. Upon rescaling according to Eq. 5, we arrive at the energy given in Equation 4.
## Appendix B Rotational averaging
The integrals
$$_0^{2\pi }E(R,\theta _1,\theta _2;K_b,\mathrm{\Omega })e^E𝑑\theta _1𝑑\theta _2\text{and}Z^{1/2}_0^{2\pi }e^E𝑑\theta _1𝑑\theta _2$$
(29)
used to compute the rotationally averaged, effective protein-protein interaction involve integration of
$$e^{\alpha \mathrm{cos}2\theta +\beta \mathrm{sin}2\theta }𝑑\theta \text{and}(\alpha \mathrm{cos}2\theta +\beta \mathrm{sin}2\theta )e^{\alpha \mathrm{cos}2\theta +\beta \mathrm{sin}2\theta }𝑑\theta .$$
(30)
The first integral in Eq. 30 can be computed in closed form by substituting the exponents with their Bessel function expansions
$$\begin{array}{c}e^{\alpha \mathrm{cos}2\theta }=I_0(\alpha )+2\underset{n=1}{\overset{\mathrm{}}{}}i^nI_n(\alpha )\mathrm{cos}2n\theta \hfill \\ e^{\beta \mathrm{sin}2\theta }=I_0(\beta )+2\underset{n=1}{\overset{\mathrm{}}{}}(1)^nI_{2n}(\beta )\mathrm{cos}4n\theta 2\underset{n=1}{\overset{\mathrm{}}{}}i^{2n+1}I_{2n+1}(\beta )\mathrm{sin}2(2n+1)\theta \hfill \end{array}$$
(31)
and integrating term by term. The cross-terms of the product of the two equations in Eq. 31 involve single powers of $`\mathrm{cos}`$ and $`\mathrm{sin}`$ and vanish upon integration. We are left with
$$Z^{1/2}=2\pi I_0(\alpha )I_0(\beta )+4\pi \underset{n=1}{\overset{\mathrm{}}{}}(1)^nI_{2n}(\alpha )I_{2n}(\beta ).$$
(32)
An analytic continuation of the sum formula,
$$J_0(\sqrt{\alpha ^2+\beta ^22\alpha \beta \mathrm{cos}\phi })=J_0(\alpha )J_0(\beta )+2\underset{n=1}{\overset{\mathrm{}}{}}J_n(\alpha )J_n(\beta )\mathrm{cos}n\phi ,$$
(33)
at $`\phi =\pi /2`$ simplifies Eq. 32 to,
$$Z^{1/2}=2\pi I_0(\xi ),\xi \sqrt{\alpha ^2+\beta ^2}.$$
(34)
Finally, the second integral in Eq. 30 can be computed by taking derivatives of $`Z^{1/2}`$:
$$(\alpha \mathrm{cos}2\theta +\beta \mathrm{sin}2\theta )e^{\alpha \mathrm{cos}2\theta +\beta \mathrm{sin}2\theta }𝑑\theta =\left(\alpha \frac{}{\alpha }+\beta \frac{}{_\beta }\right)Z^{1/2}.$$
(35)
Using these results, we arrive at the rotationally averaged energy $`E_{\mathrm{e}\mathrm{f}\mathrm{f}}`$ given by Eq. 9. For large separation distances $`R`$, the effective interaction $`U_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)E_{\mathrm{e}\mathrm{f}\mathrm{f}}(R)E_{\mathrm{e}\mathrm{f}\mathrm{f}}(\mathrm{})`$ defined in Eq. 11 can be expanded as in Eq. 15 where the coefficients are given by
$$A_24K_b2\mathrm{\Delta }^2K_b\frac{}{\xi }\left(\frac{I_1(\xi )}{I_0(\xi )}\right)_{\mathrm{\Delta }K_b}2\mathrm{\Delta }\frac{I_1(\mathrm{\Delta }K_b)}{I_0(\mathrm{\Delta }K_b)}$$
(36)
and
$$\begin{array}{c}\hfill A_42\mathrm{\Delta }^2\frac{}{\xi }\left(\frac{I_1(\xi )}{I_0(\xi )}\right)_{\mathrm{\Delta }K_b}\frac{\mathrm{\Delta }}{K_b}\frac{I_1(\mathrm{\Delta }K_b)}{I_0(\mathrm{\Delta }K_b)}\mathrm{sin}^22\mathrm{\Omega }\mathrm{\Delta }^2\left[K_b\frac{^2}{\xi ^2}+\frac{}{\xi }\right]\left(\frac{I_1(\xi )}{I_0(\xi )}\right)_{\mathrm{\Delta }K_b}\mathrm{cos}^22\mathrm{\Omega }.\end{array}$$
(37) |
warning/0001/hep-ph0001208.html | ar5iv | text | # MZ-TH/99-56 December 99 CP Violation with Three Oscillating Neutrino Flavours
## 1 Measuring leptonic CP violation with neutrino beams from a muon collider
Muon storage rings, muon colliders, and their physics potential are being studied intensively at FNAL and at CERN . In particular, a muon storage ring at some 20 GeV, being the first step in these projects, would serve to produce intense neutrino beams of unique quality. This possibility has received much attention recently. The straight sections of a muon collider would serve as sources of $`\overline{\nu }_\mu `$ and of $`\nu _e`$ with energy spectra perfectly calculable from muon decay, when positive muons are stored, and, similarly, of $`\nu _\mu `$ and $`\overline{\nu }_e`$ with well-known energy spectra when negative muons are stored.
In this work we explore the possibilities of studying CP violation in the leptonic sector, at a neutrino factory of this kind, by comparing the oscillation probabilities of CP-conjugate channels $`\nu _i\nu _j`$ and $`\overline{\nu _i}\overline{\nu }_j`$ with $`(ij)`$. The most suitable channels for studying CP violation are $`\nu _e\nu _\mu `$ and $`\overline{\nu }_e\overline{\nu }_\mu `$, as well as their T conjugate partners $`\nu _\mu \nu _e`$ and $`\overline{\nu }_\mu \overline{\nu }_e`$. In these channels, unlike the $`\nu _\mu \nu _\tau `$ channels, the CP violating part of the oscillation probability is not hidden by the CP conserving part , so that large asymmetries between CP-conjugate channels may arise, provided the leptonic CKM matrix allows for large violation of CP. Among these, the channels $`\nu _e\nu _\mu `$ and $`\overline{\nu }_e\overline{\nu }_\mu `$ seem to be the most promising because it may be easier to disentangle negative from positive muons, in a large detector of high density, than to disentangle electrons from positrons.
As a matter of example, we study the case of neutrino beams from stored muons with energy $`E_\mu =20`$ GeV, and an experimental arrangement where they travel over a distance of some 730 km, i.e. from CERN to the Gran Sasso laboratory, or, likewise, from FNAL to the Soudan mine. But we also give the scaling behaviour of the effects with either energy or distance. Within the range of squared mass differences and mixing angles, matter effects are important, see also . However, unlike the case where neutrinos traverse the Earth from the antipode , they are easy to cope with because the density in the Earth’s crust is essentially constant.
The $`\nu _e\nu _\mu `$ oscillation probability can be measured as follows: Suppose the electron neutrinos are produced by the decay of a number $`N_{\mu ^+}`$ of positive muons in the straight section of the storage ring pointing to the detector. The $`\nu _\mu `$ which appear when there is oscillation of $`\nu _e`$ into $`\nu _\mu `$, are detected by their charged current interaction in the detector. The number of observed muon neutrinos, $`n_{\nu _\mu }`$ is given by,
$`n_{\nu _\mu }=N_{kT}\mathrm{\hspace{0.33em}10}^9N_A{\displaystyle F_{\nu _e}\sigma _{\nu _\mu }P(\nu _e\nu _\mu )𝑑E}`$ (1)
where $`F_{\nu _e}`$ is the forward flux of electron neutrinos from a number $`N_{\mu ^+}`$ of positive muon decays, $`\sigma _{\nu _\mu }`$ is the charged current cross section per nucleon and $`P(\nu _e\nu _\mu )`$ is the oscillation probability for neutrinos traveling inside the Earth taking into account matter effects. $`N_{kT}`$ is the size of the detector in kilotons. An analogous formula holds in the case of anti-neutrinos.
Adopting the sample design configuration for muon production, cooling, acceleration and storage described by Geer , the number of available muons of either sign is approximately $`810^{20}`$ per year, for muons stored at an energy of 20 GeV. Of these, one fourth decay in a straight section directed towards the neutrino detector with a 10 kT target some $`730`$ km downstream, yielding about $`210^{20}`$ neutrinos per year and an identical number of anti-neutrinos. We use these numbers in what follows, and refer the reader to for details of the design of neutrino beams.
Let us begin by computing the number of produced muon neutrinos. Experimental cuts needed to eliminate background as well as detecting efficiencies will be included later on. The neutrino fluxes at a neutrino factory have simple analytical forms that follow from the well-known formulae for muon decay. Let $`x=E_\nu /E_\mu `$ be the fractional neutrino energy. For unpolarized muons of either charge, and neglecting corrections of order $`m_\mu ^2/E_\mu ^2`$ the normalized fluxes of forward moving electron neutrinos are
$`g_{\nu _e,\overline{\nu }_e}(x)=12x^2(1x)`$ (2)
and, for each neutrino type, the flux in the forward direction due to $`N_\mu `$ decaying muons is
$`F={\displaystyle \frac{d^2N_\nu }{dxd\mathrm{\Omega }}}|_{\theta 0}={\displaystyle \frac{E_\mu ^2N_\mu }{\pi m_\mu ^2L^2}}g_\nu (x)`$ (3)
The above expressions are valid for a detector placed in the forward direction whose transverse dimensions are much smaller than the beam’s transverse size $`(Lm_\mu /2E_\mu )`$.
We assume that the interaction cross sections due to charged current interactions scale linearly with the energy even in the very low energy part of the spectrum
$`\sigma _{\nu _e}=.6710^{38}\text{cm}\text{2}\text{ }E\text{ (GeV)}`$ (4)
$`\sigma _{\overline{\nu }_e}=.3410^{38}\text{cm}\text{2}\text{ }E\text{ (GeV)}`$ (5)
Regarding matter effects, let us remind the reader of the fact that of all neutrino species only $`\nu _e`$ and $`\overline{\nu }_e`$ have elastic scattering amplitudes on electrons due to charged current interaction. This, as is well known, induces effective “masses” $`\mu =\pm 2E_\nu a`$, where the upper sign refers to the electron neutrino, the lower sign to the corresponding anti-neutrino, and where $`a=\sqrt{2}G_Fn_e`$, $`n_e`$ being the electron density.
Matter effects are important provided the interaction term $`\mu `$,
$`\mu =7.710^5\text{eV}^2\left({\displaystyle \frac{\rho }{\text{gr/cm}\text{3}}}\right)\left({\displaystyle \frac{E_\nu }{\text{GeV}}}\right)`$ (6)
is comparable to, or bigger than, the quantity $`\mathrm{\Delta }_{m_{ij}^2}=m_i^2m_j^2`$ for some mass difference and neutrino energy.
CP related observables often involve the comparison between measurements in two charge-conjugate modes of the factory. One example of an asymmetry is
$`a_{CP}^{tot}={\displaystyle \frac{P(\nu _e\nu _\mu )F_{\nu _e}\sigma _{\nu _e}𝑑EP(\overline{\nu }_e\overline{\nu }_\mu )F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}𝑑E}{P(\nu _e\nu _\mu )F_{\nu _e}\sigma _{\nu _e}𝑑E+P(\overline{\nu }_e\overline{\nu }_\mu )F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}𝑑E}}`$ (7)
or in other terms,
$`a_{CP}^{tot}={\displaystyle \frac{n_{\nu _\mu }/N_{\mu ^+}n_{\overline{\nu }_\mu }/N_\mu ^{}}{n_{\nu _\mu }/N_{\mu ^+}+n_{\overline{\nu }_\mu }/N_\mu ^{}}}`$ (8)
In vacuum this quantity $`a_{CP}^{tot}`$ would be a pure CP odd observable. The voyage through our CP uneven planet, however, induces a nonzero asymmetry even if CP is conserved, since $`\nu _e`$ and $`\overline{\nu }_e`$ are affected differently by the electrons in the Earth . Therefore, to obtain the genuine CP odd quantity of interest, the matter effects must be subtracted with sufficient precision.
For this purpose, we compute the matter asymmetry in the absence of CP violation, or fake CP asymmetry, by
$`a_{CP}(\delta =0)={\displaystyle \frac{P(\nu _e\nu _\mu )_{\delta =0}F_{\nu _e}\sigma _{\nu _e}dEP(\overline{\nu }_e\overline{\nu }_\mu )_{\delta =0}F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}dE}{P(\nu _e\nu _\mu )_{\delta =0}F_{\nu _e}\sigma _{\nu _e}dE+P(\overline{\nu }_e\overline{\nu }_\mu )_{\delta =0}F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}dE}}`$ (9)
where we take into account matter effects but set $`\delta =0`$ in the transition probabilities.
The total asymmetry $`a_{CP}^{tot}`$ that will be found in an experiment of the type described above, is a function of $`a_{CP}(\delta =0)`$, eq. (9), and of the asymmetry in vacuum (taking due account of CP violation)
$`a_{CP}^{vac}={\displaystyle \frac{P^{vac}(\nu _e\nu _\mu )F_{\nu _e}\sigma _{\nu _e}𝑑EP^{vac}(\overline{\nu }_e\overline{\nu }_\mu )F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}𝑑E}{P^{vac}(\nu _e\nu _\mu )F_{\nu _e}\sigma _{\nu _e}𝑑E+P^{vac}(\overline{\nu }_e\overline{\nu }_\mu )F_{\overline{\nu }_e}\sigma _{\overline{\nu }_e}𝑑E}}`$ (10)
where $`P^{vac}(\nu _e\nu _\mu )`$ and $`P^{vac}(\overline{\nu }_e\overline{\nu }_\mu )`$ are the oscillation probabilities in vacuum. Provided $`a_{CP}(\delta =0)`$ is not too large.
$`a_{CP}^{vac}a_{CP}^{tot}a_{CP}(\delta =0)`$ (11)
is a good approximation. In any case, the error one makes in calculating $`a_{CP}`$ by means of eq.(11) is smaller than the uncertainties on $`a_{CP}(\delta =0)`$ itself. In addition, the error can be estimated by calculating the T-odd asymmetry , for each neutrino energy,
$`a_T(E_\nu ,\delta )={\displaystyle \frac{P(\nu _e\nu _\mu )P(\nu _\mu \nu _e)}{P(\nu _e\nu _\mu )+P(\nu _\mu \nu _e)}}`$ (12)
where a nonzero value cannot be induced by matter effects. This also means that $`a_T`$ a cleaner quantity in testing T violation than is $`a_{CP}`$ for CP violation.
An important component of any study of muon appearance due to $`\nu _e\nu _\mu `$ oscillations is the event selection strategy for the $`\mu `$’s produced from charged current interactions of the $`\nu _\mu `$’s. For neutrino experiments using a muon storage ring, the detailed prescription for event selection can be formulated only after the detector design is specified. There are, however, some basic issues concerning the signal and the backgrounds which all experiments are likely to be concerned with. On general grounds the background to a wrong sign muon signal is associated with the numerous decay processes that can produce fake muons: pions maskerading as muons, muonic charged currents (here one would also have to miss the right sign muon) or electronic charged currents (here one has also in the decay of the latter right sign muons), to name only a few. Without referring to a specific detector and the corresponding simulation toil we trust the experimental proficiency by setting an overall detection efficiency of 30% and by making a cut $`E_\nu >`$ 5 GeV to eliminate inefficiently observed low energy interactions.
## 2 Who mixes, two, three, or four flavours?
With growing evidence for non vanishing neutrino masses, experimental studies of neutrino oscillations, and their analysis in terms of three (or more) flavours, have become popular and will continue to be of central significance for lepton physics in the future.
The easiest way to describe any individual case of oscillations is to use a scheme where only two neutrino flavours are allowed to mix. Indeed, much work was done on analyses of neutrino oscillations in terms of two flavours but, as was pointed out by us and by others, the results for the squared mass differences may be misleading when applied to the real lepton world which contains three flavours. We summarize the situation regarding the squared mass differences as follows.
In models involving three oscillating flavours one often relies on squared mass differences which are taken from analyses of individual oscillation experiments in the framework of a two-flavour scheme. For instance, if one assumes the solar, the atmospheric and the LSND oscillations to be governed by just a single oscillation “frequency”, $`\mathrm{\Delta }^2`$, then the characteristic frequencies of the three oscillations, i.e.
$`\mathrm{\Delta }_{\text{solar}}^2=10^{10}\text{eV}^2\text{or}10^5\text{eV}^2`$
$`10^3\text{eV}^2\mathrm{\Delta }_{\text{atmospheric}}^210^2\text{eV}^2`$
$`10^1\text{eV}^2\mathrm{\Delta }_{\text{LSND}}^210^1\text{eV}^2`$
cannot be reconciled with just three neutrino mass eigenstates. Therefore, in order to simultaneously accommodate all three oscillations as observed, under the assumption stated above, we must introduce (at least) a fourth neutrino. Since we know from the width of the $`Z^0`$ boson that only three neutrino species have normal weak interactions, this extra, fourth neutrino must be sterile.
As there is no other, direct evidence for the existence of one or more sterile neutrinos, one is lead to conclude that assuming all observed oscillation phenomena to involve but a single $`\mathrm{\Delta }^2`$ is erroneous. Suppose, instead, that there are only three neutrinos, with masses such that
$`m_3^2m_2^2\mathrm{\Delta }M^2m_2^2m_1^2\mathrm{\Delta }m^2`$ (13)
Then, as was shown in , it is possible to explain the LSND result as an oscillation involving $`\mathrm{\Delta }M^2`$, the flavour conversion of solar neutrinos as one involving $`\mathrm{\Delta }m^2`$ and the atmospheric neutrino anomaly as a mixture of both frequencies. In contrast to these findings, an analysis of the atmospheric data assuming (erroneously) that only one $`\mathrm{\Delta }^2`$ is involved would find a value intermediate between those corresponding to the LSND and solar effects, as observed.
In calculating observable effects of CP violation in neutrino oscillations we assume a scenario with three flavours, where the two squared mass differences obey the inequality (13) and lie in the range
$$10^4\text{ eV}^2\mathrm{\Delta }m^210^3\text{ eV}^2,\mathrm{\Delta }M^20.3\text{ eV}^2.$$
(14)
If there are three Dirac neutrino types, then the flavour eigenstates are related to the mass eigenstates by a $`3\times 3`$ unitary matrix
$`U=\left(\begin{array}{ccc}c_{12}c_{13}& s_{12}c_{13}& s_{13}e^{i\delta }\\ s_{12}c_{23}c_{12}s_{23}s_{13}e^{i\delta }& c_{12}c_{23}s_{12}s_{23}s_{13}e^{i\delta }& s_{23}c_{13}\\ s_{12}s_{23}c_{12}c_{23}s_{13}e^{i\delta }& c_{12}s_{23}e^{i\delta }s_{12}c_{23}s_{13}& c_{23}c_{13}\end{array}\right)`$ (15)
where $`c_{ij}=\mathrm{cos}\theta _{ij}`$ and $`s_{ij}=\mathrm{sin}\theta _{ij}`$. If the neutrinos are Majorana particles, there are two extra phases, but these do not affect oscillations .
The transition probability in vacuum for a neutrino changing from $`\nu _i`$ ($`\overline{\nu }_i`$) to $`\nu _j`$ ($`\overline{\nu }_j`$) is given by the sum and the difference of CP-even and CP-odd pieces, respectively,
$`P(\nu _i\nu _j)`$ $`=`$ $`\text{}(\nu _i\nu _j)+\text{}(\nu _i\nu _j)`$ (16)
$`P(\overline{\nu }_i\overline{\nu }_j)`$ $`=`$ $`\text{}(\nu _i\nu _j)\text{}(\nu _i\nu _j),`$ (17)
where
$`\text{}(\nu _i\nu _j)`$ $`=`$ $`\delta _{ij}4\text{Re}J_{12}^{ji}\mathrm{sin}^2\mathrm{\Delta }_{12}4\text{Re}J_{23}^{ji}\mathrm{sin}^2\mathrm{\Delta }_{23}4\text{Re}J_{31}^{ji}\mathrm{sin}^2\mathrm{\Delta }_{31},`$
$`\text{}(\nu _i\nu _j)`$ $`=`$ $`8\sigma _{ij}J\mathrm{sin}\mathrm{\Delta }_{12}\mathrm{sin}\mathrm{\Delta }_{23}\mathrm{sin}\mathrm{\Delta }_{31},`$
with $`J`$ the Jarlskog invariant and
$`J_{kh}^{ij}`$ $``$ $`U_{ik}U_{kj}^{}U_{jh}U_{hi}^{}`$
$`\mathrm{\Delta }_{ij}`$ $``$ $`\mathrm{\Delta }m_{ij}^2L/4E`$ (19)
$`\sigma _{ij}`$ $``$ $`{\displaystyle \underset{k}{}}\epsilon _{ijk}`$
## 3 Results for CP asymmetries with three flavours
As stated above we henceforth assume a scenario with three flavours of neutrinos characterized by the squared mass differences (14) and the strong mixing found in , solution I, where
$`\theta _{12}35.5^0,\theta _{23}27.3^0,\theta _{13}13.1^0.`$ (20)
Although a detailed comparison might need further, refined analysis, this range of squared mass differences and the set of mixing angles (20) describes all observed neutrino anomalies in an overall and satisfactory manner. Here we show that this same set of parameters predicts a CP asymmetry which may well be large enough to be detectable with neutrino beams from a muon storage ring as described in sect. 1.
We organize the discussion of our results as follows: We first present our main result for the expected asymmetry. Next, the role of matter effects is illustrated by some examples, followed by a comparison of CP violating asymmetries with time reversal violating asymmetries. We then turn to a comparison with previous results and show why we find asymmetries which are sizeably larger than the ones estimated previously. Finally, some remarks on efficiencies and statistical uncertainties are added.
Let $`n`$ be the number of muons, $`\overline{n}`$ the number of antimuons detected in one year’s time in a detector placed at 732 km from the collider. Assuming for a moment an efficiency of $`100\%`$, we find the asymmetry $`(n\overline{n})/(n+\overline{n})`$ shown in Fig. 1, as a function of the neutrino energy $`E_\nu `$. Part (a) of the figure refers to the lower limit $`\mathrm{\Delta }m^2=10^4`$ eV<sup>2</sup>, part (b) refers to the upper limit $`\mathrm{\Delta }m^2=10^3`$ eV<sup>2</sup> of (14). The solid line corresponds to setting $`\delta =0`$ (no CP violation), the dashed line shows the full asymmetry, assuming $`\delta =\pi /2`$. As a matter of example Fig. 2 shows the absolute numbers $`n`$ and $`\overline{n}`$, obtained in one year of running, under the same assumptions as before and for $`\mathrm{\Delta }m^2=10^3`$ eV<sup>2</sup>.
*The role of matter effects:* Clearly, the smaller matter effects the cleaner the measurement of the effects of genuine CP violation will be. As the asymmetry due to charged-current matter interaction grows faster than the CP asymmetry, as a function of the baseline, intermediate distances between collider and neutrino detector are preferred over long distances. We illustrate this observation quantitatively by defining the asymmetry
$$a_{CP}(E_\nu ,\delta )=\frac{P(\nu _e\nu _\mu )P(\overline{\nu }_e\overline{\nu }_\mu )}{P(\nu _e\nu _\mu )+P(\overline{\nu }_e\overline{\nu }_\mu )}$$
(21)
and by calculating the ratio $`a_{CP}(E_\nu ,0)/[a_{CP}(E_\nu ,\pi /2)a_{CP}(E_\nu ,0)]`$, as a function of $`E_\nu `$. Fig. 3 shows this quantity for a baseline of 732 km, part (a), and a baseline of 7332 km, part (b), corresponding to the distance from FNAL to the Gran Sasso. In the case of the very long baseline, the CP asymmetry is completely swamped by matter effects.
*CP- versus T-asymmetry:* We also computed the T-odd asymmetry (12) for the example $`\mathrm{\Delta }m^2=10^4`$ eV<sup>2</sup> and the shorter baseline $`L=732`$ km and compared it to the CP asymmetry (21) from which the matter effects were subtracted. We found them to agree within reasonable limits, thus corroborating the approximation (11).
*Comparison with previous results:* The authors of ref. who use the following sample set of parameters
$$\mathrm{\Delta }m_{12}^2=10^4\text{ eV}^2,\mathrm{\Delta }m_{23}^2=10^3\text{ eV}^2,\mathrm{sin}^2\theta _{12}=.5,\theta _{23}=45^0,\theta _{13}=13^0,$$
(22)
find CP violating effects which are markedly smaller than the ones we showed above. The explanation for this difference is simple: The most noticeable difference between the set (22) and ours is the value of $`\mathrm{\Delta }m_{23}^2`$. Assuming the values (22) both the (12)- and the (13)-channels are strongly affected by matter effects, the effective sines $`(\mathrm{sin}\theta _{12})_{matter}`$ and $`(\mathrm{sin}\theta _{13})_{matter}`$ as defined in quickly tend to 1 as the parameter $`\mu `$, eq. (6), increases. Consequently, the simultaneous interplay of all three flavours and, hence, the visibility of CP violation decrease. In contrast to this situation our values of squared mass differences are such that only $`(\mathrm{sin}\theta _{12})_{matter}`$ is affected while $`(\mathrm{sin}\theta _{13})_{matter}`$ and, of course, $`(\mathrm{sin}\theta _{23})_{matter}`$ remain unaffected. It is convenient to define an effective mixing matrix $`V_{ik}`$ which is obtained from (15) by replacing the sines and cosines by the matter affected sines and cosines, $`(\mathrm{sin}\theta _{ik})_{matter}`$, etc.
To illustrate the comparison Fig. 4(a) shows the pertinent, effective matrix elements for our set of parameters, eqs. (14) and (20), as a function of $`\mu `$, while Fig. 4(b) shows the same matrix elements for the set (22). In the latter case, both $`V_{11}`$ and $`V_{12}`$ tend to zero with $`\mu `$ increasing to its value in the Earth’s crust.
*Efficiencies and statistical uncertainties:* The discussion of statistical uncertainties is straightforward. In order to exclude the possibility that a measurement of a non vanishing genuine CP violation is due to a statistical fluctuation, the measured value must be larger than $`n\delta (a_{CP})_{\text{stat}}`$, where $`\delta (a_{CP})_{\text{stat}}`$ is the $`1\sigma `$ statistical error on $`a_{CP}`$ in the absence of CP violation, and $`n`$ is the number of standard deviations we require in order to be happy with our result. Since in absence of CP violation the expectations of $`n_{\nu _\mu }/N_{\mu ^+}`$ and $`n_{\overline{\nu }}/N_\mu ^{}`$ are equal, we get
$`\delta (a_{CP})_{\text{stat}}=\left({\displaystyle \frac{1}{4n_{\nu _\mu }}}+{\displaystyle \frac{1}{4n_{\overline{\nu }_\mu }}}\right)^{1/2}`$ (23)
where $`n_{\nu _\mu }`$ ($`n_{\overline{\nu }_\mu }`$) is the expected number of $`\nu _\mu `$ ($`\overline{\nu }_\mu `$) interactions seen in the detector.
Regarding the contribution of the background to the statistical error, and according to the current estimates, the main source of background will be due to charm production in the charged current neutrino interactions in the detector
$`\mu ^{}\nu _\mu `$ CC interaction $`\mu ^{}_{\text{ }\text{lost}}`$
$`c`$ c decay $`\mu ^+_{\text{ }\text{found}}`$
Clearly, this background only affects the signal corresponding to the $`\overline{\nu _e}\overline{\nu }_\mu `$ oscillations (because we expect there $`\mu ^+`$ appearance), and therefore, this background “noise” should be subtracted appropriately in the counting of $`n_{\overline{\nu }_\mu }`$ in Eq.(8). Such a subtraction introduces a further source of statistical error. Using the estimate that
$`n_{\overline{\nu }_\mu }_{\text{back}}10^5n_{\overline{\nu }_\mu }_{P=1}`$ (24)
where $`n_{\overline{\nu }_\mu }_{P=1}`$ is the number of antimuon neutrino interactions that would be seen if all the initial antielectron neutrinos oscillated into antimuon neutrinos, we find
$`\delta (a_{CP})_{\text{stat}}=\left({\displaystyle \frac{1}{4n_{\nu _\mu }}}+{\displaystyle \frac{1}{4n_{\overline{\nu }_\mu }}}+{\displaystyle \frac{10^5n_{\overline{\nu }_\mu }_{P=1}}{4n_{\overline{\nu }_\mu }^2}}\right)^{1/2}`$ (25)
as the complete expression for the statistical error. It is important to notice that even with only one year of data taking and a modest 30% detecting efficiency, the statistical error will be small enough to rule out the possibility of attributing to a statistical fluctuation a measurement of a non vanishing CP violation.
## 4 Conclusions
In summary, the two rather different “frequencies” (13), together with the strong mixing of all three flavours that describe the solar neutrino deficit, the atmospheric oscillations, and the LSND anomaly, lead to relatively large CP and T violating asymmetries in neutrino oscillations. With the set of parameters (14) and (20) the full interference of all three flavours is well developed and is only moderately damped by matter effects. We have also tried the other solutions to the neutrino anomalies we had found in but find no more than 20% changes in the asymmetries. Among these the CP asymmetry seems large enough to be measurable with neutrino beams from a 20 GeV muon storage ring and with a detector at some 730 km from the source, corresponding to the distance of the Gran Sasso laboratory from CERN or, likewise, of the Soudan mine from FNAL.
Acknowledgements
We are very grateful to Karl Jakobs for enlightening discussions. Financial support from the DFG is also acknowleged. |
warning/0001/hep-ph0001320.html | ar5iv | text | # Oblique Parameter Constraints on Large Extra Dimensions
## I Introduction
The possible existence of extra dimensions has been a fascinating idea in physics ever since Kaluza-Klein theory was proposed . Consistent string theories demand the existence of extra dimensions . However, if the string scale ($`M_S`$) is as high as the grand unification scale ($`M_{\mathrm{GUT}}10^{16}`$ GeV) or the Planck scale ($`M_{\mathrm{Pl}}10^{19}`$ GeV), as is the case for a weakly coupled heterotic string, then the length scale of the compactified extra dimensions $`(R1/M_S)`$ would be too small to be appreciable experimentally. Recent developments in string theory indicate that the string scale can be much lower than the Planck scale and even close to the electroweak scale . This possibility provides new avenues towards many theoretical issues such as alternative solutions to the gauge hierarchy problem , fermion mass and flavor mixings , and new inflationary cosmological models . More importantly, such a scenario may lead to a rich phenomenology and is thus experimentally testable at low energies .
Assume that there are $`n`$ extra dimensions in which only gravity can propagate while Standard Model (SM) fields are confined to four dimensional spacetime. The large value of the Planck scale can be understood by Gauss’ law from the relation
$$M_{\mathrm{Pl}}^2R^nM_S^{n+2},$$
(1)
where the string scale $`M_S`$ is taken to be the Planck mass in $`4+n`$ dimensions. For $`M_S𝒪`$ (1 TeV), $`R`$ can range from 1 fm to 1 mm for $`n=6`$ to 2 . There are no direct gravitational tests sensitive to those small scales yet . Such large extra dimensions manifest themselves only through interactions involving the Kaluza-Klein (KK) modes of the gravitons with enhanced coupling strength after summing over the many contributing light KK states. The effective theory governing the graviton couplings to matter was described in . Phenomenological studies showed that future collider experiments can provide constraints on $`M_S`$ typically of order 1 TeV, depending on the collider center of mass energies . Astrophysical (cosmological) considerations have been used to impose lower bounds on $`M_S`$ to be about 30 (100) TeV for $`n=2`$ and very weak bounds for higher dimensions .
In this paper we consider radiative corrections to the masses of the electroweak gauge bosons $`(W,Z)`$ arising from the exchange of massive spin-2 KK gravitons. The motivation for this study is two-fold. Firstly, a rigorously renormalizable theory of gravity does not exist. It would be interesting to explore to what extent the formalism in Ref. is finite against radiative corrections in the sense of an effective field theory. An early attempt in Ref. for a scalar self-energy correction showed that the radiative corrections are proportional to the scalar mass, instead of the ultraviolet cutoff. Another one-loop calculation for the muon $`g2`$ also reached finite results . Secondly, if one is able to compute radiative effects from KK gravitons, one may hope to examine constraints on new physics characterized by the string scale from precision electroweak measurements. In fact, a recent paper appeared to estimate the $`\rho `$ parameter , which was found to be both ultraviolet and infrared divergent. However, our results arrive at completely different conclusions.
The rest of the paper is organized as follows. In Section II we compute one-loop self-energy diagrams for a gauge boson from exchange of massive spin-2 KK gravitons. To do so, we need to derive the four–point couplings of two gravitons and two gauge bosons which are beyond that given in Ref. and have been left out in the previous calculations . In Section III we adopt the $`STU`$-formalism and constrain the string scale based on current experimental values of the oblique parameters. Computing oblique parameters has an advantage over calculating $`W(Z)`$ mass corrections directly since they are free of certain technical and conceptual uncertainties. We discuss our results and conclude in Section IV. In two appendices we present the relevant Feynman rules and describe the regularization of the infrared divergences.
## II Self-energy corrections to the gauge bosons
When the typical energy scale of interest is much smaller than the string scale $`M_S`$, the compactified higher dimensional theory can be described by an effective theory where only relevant light degrees of freedom are retained, namely, those related to the so called Kaluza-Klein states. In principle, if only gravity is of higher dimensional origin, spin-0, 1 and 2 KK states can arise . The spin-1 states decouple from the Standard Model fields and make no contribution to the self-energies. Naively, one expects the spin-0 contributions to be of the same order of magnitude as that of the spin-2 states since they have the same coupling strength. However, the properties of the spin-0 states are model-dependent since there is no symmetry to protect their masses. In the following we shall concentrate on the spin-2 KK states (which we will call KK gravitons).
The physical process that we consider in this paper is the radiative corrections to the weak gauge boson masses at the one-loop level from these KK gravitons. We focus our attention on the transverse part of the self-energy $`\mathrm{\Pi }_{XY}(p^2)`$ between gauge bosons $`X`$ and $`Y`$, as it is the only part relevant to us. These self-energies are written as
$$\mathrm{\Pi }_Z(p^2)\mathrm{\Pi }_{ZZ}^T(p^2)\mathrm{and}\mathrm{\Pi }_W(p^2)\mathrm{\Pi }_{WW}^T(p^2).$$
(2)
Note that both $`\mathrm{\Pi }_{Z\gamma }`$ and $`\mathrm{\Pi }_{\gamma \gamma }`$ are identically zero, following from the simple fact that they must be proportional to the photon mass, as required by the gravitational nature of the KK graviton-matter interactions.
There are two diagrams contributing to the self-energy of a gauge boson, as shown in Fig. 1. The first (seagull diagram) involves a four–point coupling of two KK gravitons and two gauge bosons, and the second (rainbow diagram) involves two three–point couplings of a KK graviton and two gauge bosons. These two diagrams are at the same order in the gravitational coupling $`\kappa ^2=16\pi G_N`$, where $`G_N`$ is the usual four-dimensional Newton’s constant.
To obtain the four-point vertex Feynman rule required for the seagull diagram, one needs to expand the graviton-matter interaction Lagrangian to order $`\kappa ^2`$, which is beyond the order of the expansion given in . We provide this Feynman rule in Appendix A. The corresponding seagull diagram is purely transverse, proportional to $`(\eta _{\mu \nu }p_\mu p_\nu /p^2)\mathrm{\Pi }_S(p^2)`$, with
$$\mathrm{\Pi }_S(p^2)=\frac{\kappa ^2p^2}{16\pi ^2}\underset{\stackrel{}{n}}{}_0^{\mathrm{}}\frac{dk_E^2k_E^2}{k_E^2+m_\stackrel{}{n}^2}\left(\frac{k_E^4}{12m_\stackrel{}{n}^4}+\frac{3k_E^2}{4m_\stackrel{}{n}^2}+\frac{3}{2}\right),$$
(3)
where we have included a factor $`1/2`$ to avoid double-counting when we sum over the KK modes and the subscript $`E`$ indicates that the variable is in Euclidean space. The summation over the KK modes in a tower can be written as an integration because of the near-degeneracy of the KK states,
$$\underset{\stackrel{}{n}}{}f(m_\stackrel{}{n})=_0^{\mathrm{}}𝑑m_\stackrel{}{n}^2\rho (m_\stackrel{}{n})f(m_\stackrel{}{n})$$
(4)
where
$$\rho (m_\stackrel{}{n})=\frac{R^nm_\stackrel{}{n}^{n2}}{(4\pi )^{n/2}\mathrm{\Gamma }(n/2)}$$
(5)
is the KK state density. By convention for the torus compactification , the relation between the four-dimensional Newton’s constant and the $`(4+n)`$-dimensional string scale $`M_S`$ is given by
$$\kappa ^2R^n=8\pi (4\pi )^{n/2}\mathrm{\Gamma }(n/2)M_S^{(n+2)}.$$
(6)
Using Eqs. (4), (5) and (6), we obtain
$$\frac{\kappa ^2}{16\pi ^2}\underset{\stackrel{}{n}}{}f(m_\stackrel{}{n})=\frac{1}{2\pi M_S^{n+2}}_0^{\mathrm{}}𝑑m_\stackrel{}{n}^2m_\stackrel{}{n}^{n2}f(m_\stackrel{}{n}^2).$$
(7)
The above sum is divergent in the ultraviolet for $`n2`$ unless $`f(m_\stackrel{}{n})`$ falls off very rapidly with $`m_\stackrel{}{n}^2`$. Since the effective theory is only expected to be valid below the string scale, we introduce an explicit cutoff $`\lambda M_S`$ to regularize the mass sum, where $`\lambda 𝒪(1)`$ and parameterizes the sensitivity to the cutoff.<sup>*</sup><sup>*</sup>*This bad ultraviolet behavior can be remedied in models of fluctuating branes and models of fermions located in different branes . We shall use the same cutoff for the momentum integral as for the mass summation. We believe that this regulator will best reflect the ultraviolet behavior of the divergence. Our final result for the seagull diagram is
$$\mathrm{\Pi }_S(p^2)=\frac{\lambda ^{n+2}p^2}{12\pi }_0^1_0^1𝑑x𝑑y\frac{xy^{3+\frac{n}{2}}}{x+y}\left(x+3y\right)\left(x+6y\right),$$
(8)
where $`x=(k_E/\lambda M_S)^2`$ and $`y=(m_\stackrel{}{n}/\lambda M_S)^2`$.
The necessary Feynman rules for calculating the rainbow diagram have been derived in and are summarized in Appendix A. The complete expression for the general $`\mathrm{\Pi }_R(p^2)`$ is lengthy and unilluminating. We will instead present the special cases relevant to our current consideration, $`\mathrm{\Pi }_R(0)`$ and $`\mathrm{\Pi }_R(m^2)`$. The former can be written as
$`\mathrm{\Pi }_R(0)`$ $`=`$ $`{\displaystyle \frac{\lambda ^{n+2}m^2}{24\pi }}{\displaystyle _0^1}{\displaystyle _0^1}dxdy{\displaystyle \frac{xy^{3+\frac{n}{2}}}{\left(r+x\right)\left(x+y\right)}}\times `$ (10)
$`\left[x\left(x+y\right)\left(x+13y\right)+r\left(4x^2+26xy+52y^2\right)\right],`$
where we have introduced a dimensionless mass ratio between the gauge boson mass and the ultraviolet cutoff $`r=(m/\lambda M_S)^2`$. Similarly,
$`\mathrm{\Pi }_R(m^2)`$ $`=`$ $`{\displaystyle \frac{\lambda ^{n+2}m^2}{24\pi }}{\displaystyle _0^1}{\displaystyle _0^1}{\displaystyle _0^1}𝑑x𝑑y𝑑z{\displaystyle \frac{xy^{3+\frac{n}{2}}F(x,y,z,r)}{[x+y(1z)+rz^2]^2}},`$ (11)
where $`z`$ is a Feynman parameter and
$`F(x,y,z,r)`$ $`=`$ $`4r^3\left(2+z\right)^2z^4+r^2z^2\left[16y(2+z)+x(24+52z21z^2)\right]`$ (14)
$`+r\left[4xy(28z+z^2)+4y^2(4+9z^2)+x^2(114z+15z^2)\right]`$
$`x(x^2+4xy+23y^2).`$
It is important to note that the on-shell mass corrections are proportional to the gauge boson mass squared $`m^2`$, similar to the results obtained in Ref. for a scalar mass correction, and not to $`M_S^2`$ as found in . Consequently, there is no hard quadratic dependence on the cutoff. On the other hand, the proportionality factor $`\lambda ^{n+2}`$ appearing in the above expressions implies that the results are rather sensitive to the precise value of the cutoff in comparison to the string scale $`M_S`$. This is an intrinsic uncertainty for any process involving virtual KK graviton exchanges in an effective theory.
Before we end this section, a few remarks related to the potential uncertainties of the results are in order:
1. Eqs. (8), (10) and (11) appear to be infrared divergent as a consequence of the pole at $`y=0`$ for massless gravitons. This must be an artifact of the fixed-order perturbative calculation since gravity is known to be infrared safe. We will regularize the infrared singularity using the principles of dimensional regularization and perform the minimal subtraction. The procedure is described in Appendix B.
2. It is interesting to investigate the limit for a large value of the cutoff $`\lambda M_S\mathrm{}`$, or equivalently $`r0`$. As an illustration we consider $`n=2`$. The seagull diagram is independent of $`r`$ and gives
$$\mathrm{\Pi }_S(0)=0,\mathrm{\Pi }_S(m^2)=\frac{7\lambda ^4m^2}{18\pi },$$
(15)
and the rainbow diagram gives
$$\underset{r0}{lim}\mathrm{\Pi }_R(0)=\frac{\lambda ^4m^2}{72\pi },\underset{r0}{lim}\mathrm{\Pi }_R(m^2)=\frac{29\lambda ^4m^2}{72\pi }.$$
(16)
Therefore, the total self-energy correction $`\mathrm{\Pi }(p^2)=\mathrm{\Pi }_R(p^2)+\mathrm{\Pi }_S(p^2)`$ takes the following values at $`p^2=0`$ and $`p^2=m^2`$:
$$\underset{r0}{lim}\mathrm{\Pi }(0)=\underset{r0}{lim}\mathrm{\Pi }(m^2)=\frac{\lambda ^4m^2}{72\pi }.$$
(17)
One would hope that the self-energy corrections vanish as the string scale is set to infinity, as required by the decoupling theorem. This has been explicitly shown not to be the case by the naive results in Eq. (17). The problem lies in the fact that an unknown cosmological constant can also contribute to the gauge boson self-energies via gravitational interactions. A non-zero cosmological constant term is of the form
$$\mathrm{\Lambda }d^4x\sqrt{g},$$
(18)
and would lead to an additional contribution to the self-energy by tadpole diagrams,
$$\mathrm{\Pi }_\mathrm{\Lambda }(p^2)m^2\frac{\mathrm{\Lambda }\lambda ^{n2}}{M_S^4}_0^1y^{\frac{n}{2}2}𝑑y.$$
(19)
This term can drastically change the previously calculated self energies. If $`\mathrm{\Lambda }(\lambda M_S)^4`$, Eq. (19) could be at the same order of Eq. (17) and could thus provide the appropriate counter-term. However, due to the lack of a consistent way of determining the cosmological constant, the precise value of the term in Eq. (19) is not known. As a result, the self-energies and subsequently, corrections to the $`W`$ and $`Z`$ boson masses have inherent uncertainties.
3. In our calculations for the gauge boson self-energies, we have adopted the momentum-cutoff scheme to regularize the divergent mass sum and the momentum integral, since we consider this scheme a most direct reflection of the ultraviolet behavior. We anticipate that the physics results would not depend upon the specific regularization scheme.
It is important to emphasize that the above potential ambiguities are of no concern to us if we adopt the $`STU`$-formalism since the oblique parameters are manifestly finite . The regulator independence built into the definition of the oblique parameters make the infrared and ultraviolet divergences as well as the irrelevant constants drop out when taking the difference of the appropriate combination of self-energies, as we will present next.
## III Oblique parameters
The studies of electroweak radiative corrections have proven to be powerful in constraining new physics beyond the SM. A convenient parameterization of new physics from a higher scale is the $`STU`$-formalism . The $`S`$, $`T`$ and $`U`$ parameters can be obtained by evaluating the self-energy corrections at the energy scales $`m_Z`$ and 0. We write the oblique parameters asThese definitions are the same as in Ref. . In the case under consideration, they are identical to those originally introduced in Ref. .
$`\alpha S`$ $`=`$ $`4s^2c^2{\displaystyle \frac{\mathrm{\Pi }_Z(m_Z^2)\mathrm{\Pi }_Z(0)}{m_Z^2}},`$ (20)
$`\alpha T`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }_W(0)}{m_W^2}}{\displaystyle \frac{\mathrm{\Pi }_Z(0)}{m_Z^2}},`$ (21)
$`\alpha (S+U)`$ $`=`$ $`4s^2{\displaystyle \frac{\mathrm{\Pi }_W(m_W^2)\mathrm{\Pi }_W(0)}{m_W^2}},`$ (22)
where $`s`$ and $`c`$ are the sine and cosine of the weak mixing angle and $`\alpha `$ is the fine structure constant, all measured at $`m_Z`$.
The $`S`$ ($`S+U`$) parameter measures the difference in the contribution of new physics to neutral (charged) current processes at different energy scales. $`U`$ is generally small. The $`T`$ parameter serves as a comparison between the new contributions to the neutral and charged current processes at low energy, proportional to $`\mathrm{\Delta }\rho `$. Comparing experimental data mainly from the LEP and SLC to SM predictions with $`m_H=300`$ GeV leads to the bounds ,
$`S`$ $`=`$ $`0.30\pm 0.13,`$ (23)
$`T`$ $`=`$ $`0.14\pm 0.15,`$ (24)
$`U`$ $`=`$ $`\mathrm{\hspace{0.17em}0.15}\pm 0.21.`$ (25)
We perform an oblique correction analysis using the calculations of the previous section and the SM parameters
$`m_Z=91.1867\mathrm{GeV},m_W=80.315\mathrm{GeV}\mathrm{and}s^2=0.232.`$ (26)
We assume the central values of Eq. (24) to be the SM predictions and attribute the error bars to the physics contribution of our current interest. Our numerical results are presented in Figs. (2)-(4) where we have set $`\lambda =1`$.
In Fig. 2, we have plotted $`\mathrm{\Delta }S`$, the excess contribution to the SM value of the $`S`$ parameter arising from the KK graviton exchanges of the one-loop self-energy diagrams, for $`n=2`$ (positive) and $`n=3,4`$ (negative). These values must lie within the error bars stated in Eq. (24) to not be in conflict with precision electroweak measurements. We see that for $`n=2`$, this leads to a lower bound $`M_S>1.55`$ TeV and for $`n=3`$, $`M_S>600`$ GeV at the $`1\sigma `$ level. There is no constraint for $`n4`$ from the $`S`$ parameter and we have therefore neglected to plot the cases $`n=5,6`$.
Figure 3 shows that the $`T`$ parameter imposes constraints on $`M_S`$ for all $`n`$. This seems to be counter-intuitive at first sight since the interactions of KK gravitons with matter should respect the custodial $`SU(2)`$ symmetry and thus lead to a null contribution to $`T`$. However, the mass difference of the gauge bosons make their couplings to the gravitons different, resulting in a negative graviton contribution to $`T`$. In fact, the constraints for $`n3`$ are more stringent than that obtained from the $`S`$ parameter. For instance, the lower limit on $`M_S`$ is raised to 1.25 (0.75) TeV for $`n=3`$ (6).
The $`U`$ parameter is generally small and it is true here as well. On the other hand, the error bars on $`U`$ are still rather large and we would not expect improvement by it. Indeed, from Fig. 4, we see that the $`U`$ parameter places no constraint whatsoever.
## IV Discussion and Conclusion
We have obtained significant constraints on the string scale $`M_S`$, by considering the $`S`$, $`T`$ and $`U`$ parameters from precision electroweak data. One may wonder if the $`STU`$-formalism is suitable since gravitons can be lighter than gauge bosons. We consider our treatment appropriate because as a collective contribution, the relevant effects are characterized by the string scale at about a TeV. A more subtle question is whether the non-oblique corrections would also be as important and how they can be incorporated. In fact, leading KK graviton corrections appear even at tree level and significant effects were found in the literature . We thus view our approach to single out the oblique corrections as the next-to-leading contribution as reasonable.
As we see from Eq. (17), the self-energy corrections remain finite in the limit $`M_S\mathrm{}`$. As explained in the paragraph following it, it is possible to achieve “decoupling behavior” by introducing a cosmological constant as a counter-term. By no means is this necessary as far as the oblique parameters are concerned. The effect of the heavy states does decouple in $`S`$, $`T`$ and $`U`$ as seen from the relation Eq. (17),
$$\underset{M_S\mathrm{}}{lim}S,T,U=0.$$
(27)
Altough demonstrated specifically for the case $`n=2`$, we have confirmed this to be true for all $`n`$. It is reassuring to obtain the “decoupling” relations, which imply that the radiative corrections based on this effective field theory of KK gravitons are under control.
For all of our numerical analysis, the choice $`\lambda =1`$ corresponds to taking the cutoff scale as the string scale. This is the same as the choice in most phenomenological studies . There is the intrinsic uncertainty due to the ratio parameter $`\lambda `$, although it is reasonable to assume it to be of order unity. More definitive results will depend on details of a string model near the string scale.
In summary, we obtained significant lower bounds on $`M_S`$ from the $`S`$ and $`T`$ parameters. For one and two standard deviations, they are
$$\begin{array}{ccccc}n& & M_S(\mathrm{GeV})\mathrm{at}1\sigma & & M_S(\mathrm{GeV})\mathrm{at}2\sigma \\ & & & & \\ 2& & 1550& & 1100\\ 3& & 1250& & 850\\ 4& & 950& & 650\\ 5& & 900& & 600\\ 6& & 750& & 500\end{array}$$
(28)
These results are comparable to that inferred from current LEP II experiments and are slightly weaker than those anticipated at future runs of LEP and the Tevatron.
###### Acknowledgements.
We thank C. Goebel, J. Lykken and D. Zeppenfeld for discussions. This work was supported in part by a DOE grant No. DE-FG02-95ER40896 and in part by the Wisconsin Alumni Research Foundation.
## A
In this appendix we summarize the Feynman rules used in the self-energy calculations.
The propagator for the massive spin-2 KK states $`\stackrel{~}{h}_{\mu \nu }^\stackrel{}{n}`$ is
$$i\mathrm{\Delta }_{\{\mu \nu ,\stackrel{}{n}\},\{\rho \sigma ,\stackrel{}{m}\}}^{\stackrel{~}{h}}(k)=\frac{i\delta _{\stackrel{}{n},\stackrel{}{m}}B_{\mu \nu ,\rho \sigma }(k)}{k^2m_\stackrel{}{n}^2+i\epsilon },$$
(A1)
where
$`B_{\mu \nu ,\rho \sigma }(k)`$ $`=`$ $`\left(\eta _{\mu \rho }{\displaystyle \frac{k_\mu k_\rho }{m_\stackrel{}{n}^2}}\right)\left(\eta _{\nu \sigma }{\displaystyle \frac{k_\nu k_\sigma }{m_\stackrel{}{n}^2}}\right)+\left(\eta _{\mu \sigma }{\displaystyle \frac{k_\mu k_\sigma }{m_\stackrel{}{n}^2}}\right)\left(\eta _{\nu \rho }{\displaystyle \frac{k_\nu k_\rho }{m_\stackrel{}{n}^2}}\right)`$ (A3)
$`{\displaystyle \frac{2}{3}}\left(\eta _{\mu \nu }{\displaystyle \frac{k_\mu k_\nu }{m_\stackrel{}{n}^2}}\right)\left(\eta _{\rho \sigma }{\displaystyle \frac{k_\rho k_\sigma }{m_\stackrel{}{n}^2}}\right).`$
The three-point vertex is shown in Fig. 5 and the corresponding Feynman rule is
$$i\frac{\kappa }{2}\delta ^{ab}\left[(m^2+k_1\text{.}k_2)C_{\mu \nu ,\rho \sigma }+D_{\mu \nu ,\rho \sigma }(k_1,k_2)\right],$$
(A4)
where
$$C_{\mu \nu ,\rho \sigma }=\eta _{\mu \rho }\eta _{\nu \sigma }+\eta _{\mu \sigma }\eta _{\nu \rho }\eta _{\mu \nu }\eta _{\rho \sigma }$$
(A5)
is the tensor that appears in the massless graviton propagator in the de Donder gauge, and
$$D_{\mu \nu ,\rho \sigma }(k_1,k_2)=\eta _{\mu \nu }k_{1\sigma }k_{2\rho }[\eta _{\mu \sigma }k_{1\nu }k_{2\rho }+\eta _{\mu \rho }k_{1\sigma }k_{2\nu }\eta _{\rho \sigma }k_{1\mu }k_{2\nu }+(\mu \nu )].$$
(A6)
The Feynman rule for the four-point vertex as shown in Fig. 5 is
$$i\frac{\kappa ^2}{4}\delta ^{ab}\left[(m^2+k_1\text{.}k_2)C_{\mu \nu ,\rho \sigma \omega \tau }+H_{\mu \nu \rho \sigma \omega \tau }(k_1,k_2)+I_{\mu \nu \rho \sigma \omega \tau }(k_1,k_2)\right]$$
(A7)
where
$$C_{\mu \nu ,\rho \sigma \omega \tau }=\frac{1}{2}[\eta _{\mu \omega }C_{\rho \sigma ,\nu \tau }+\eta _{\sigma \omega }C_{\mu \nu ,\rho \tau }+\eta _{\rho \omega }C_{\mu \nu ,\sigma \tau }+\eta _{\nu \omega }C_{\mu \tau ,\rho \sigma }\eta _{\omega \tau }C_{\mu \nu ,\rho \sigma }+(\omega \tau )],$$
(A8)
$`H_{\mu \nu \rho \sigma \omega \tau }(k_1,k_2)`$ $`=`$ $`[C_{\mu \nu ,\rho \tau }k_{1\sigma }k_{2\omega }+C_{\mu \nu ,\rho \omega }k_{1\tau }k_{2\sigma }C_{\mu \nu ,\omega \tau }k_{1\rho }k_{2\sigma }+(\rho \sigma )]`$ (A10)
$`[(\mu ,\nu )(\rho ,\sigma )],`$
and
$`I_{\mu \nu \rho \sigma \omega \tau }(k_1,k_2)=`$ (A11)
$`\{[(C_{\sigma \omega ,\nu \tau }\eta _{\sigma \tau }\eta _{\nu \omega })k_{1\mu }k_{2\rho }+(C_{\nu \omega ,\sigma \tau }\eta _{\sigma \omega }\eta _{\nu \tau })k_{1\rho }k_{2\mu }+(\mu \nu )]+(\rho \sigma )\}+C_{\mu \nu ,\rho \sigma }k_{1\tau }k_{2\omega }.`$ (A12)
## B
We provide details of the treatment of the infrared singularity for the case $`n=2`$. From the expression for the self-energy contribution from the seagull diagram, Eq. (8), we extract the integral
$`{\displaystyle _0^1}{\displaystyle \frac{xy^{3+\frac{n}{2}}}{x+y}}\left(x+3y\right)\left(x+6y\right)𝑑y`$ $`=`$ $`\left({\displaystyle \frac{2x^2}{n4}}\right){}_{2}{}^{}F_{1}^{}[1,(4+n)/2,(2+n)/2,1/x]`$ (B2)
$`+18\left({\displaystyle \frac{x}{n2}}+{\displaystyle \frac{1}{n}}{}_{2}{}^{}F_{1}^{}[1,n/2,(2+n)/2,1/x]\right)`$
where
$${}_{2}{}^{}F_{1}^{}[a,b,c,D]=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a)_k(b)_k}{(c)_k}\frac{D^k}{k!}$$
(B3)
is the hypergeometric function and $`(a)_k`$ is the Pochhammer symbol
$$(a)_k=\mathrm{\Gamma }(a+k)/\mathrm{\Gamma }(a).$$
(B4)
This result is convergent only for $`n>4`$. For $`n=2+ϵ`$, we analytically continue the above equation and expand it in a power series in $`ϵ`$,
$$\frac{16x}{ϵ}x\left[x\mathrm{ln}(1+x^1)\right]+𝒪(ϵ).$$
(B5)
Integrating the $`𝒪(ϵ^0)`$ term of the above expression over $`x`$, we obtain
$$_0^1𝑑xx\left[x\mathrm{ln}(1+x^1)\right]=\frac{14}{3}.$$
(B6)
We can repeat the procedure for $`\mathrm{\Pi }_R(0)`$. The $`𝒪(ϵ^0)`$ term of the integral in Eq. (10) is
$`30r^2\left[Li_2({\displaystyle \frac{1}{1r}})Li_2({\displaystyle \frac{2}{1r}})+Li_2(r^1)\right]+3r^2\left[r\mathrm{ln}{\displaystyle \frac{r}{r+1}}10\mathrm{ln}{\displaystyle \frac{2(r+1)}{r1}}\right]`$ (B7)
$`+{\displaystyle \frac{3}{2}}r\left(2r1+40\mathrm{ln}2\right){\displaystyle \frac{1}{3}},`$ (B8)
where $`Li_2`$ is the dilogarithm function. The same procedure can be carried out for the integral in $`\mathrm{\Pi }_R(m^2)`$, but the expressions are exceedingly cumbersome and not at all illuminating. We do not show them here.
The point we wish to emphasize is that by performing dimensional regularization in the number of extra spacetime dimensions $`n`$, we are able to systematically isolate the infrared divergence and calculate finite quantities with a sensible physical meaning. Using the results in this Appendix, we showed that the oblique corrections vanish identically in the limit $`M_S\mathrm{}`$, as one would expect.
We have demonstrated this procedure only for $`n=2`$, but it can be carried out for any $`n`$. For example, for $`n=3`$, we would substitute $`n=ϵ+3`$ and so on. |
warning/0001/cond-mat0001176.html | ar5iv | text | # Extensive infrared spectroscopic study of CuO: signatures of strong spin-phonon interaction and structural distortion
## I Introduction
Since 1986 the interest to cupric oxide CuO has been mostly governed by its close relation to the problem of high-$`T_c`$ superconductivity. In addition to the role of the parent compound of all the high-$`T_c`$ materials with CuO<sub>2</sub> planes, it has a number of physical and chemical features common to several undoped antiferromagnetic (AFM) cuprates, (e.g. La<sub>2</sub>CuO<sub>4</sub>, YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6</sub>): similar copper coordination and electronic state, Cu-O distances, values of localized magnetic moments, superexchange constants, low-dimensionality of magnetism etc.
CuO, however, is a quite interesting system in its own right. Although Cu<sup>2+</sup> ions are expected to be in the $`3d^9`$ state with one 3$`d`$-hole per atom, this transition metal (TM) oxide is a strongly correlated insulator of the ”charge-transfer” type according to the theory of Zaanen, Sawatzky and Allen ; the holes are well localized forming local magnetic moments. CuO undergoes a 2-stage magnetic transition: at $`T_{N1}`$ = 230 K an incommensurate magnetic structure is observed, while at $`T_{N2}`$ = 213 K magnetic moments order parallel to the b-axis antiferromagnetically along the $`[10\overline{1}]`$ chains and ferromagnetically along the chains . From the analysis of the spin-wave velocity it was found that the exchange constant along the $`[10\overline{1}]`$ chains (60 - 80 meV) is several times larger than this value along any other direction. The anomalous temperature dependence of the magnetic susceptibility points to low-dimensional, or, at least, highly anisotropic character of magnetic interactions and persistency of spin correlations at temperatures well above the Néel point .
Another feature of the cupric oxide is the low-symmetry monoclinic lattice, which distinguishes it from the other TM monoxides, e.g. MnO, FeO, CoO and NiO with the rock-salt structure. It is a prominent manifestation of the Jahn-Teller effect: in the high-symmetry octahedral position characteristic to the cubic structure, the Cu<sup>2+</sup> ion would have degenerate $`d_{x^2y^2}`$ and $`d_{z^2}`$ orbitals, which is energetically unfavourable, and therefore tend to displace away from the symmetry position. This tendency is so strong that CuO has not just a distorted cubic lattice, but a completely different monoclinic tenorite structure.
Several groups have reported results of infrared (IR) spectroscopic studies of powder as well as single-crystal specimens of CuO. The interpretation of infrared spectra was always embarrassed by the low crystal symmetry, especially for the case of polycrystalline samples. Kliche and Popovic have measured infrared spectra of sintered powder samples as a function of temperature and for the first time assigned strong IR-active modes to the species $`A_u`$ and $`B_u`$ by comparison of frequencies with those in PdO. They also reported an additional broad mode at about 414 cm<sup>-1</sup> the intensity of which increases drastically with cooling down below $`T_N`$ and suggested that it is a zone-boundary phonon mode which becomes IR-active because of IR absorption from AFM superstructure. It is evident now that it was a manifestation of the anomalous softening of the $`A_u^3`$ mode reported by Homes et al .
So far it was a serious problem to obtain single crystals of CuO suitable for quantitative infrared studies. Guha et al has succeeded to measure infrared polarized spectra of single crystals of CuO at room temperature and account for low-symmetry effects in data analysis. They have measured reflectivity from the (1$`\overline{1}`$0) natural face and modelled spectra by the dielectric function formulas adapted to monoclinic crystals . However, due to inconvenient crystal orientation in their experiment mixed LO-TO modes were excited, the properties of which depend on the wave vector direction.
Homes et al were the first to present single-crystal infrared spectra as a function of temperature. Again, however, only the (1$`\overline{1}`$0) crystal surface was accessible for optical experiments, and mixed LO-TO modes were actually observed. An appreciable (about 5 %) sharp softening of the $``$ 440 cm<sup>-1</sup> reststrahlen band for the E $``$ b upon cooling down was definetely registered at the Néel transition. Spectra were fitted with introduction of $`3A_u`$ and $`3B_u`$ phonon modes only. No new phonon structures at the magnetically ordered phase were reported indicating absence of a crystal superlattice below $`T_{N2}`$.
This statement sounds puzzling in a view of observation by Chen et al of five new modes at low temperatures in the Raman spectra. Authors have assigned these modes to folded phonons; as a folding mechanism, a strong spin-phonon interaction was proposed. The most intense new mode 240 cm<sup>-1</sup> hardens strongly at cooling down, which was attributed to an additional lattice rigidity due to magnetization.
There is a serious inconsistency concerning structure and parameters of IR-active phonon modes, especially at high frequencies. For instance, the deviations in resonance frequency of these modes reported by different groups are too significant to be explained by experimental errors, isotope effect, crystal non-stoiochiometry etc. In our opinion, the explanation lies in the intermediate LO-TO nature of the observed modes and corresponding uncertainty of phonon parameters, especially for the high-frequency intense modes with large LO-TO splitting. Moreover, no infrared data so far were reported where the $`A_u`$ and $`B_u`$ modes were completely separated. Parameters of the $`A_u`$ modes were extracted at best from the single-crystal spectra for $`𝐄c`$ where the $`B_u`$ modes are also present.
In this paper we aimed to resolve this uncertainty by separate measurement of the characteristics of purely TO $`A_u`$ and $`B_u`$ modes. For monoclinic crystals the only option for observation of the $`B_u`$ TO modes is to measure normal reflectivity from the (010) face (the ac-plane). To study the $`A_u`$ TO modes any crystal plane containing the b-axis, e.g. (001) face, may suffice. We succeeded to obtain these crystal faces with a sufficiently large area, allowing to perform reliable measurements and quantitative analysis of the data as described below.
## II Crystal structure and factor-group analysis
Cupric oxide CuO, unlike other TM monoxides, crystallizes in a low-symmetry monoclinic tenorite structure (Fig. 1). It is generally accepted, following Åsbrink and Norrby , that at room temperature (RT) the space group is C$`{}_{}{}^{6}{}_{2h}{}^{}`$ (C2/c); there are four CuO molecules in the unit cell with dimensions a = 4.6837 Å, b = 3.4226 Å, c = 5.1288 Å, $`\beta `$ = 99.54 and two CuO units in the primitive cell; the copper and oxygen occupy the C<sub>i</sub> and C<sub>2</sub> symmetry positions correspondingly. Each copper atom is situated in the center of the oxygen parallelogram. Each oxygen atom, in turn, has a distorted tetrahedral copper coordination. The adjacent CuO<sub>4</sub> parallelograms form two sets of ribbons propagating along the and the \[1$`\overline{1}`$0\] directions. The structure can be also considered as being composed from two types of zig-zag Cu-O chains running along the and the \[10$`\overline{1}`$\] directions (Fig. 2). The Cu-O-Cu angle is 146 in the \[10$`\overline{1}`$\] chains and 109 in the chains.
For the C$`{}_{}{}^{6}{}_{2h}{}^{}`$ space group the factor-group (FG) analysis gives the following set of the zone-center lattice modes: $`\mathrm{\Gamma }=A_g+2B_g+3A_u+3B_u+3\text{ translational}`$. Out of these, 3 modes ($`A_g+2B_g`$) are Raman-active, 6 modes (3$`A_u+3B_u`$) are IR-active. The $`A_u`$ modes are polarized along the b-axis. The dipole moments of the $`B_u`$ modes lie within the ac-plane, but due to the low symmetry their directions are not exactly determined by the crystal structure.
More recently Åsbrink and Waskowska have refined the CuO structure at 196 K and 300 K using the so called ”less significant reflections” in the X-ray data analysis and found that less symmetric space-group C$`{}_{}{}^{4}{}_{s}{}^{}`$ (Cc) is also consistent with the X-ray diffraction data for both phases. They suggested that the C$`{}_{}{}^{6}{}_{2h}{}^{}`$ space-group might result from the time-averaging or site-averaging of non-equivalent (due to valence fluctuations) atom positions of lower symmetry. Some lattice distortions, especially changes of the Cu-O distances, were clearly detected when passing from RT to 196 K. In general, one can state, that the C$`{}_{}{}^{6}{}_{2h}{}^{}`$ space-group is a good approximation to the real structure of the cupric oxide, but some minor deviations from this do not contradict to the X-ray data.
## III Experimental
### A Sample preparation and characterization
Single crystals of CuO were obtained from a CuO - PbO - Bi<sub>2</sub>O<sub>3</sub> melt. The details are described elsewhere . After cooling down the crucible contained randomly oriented large single-crystal pieces of CuO along with inclusions of other phases. From this conglomerate the largest CuO single crystals were extracted and oriented using the X-ray diffraction. As usual, natural crystal faces were presumably of (110) and (1$`\overline{1}`$0) orientation. This face orientation was used in previous papers where infrared reflectivity measurements were performed. However, for the reasons mentioned above, we aimed to obtain large enough the (001) (the ab-plane) and the (010) (the ac-plane) crystal faces. These two mutually perpendicular faces were cut on one selected single-crystal sample, which was used for measurement of all the reflectivity spectra presented in this paper. Cuts were polished with a fine 0.06 $`\mu `$m Al<sub>2</sub>O<sub>3</sub> powder. Microscopic analysis of the surface has shown that the crystal is twinned. Fortunately, one twin orientation was almost completely dominating; the domains of the alternative twin orientation form narrow stripes covering less than 5 % of surface area. Such domination was also confirmed by the X-ray Laue snapshots, where no detectable reflections corresponding to the alternative twin orientation were observed. The structure of twins is such that the b-axis direction is the same for different twin orientations; all twin reflections are within the ac-plane. So a minor (less than 5 %) contribution of other twin domains is possible for the (010) face reflectivity spectra. For the case of reflection from the (001) face, when E $``$ b, all twins contribute in the same way and twinning has no effect. The Lauegram has shown that (001) and (010) crystal faces were cut with accuracy of 1.8 and 2.1respectively. The electron Auger microscopy has shown the presence of only copper and oxygen atoms on the both crystal faces. The area of the (001) face suitable for quantitative optical measurements (i.e. containing no impurity inclusions, having the lowest fraction of the alternative twin orientation), was $``$ 3 mm<sup>2</sup>; that of the (010) face was $``$ 4 mm<sup>2</sup>.
### B Reflectance measurement
Infrared reflectivity spectra were measured from 70 to 6000 cm<sup>-1</sup> using a Bruker IFS 113v FT-IR spectrometer. The average angle of incidence was about 11. A set of different light sources, beamsplitters, polarizers and detectors were used to cover this frequency range. The mid-infrared (MIR) spectra from 400 to 6000 cm<sup>-1</sup> were measured using a globar source, KBr beamsplitter, KRS-5 polarizer and DTGS and MCT detectors. The far-infrared (FIR) region 70 - 700 cm<sup>-1</sup> was studied with the aid of the Hg lamp, a set of mylar beamsplitters, a polyethylene polarizer and the helium-cooled Si bolometer.
The polarizer was mounted in the optical path of the incident beam; no additional polarizers (analyzers) were put inbetween sample and detector. The transmission properties of the polarizers were measured independently and, when necessary, special care of the correction for the unwanted polarization leakage was taken. Polarizer rotation was performed using a computer-controlled mechanical rotator.
An original ”three-polarization” measurement technique was used involving the measurement of three reflectivity spectra per crystal face for different polarizations of almost normally incident light: vertical (0), horizontal (90) and diagonal (45). In principle, the knowledge of these spectra should be enough to calculate the reflectivity for any other polarization direction. In particular, the relation $`R(0^{})`$+$`R(90^{})`$ = $`R(45^{})`$+$`R(45^{})`$ should work. We especially checked the validity of this relation and experimentally proved that it holds with a good accuracy.
The sample was mounted with a good thermal contact in a continuous-flow cryostat (Oxford Instruments) with an automatic temperature control. Spectra were measured at temperatures 300, 250, 240, 230, 220, 210, 200, 180, 150, 100 and 7 K, so that, special attention has been paid to the range in the vicinity of $`T_{N_1}`$ = 230 K and $`T_{N_2}`$ = 213 K. The temperature setting accuracy was about 1 K.
A reference for the absolute reflectivity was provided by in situ evaporation of a gold layer on the surface and consecutive repetition of the same set of measurements for every temperature. Such a procedure has compensated errors associated with not only non-ideality of the sample face but also the thermal deformation of the cryostat cold finger. To account for a possible drift of a single-beam intensity due to source and detector instability, every sample-channel measurement was accompanied by a measurement of the intensity of the light beam passed via the second channel without reflection from the sample.
## IV Spectra treatment
### A The dispersion analysis
Let us introduce the orthogonal system of coordinates {$`\mathrm{𝐱𝐲𝐳}`$}: $`𝐱𝐚`$, $`𝐲𝐛`$, $`𝐳𝐚`$, $`𝐳𝐛`$ so that there is a slight inclination ($``$ 9.5) between axes z and c. Due to the monoclinic symmetry the whole 3D dielectric tensor $`\widehat{ϵ}`$ is the composition of two components: the scalar $`ϵ_b=ϵ_{yy}`$ along the b-axis and the 2D tensor $`\widehat{ϵ}_{ac}=\left[\begin{array}{cc}ϵ_{xx}& ϵ_{xz}\\ ϵ_{zx}& ϵ_{zz}\end{array}\right]`$ within the ac-plane ($`ϵ_{xz}=ϵ_{zx}`$ is expected without external magnetic field). The dispersion formulas are:
$$ϵ_b(\omega )=ϵ_b^{\mathrm{}}+\underset{i,A_u}{}\frac{\omega _{\text{p},i}^2}{\omega _{\text{TO},i}^2\omega ^2\text{i}\gamma _i\omega }\text{ , }$$
(1)
$`\widehat{ϵ}_{ac}(\omega )`$ $`=`$ $`\widehat{ϵ}_{ac}^{\mathrm{}}+{\displaystyle \underset{i,B_u}{}}{\displaystyle \frac{\omega _{\text{p},i}^2}{\omega _{\text{TO},i}^2\omega ^2\text{i}\gamma _i\omega }}\times `$ (2)
$`\times `$ $`\left[\begin{array}{cc}\mathrm{cos}^2\theta _i& \mathrm{cos}\theta _i\mathrm{sin}\theta _i\\ \mathrm{cos}\theta _i\mathrm{sin}\theta _i& \mathrm{sin}^2\theta _i\end{array}\right]\text{,}`$ (5)
where $`\omega _{\text{TO},i}`$ \- the transverse frequency, $`\omega _{\text{p},i}`$ \- the plasma frequency, $`\gamma _i`$ \- the linewidth of the $`i`$-th mode, $`\theta _i`$ \- the angle between the dipole moment of the $`i`$-th mode and the $`x`$-axis (for the $`B_u`$ modes only), $`ϵ_b^{\mathrm{}}`$ and $`\widehat{ϵ}_{ac}^{\mathrm{}}`$ are the high-frequency dielectric tensors. The b-axis complex reflectivity $`r_b`$ and the reflectance of the (001) plane for E $``$ b are expressed via the dielectric function:
$$r_b(\omega )=\frac{1\sqrt{ϵ_b(\omega )}}{1+\sqrt{ϵ_b(\omega )}}\text{ , }R_b(\omega )=|r_b(\omega )|^2\text{.}$$
(6)
The complex reflectivity tensor $`\widehat{r}_{ac}`$ can be expressed via the dielectric tensor $`\widehat{ϵ}_{ac}`$ by the matrix formula, which is formally analogous to (6):
$$\widehat{r}_{ac}(\omega )=(\widehat{1}\sqrt{\widehat{ϵ}_{ac}(\omega )})(\widehat{1}+\sqrt{\widehat{ϵ}_{ac}(\omega )})^1\text{,}$$
(7)
where $`\widehat{1}`$ is the unity tensor. The matrix square root naturally means, that the matrix is first reduced to the diagonal form by a proper rotation, the square root is then taken from each diagonal element, and finally it is rotated back to the initial coordinate system. The ”-1” exponent implies calculation of the inverse matrix.
The reflectance of the (010) plane depends on the direction of the incident light polarization $`𝐞=𝐄/|𝐄|`$:
$$R_{ac}(\omega ,𝐞)=|\widehat{r}_{ac}(\omega )𝐞|^2\text{.}$$
(8)
The reflectances for values 0, 45and 90 of the angle between the electric field vector and the x-axis used in the ”three-polarization” measurement scheme are expressed in terms of the components of $`\widehat{r}_{ac}`$:
$`R_{00}(\omega )`$ $`=`$ $`|r_{xx}(\omega )|^2+|r_{xz}(\omega )|^2`$ (9)
$`R_{45}(\omega )`$ $`=`$ $`(|r_{xx}(\omega )+r_{xz}(\omega )|^2+`$ (10)
$`+`$ $`|r_{xz}(\omega )+r_{zz}(\omega )|^2)/2`$ (11)
$`R_{90}(\omega )`$ $`=`$ $`|r_{xz}(\omega )|^2+|r_{zz}(\omega )|^2\text{.}`$ (12)
The phonon parameters can be obtained by fitting of the reflectance spectra using the written above formulas. To obtain the $`A_u`$ modes parameters the $`R_b(\omega )`$ spectrum is fitted. The characteristics of the $`B_u`$ modes, including unknown angles $`\theta _i`$, can be extracted by simultaneous fitting of three spectra $`R_{00}(\omega )`$, $`R_{45}(\omega )`$, $`R_{90}(\omega )`$ from the (010) plane.
### B The Kramers-Kronig analysis
For the case E $``$ b the Kramers-Kronig (KK) analysis can be performed in a usual way because one of the dielectric axes is parallel to this direction. Due to the low crystal symmetry, the directions of the two other principal dielectric axes in the ac-plane depend on the frequency, which precludes the straightforward application of the KK method to the ac-plane reflectance data. For this case we used a modified version of this technique, which allows to determine frequency dependence of all the components of the complex reflectivity tensor $`\widehat{r}_{ac}`$, provided that three reflectance spectra $`R_{00}(\omega )`$, $`R_{45}(\omega )`$ and $`R_{90}(\omega )`$ are measured in wide enough frequency range. The details of this KK method generalization are described in Ref..
For a correct implementation of the KK integration, the reflectivity in the range 6000 - 37000 cm<sup>-1</sup> was measured using the Woollam (VASE) ellipsometer system. At higher frequencies the $`\omega ^4`$ asymptotics was assumed, as usual. At low frequencies the reflectivity was extrapolated by a constant value.
## V Results and analysis
### A E $``$ b
The reflectance spectra for the (001) plane, when E $``$ b, are shown at Fig. 3. In this configuration only the $`A_u`$ TO modes should be active. Exactly three strong modes are observed: $`A_u^1`$ ($``$ 160 cm<sup>-1</sup>), $`A_u^2`$ ($``$ 320 cm<sup>-1</sup>) and $`A_u^3`$ ($``$ 400 cm<sup>-1</sup>), which confirms the FG-analysis predictions for the established for CuO crystal structure. The most drastic temperature changes take place in the range 350 - 550 cm<sup>-1</sup>, where the reststrahlen band corresponding to very intense lattice mode $`A_u^3`$ is situated. The reflectivity maximum elevates from 65 % to almost 100 %; its gravity center moves to lower frequencies upon cooling down the sample. It indicates, that this mode experiences strong softening and narrowing as a temperature is decreased.
In addition to three strong modes, at least five ”extra” structures in these reflectivity spectra are detectable ”by eye” (Fig. 3). The first structure is a dip at $``$ 425 cm<sup>-1</sup> just on the top of the reststrahlen band, which becomes visually evident below 210 K. The second is a $``$ 485 cm<sup>-1</sup> structure also on the top of the same reststrahlen band seen at 100 K and 7 K. The third feature is a $``$ 630 cm<sup>-1</sup> peak (see inset), which is very small (but observable) at 300 K, and becomes significant at low temperatures. The fourth structure is high-frequency mode $``$ 690 cm<sup>-1</sup> (see inset), which is very obvious at 7 K (690 cm<sup>-1</sup>) and 100 K (680 cm<sup>-1</sup>), still detectable at 150 K ($``$ 650 cm<sup>-1</sup>) and not seen at higher temperatures, probably, because of strong broadening. The fifth structure is seen at 165 - 170 cm<sup>-1</sup> at the top of the reststrahlen band of the $`A_u^3`$ mode (Fig. 4(a)). It is better observable at low temperatures; but even at room temperature the form of the reststrahlen band differs from a single-mode shape.
The observations ”by eye” should be accompanied by numerical analysis. The dispersion analysis of spectra was performed in two stages. On the first stage, in order to determine the characteristics of the principal modes, we have fitted the reflectivity curves with introduction of 3 modes only. The experimental and fitting curves are compared at Fig. 5(a) for the $`T`$ = 100 K. The fit quality is good enough; the deviations are observed only in the range of additional modes. A relative weakness of additional structures ensures small errors in determination of main mode parameters. The parameters of 3 principal $`A_u`$ phonon modes obtained by such a fitting as a function of temperature are shown at Fig. 6.
From Fig. 6 a conclusion can be drawn immediately that the $`A_u^1`$ and $`A_u^2`$ modes behave in a quite similar way, while the highest-frequency $`A_u^3`$ mode is significantly different. The $`A_u^1`$ and $`A_u^2`$ modes are steadily hardening with cooling ($``$ 1 %), with some increasing of the slope $`\omega _{\text{TO}}/T`$ below $`T_N`$. In contrast, the $`A_u^3`$ mode slightly softens with cooling down to the $`T_{N1}`$, then undergoes a drastic sharp softening ($``$ 10 %) in the vicinity of the transition temperature and then hardens with further cooling from $`T_{N2}`$ down to the helium temperature. The $`A_u^2`$ and $`A_u^3`$ modes are relatively narrow at 300 K and exhibit further narrowing with cooling with some dip at the transition temperature. On the contrary, the $`A_u^3`$ mode is very broad at 300 K, and broadens more with approaching the $`T_N`$. Its linewidth has a pronounced maximum inbetween $`T_{N1}`$ and $`T_{N2}`$. In the AFM phase it quickly narrows with cooling down.
On the second stage, in order to determine or, at least, estimate parameters of the mentioned ”extra” modes, additional fitting of spectra with introduction of both principal and ”extra” modes has been performed (see Fig. 7). It was possible to fit the structures at 167 cm<sup>-1</sup>, 425 cm<sup>-1</sup>, 485 cm<sup>-1</sup> and 690 cm<sup>-1</sup>. The 630 cm<sup>-1</sup> peak cannot be fitted by the usual lorentian term. The first fit was performed at 7 K, when additional structures are mostly sharp. Other spectra were fitted in sequence 100 K, 150 K, …, 300 K. At each step the oscillator parameters, corresponding to the previous temperature served as initial approximation for the least-squares fitting. Parameter confidence limits were always calculated by the ”covariant matrix” method (see error bars at Fig. 7), which takes into account possible correlation of parameters. In this way it was possible to extend curves of modes 167 cm<sup>-1</sup>, 425 cm<sup>-1</sup> and 480 cm<sup>-1</sup> up to room temperature; the 690 cm<sup>-1</sup> mode was fitted only for $`T`$ 150 K, because at higher temperatures fitting could not give reasonable values of parameters for this mode. The errors in determination of additional modes are relatively larger than those for the principal ones, which is natural, in view of their small intensity. In general, errors are smaller at lower temperatures.
The temperature dependence of TO frequency additional mode 167 cm<sup>-1</sup> is more or less typical for phonons (the value of 171 cm<sup>-1</sup> at 7 K is an artefact of the dispersion analysis with no physical meaning). On the contrary, the 425 cm<sup>-1</sup> and 485 cm<sup>-1</sup> modes demonstrate a puzzling temperature dependence, similar to that of the $`A_u^3`$ mode 400 cm<sup>-1</sup>. First of all, there is no indication of disappearance of these modes above the AFM transition, although they are not clearly seen ”by eye”. Nevertheless, the transition strongly affects parameter values of these modes. Both modes strongly soften above $`T_N`$ and strongly harden below $`T_N`$. The modes are narrowing below $`T_N`$ (which facilitates their visual observation), and are very broad at higher temperatures with a strong maximum at the transition point. One can content only by a qualitative conclusions, because the temperature dependencies of these modes are masked by large error bars. It can be explained by some correlation between parameters of these modes with ones of the $`A_u^3`$ mode. It is a typical problem for several broad closely located modes. Therefore, it is unreasonable to attribute physical meaning to increasing of the plasma frequency of these modes above 200 K: their plasma frequencies are just subtracted from the plasma frequency of the $`A_u^3`$ mode without significant influence on the fit quality. For the 690 cm<sup>-1</sup> mode the dispersion analysis results confirm visual observations: anomalously strong hardening at cooling and strong broadening with heating. The latter, probably, precludes satisfactory fitting of this mode at higher temperatures.
The curves of the b-axis optical conductivity $`\sigma _b(\omega )=\omega \text{Im}ϵ_b(\omega )/4\pi `$, obtained by the KK transform of the reflectivity spectra for each temperature, are shown at Fig. 8. The conductivity is very illustrative to show a remarkable difference between the $`A_u^1`$ and $`A_u^2`$ narrow modes and the $`A_u^3`$ mode exhibiting a puzzling temperature transformation. The sharpness of the $`A_u^1`$ and $`A_u^2`$ modes and the absence of the $`B_u`$ modes contribution to this spectrum confirms a good sample quality and its proper orientation. Some deformations of the lineshapes for $`T`$ = 100 K and 7 K (in particular, a small negative value of conductivity just below the mode frequencies) are most likely the results of uncertainties of the KK method, which is very sensitive to experimental errors in the regions where the reflectivity approaches 0 or 1. We shall not attribute any physical meaning to this effect.
### B E $``$ ac
The reflectance spectra from the (010) plane for three polarizations of the incident light ($`R_{00}`$, $`R_{45}`$ and $`R_{90}`$) are shown at Fig. 9. The $`B_u`$ TO modes are expected to appear in these spectra. The narrow mode $`B_u^1`$ ($``$ 145 cm<sup>-1</sup>) is clearly observed in all polarizations. Its intensity depends, of course, on the light polarization, i.e. the angle between the electric field vector and the mode dipole moment. Strong reststrahlen bands are seen in the 450 - 600 cm<sup>-1</sup> range. The shape and the center position of the band is polarization-dependent, which is consistent with a suggestion that it is actually formed by at least two high frequency intense modes. In the $`R_{00}`$ (E $``$ a) and $`R_{45}`$ spectra one can observe some minor contribution from the $`A_u^1`$ (162 cm<sup>-1</sup>) and $`A_u^2`$ (323 cm<sup>-1</sup>) modes. The possible reason is some misorientation of the sample. A spike at about 130 cm<sup>-1</sup> is of apparatus origin and should be ignored.
As is in the case E $``$ b, some extra structures are seen. The first is a broad band in the range 380 - 425 cm<sup>-1</sup>. It is especially pronounced for E $``$ a, less evident for intermediate polarization and absent for E $``$ a (see the left insets in Fig. 9). A dip at 425 cm<sup>-1</sup>, indoubtedly correlates with the dip at the same frequency at reflectance for E $``$ b. The second is a pronounced structure consisting of a peak at $``$ 480 cm<sup>-1</sup> and a dip at $``$ 485 cm<sup>-1</sup> in the $`R_{90}`$ spectrum, existing at all temperatures. In spite of proximity of this frequency to additional structure at Fig. 3, a completely different temperature dependence makes one to separate these modes. The third structure is a dip at $``$ 507 cm<sup>-1</sup> on the top of the reststrahlen band of $`R_{00}`$ and $`R_{45}`$ (see the right insets in Fig. 9). This frequency is very close to the LO frequency of the $`A_u^3`$ mode, therefore it is most likely the $`A_u^3`$ mode which is seen in this spectrum for the same reason as the $`A_u^1`$ and $`A_u^2`$ modes are. The fourth structure is seen in the vicinity of the $`B_u^1`$ mode (Fig. 10). The shape of this mode is such that it is worth to suggest, that it is actually composed of two different modes (see, especially, Fig. 10(c).
For each temperature a fitting procedure with introduction of 3 oscillators, corresponding to principal $`B_u`$ modes has been performed. The phonon polarization angles were adjusted along with other phonon parameters. Spectra $`R_{00}`$, $`R_{45}`$ and $`R_{90}`$ were fitted at the same time. In the spirit of the FG-analysis prediciton, three lattice modes were introduced for spectra fitting: one low-frequency, and two high-frequency modes. The fit quality for T = 100 K is seen on the Fig. 5(b)-(d). One can see that the $`B_u^1`$ 145 cm<sup>-1</sup> mode as well as the general shape of the reststrahlen band at 450 - 600 cm<sup>-1</sup> are satisfactorily reproduced, confirming a suggestion that only 3 strong phonon modes are present. A bump at about 620 cm<sup>-1</sup> is fitted without invoking of additional Lorentzians: it results from the non-collinearity of the mode and the incident radiation polarizations.
The temperature dependence of the $`B_u`$ phonon parameters is presented at Fig. 11. Unlike the case of the $`A_u`$ modes, there is no significant difference in the temperature dependence of parameters of the low- and high-frequency $`B_u`$ modes. All modes are monotonically hardening with cooling down, with some positive kink at $`T_N`$ for the $`B_u^3`$ mode and negative kink for the $`B_u^2`$ mode. The plasma frequencies of all modes have slight maximum at the transition temperature. Two high frequency modes are much more intense than the $`B_u^1`$ mode. The linewidth of all the modes doesn’t decrease with cooling down. Instead, it increases at low temperatures. One should be careful in a straightforward interpretation of this result, because the lineshape is not described perfectly, especially one of the high-frequency modes. The true linewidth is better seen from the optical conductivity graph (see below). The oscillator polarization angles do not significantly change with frequency. The maximum angle change is 4 - 5, while the change of the relative angle between different oscillators polarization is less than 2 . One can state that within experimental errors oscillators angles almost doesn’t change.
To investigate the true shape of the principal phonon modes and additional structures, the KK analysis of the ac-plane data for the mentioned set of temperatures was implemented in the extended form . At Fig. 12 all the components of the optical conductivity tensor $`\widehat{\sigma }_{ac}(\omega )=\omega \text{Im}\widehat{ϵ}_{ac}(\omega )/4\pi `$ are plotted for selected temperatures. Note that the off-diagonal component $`\sigma _{xz}`$ may have any sign unlike the diagonal components $`\sigma _{xx}`$ and $`\sigma _{zz}`$ which must be positive.
The polarization angle of the $`B_u^2`$ mode is 35, which is close to the direction of the chains. The $`B_u^3`$ mode is almost orthogonal to the $`B_u^2`$ mode: the angle is -55, which is close to the \[10$`\overline{1}`$\] chains direction. Therefore, with some approximation one can state, that the $`B_u^2`$ and $`B_u^3`$ modes are stretches of the and \[10$`\overline{1}`$\] chains correspondingly, which is in agreement with several lattice dynamical calculations. Note, that such orthogonality is not determined by the crystal symmetry. The polarization angles of both modes are almost temperature-independent.
## VI Discussion
### A Comparison with previous results
A comparison of our data with previosly reported results of IR studies of CuO cannot be direct, because we measured spectra, where the $`A_u`$ and $`B_u`$ modes are separated and excited in purely transverse regime. All the quantitative and even qualitative deviations with previous data (see below) can be ascribed to a different way of spectra measurement and analysis.
At the Table I the phonon frequencies at room temperature previously obtained by means of infrared spectroscopy as well as neutron scattering are collected. It is seen that the most serious discrepancy between reported values of phonon frequencies takes place for modes $`A_u^3`$, $`B_u^2`$ and $`B_u^3`$, which are very intense and manifest the largest LO-TO frequency splitting. Pure TO mode should have the lowest possible frequency, which is in nice agreement with the current result: our reported frequencies are smaller than those reported in other IR spectroscopy papers. One should mention a much better agreement between our data and the results of neutron scattering experiments, where characteristics of pure TO modes were determined as well.
The softening of the $`A_u^3`$ mode was reported by Homes et al , who observed a sudden drop of the phonon frequency at the Néel transition from 450 to 430 cm<sup>-1</sup>, i.e. $``$ 5 %. Our data qualitatively confirm this interesting result. We observe even stronger effect: the TO frequency drops from 410 to 370 cm<sup>-1</sup> ($``$ 10 %), which is twice as large as was reported (Fig. 6).
Ref. , was thus far the only paper, to our knowledge, where IR spectra of single-crystal CuO at low temperatures were studied. Authors didn’t report any new IR-active modes below the Néel temperature; it was implied, that no additional lines are present at higher temperatures either. However, an absence of extra IR-active modes is quite strange, if one compares it with an observation of at least 5 ”unexpected” lines in the Raman spectra. Conversely, according to our data, several ”extra” modes are present in IR spectra in the whole 7 - 300 K frequency range. We believe, that better orientation of the wave vector and polarization of the incident radiation may facilitate observation and analysis of minor IR-active modes.
### B The $`A_u^3`$ mode anomaly
The $`A_u^3`$ mode among other principal IR-active modes behaves in the most anomalous way. It is demonstrated at Fig. 13, where relative RT-normalized TO frequencies and relative linewidths of all 6 principal modes are plotted together on the same graph. A close relation between magnetic ordering transition at 213 - 230 K and temperature transformations of this mode is without any doubt.
At 300 K and, especially, at the transition temperature the mode is anomalously broad (12 - 15 %), which indicates, that it is strongly coupled to some system quasiparticles. The most probable candidates are the low-energy magnetic excitations. If one propose, that there is a strong coupling between spin excitations and the $`A_u^3`$ lattice mode, the temperature transformations of the mode can be explained by reconstruction of the magnon spectrum upon cooling down below the Néel temperature. It is well established, that spin correlations in the AFM \[10$`\overline{1}`$\] chains are present well above the Néel temperature. At high temperatures one can consider the magnetic interaction to be of quasi 1D character. In the 1D Heisenberg AFM chains with S=1/2 and exchange $`J`$ there is a continuum of triplet excitations with lower and upper boundary curves $`ϵ_{\text{min}}(q)=(\pi J/2)\mathrm{sin}(q)`$ and $`ϵ_{\text{max}}(q)=\pi J\mathrm{sin}(q/2)`$. A large linewidth of the $`A_u^3`$ mode can be explained by its interaction with continuum of magnetic excitations. Below the Néel point an exchange interaction in another directions gives rise to long-range magnetic ordering and a continuum of spin excitations should collapse to the magnon dispersion curves, which, in turn, results in narrowing of the phonon mode.
The reason for an exceptionally strong interaction of the $`A_u^3`$ mode with spin excitations one should look for in analysis of its eigenvector. Several lattice dynamical calculations were performed which yielded the lattice mode eigenvectors. According to results of all the calculations this mode is characterized by the largest displacement of the oxygen atoms along the b-axis. It results also in a large dipole moment of this mode. As a consequence, the Cu-O-Cu chain angle experiences the largest variation, when the $`A_u^3`$ mode is excited. The copper spins are coupled via the superexchange interaction, which is very sensitive to the Cu-O-Cu angle. In the \[10$`\overline{1}`$\] chains the angle is equal to 146 (Fig. 2), which is close to the 180 superexchange. It gives a negative exchange constant and AFM interaction. However, the 90 superexchange is expected to be positive (an indirect confirmation is the ferromagnetic exchange along the chain with angle equal to 109), and there exists some intermediate angle, where the superexchange is changing sign. Therefore, the motion of atoms, corresponding to excitation of the $`A_u^3`$ mode can significantly vary the value of the superexchange coupling constant. A strong chain bending probably could alternate the exchange sign.
Due to energy and momentum conservation a zone-center phonon can couple to a pair of magnons (bi-magnon) having opposite wave-vectors and one-half frequency of the phonon. Another option is interaction with a single zone-center optical magnon of the same energy. The magnon dispersion curves has been studied by inelastic neutron scattering . One acoustic and one optical branch were observed; the energy of optical magnon at the $`\mathrm{\Gamma }`$ point is 5.6 THz (187 cm<sup>-1</sup>), which is very close to one-half of the $`A_u^3`$ mode frequency (370 - 380 cm<sup>-1</sup> at low temperatures), while no optical magnons near 370 - 410 cm<sup>-1</sup> were observed. Therefore, the ”bi-magnon” or, in particular ”optical bi-magnon” scenario of resonant phonon-magnon coupling is the most probable. A more detailed theory of this effect has to be fabricated.
### C Other signatures of spin-phonon interaction
The anomalous properties of the $`A_u^3`$ mode is not the only manifestation of the spin-phonon interaction in CuO (although the most prominent one). Other modes also demonstrate some anomalies, most likely related to the magnetic ordering.
One type of anomalies is unusually strong hardening of some ”extra” modes below $`T_N`$. The most outstanding example is reported in this work hardening of the 690 cm<sup>-1</sup> mode from 650 cm<sup>-1</sup> at 150 K to 690 cm<sup>-1</sup> at 7 K and strong hardening of the the Raman-active 240 cm<sup>-1</sup> mode found by Chen et al . To our mind, these effects are related. We would follow the idea of explanation given in Ref. . The phenomenon can be viewed as a non-resonant spin-phonon interaction, when spin products in the Heisenberg Hamiltonian can be approximated by their effective averages. In this case the temperature dependence of the phonon frequency is expressed by:
$$\omega _n^2(T)=\omega _{n0}^2+\underset{ij}{}C_{ij}^n𝐒_i𝐒_j(T).$$
(13)
The $`C_{ij}^n`$ coefficients are characteristics of spin-phonon interaction; in principle, they may have any sign, depending on the eigenvector of the particular phonon mode. The mode strong hardening is supposedly due to the second ”spin” term: magnetic ordering gives additional lattice rigidity.
The same phenomenon is possibly responsible for another anomaly, namely, the change of slope $`\omega _{\text{TO}}/T`$ of the principal phonon modes at the Néel point (see Figs. 6 and 11). It is seen that for the modes $`A_u^1`$, $`A_u^2`$ and $`B_u^3`$ the slope is higher in the AFM phase, while for the $`B_u^2`$ mode the slope is higher at $`T>T_N`$; for the $`B_u^1`$ mode it doesn’t change at all. Such a differences can be explained by different values and signs of coefficients $`C_{ij}^n`$.
### D ”Extra” zone-center modes and zone folding
Activation of additional phonon modes in infrared and Raman spectra is usually a signature of unit cell multiplication and zone folding. Such activation is in effect in CuO. As is stated above, several ”extra” IR-active modes are observed. In the Raman spectra 5 new modes were detected at low temperatures . Authors related their frequencies to the phonon dispersion curves obtained by inelastic neutron scattering at the zone boundary point. In Ref. this point was referred to as $`Z^{}`$. We follow the original notation and designate it by X (see Fig. 4 in Ref. ).
At Fig. 14 the Brillouin zones corresponding to several schemes of the unit cell multiplication are drawn (a projection to the ac-plane). In these schemes different symmetry points fold to the zone center. It is easy to prove, that folding of the X-point to the zone center is equivalent to disappearance of the non-trivial translation (a+b)/2, or base-centering of the space-group. In other words, activation of the phonons from the X point requires that the unit cell should become the primitive one.
In the Table II ”extra” IR active and Raman-active mode frequencies are collected. Each mode frequency is related to close phonon energy (if any) in symmetry points X(1, 0, 0), A(1/2, 0, -1/2), B(1/2, 0, 0) or C(0, 1/2, 0) at 300 K. To make comparison more reliable, the ”extra” modes frequencies are taken at temperature as close as possible to 300 K. One can see, that many modes (all Raman-active and several IR-active ones) have close analogs at the X point. At the same time, other IR-active modes (147 cm<sup>-1</sup>, 165 cm<sup>-1</sup>, 475 cm<sup>-1</sup>, 480 cm<sup>-1</sup>) could be phonons from the A point, or, in some cases, from the B or C points. The 690 cm<sup>-1</sup> mode is a special case: due to strong hardening of this mode with cooling and failure to observe it at high temperatures, the comparison to the 300 K dispersion data is impossible. In addition, the mode energy is higher than the upper limit of the reported frequency region in the neutron scattering experiments (20 THz, or 667 cm<sup>-1</sup>). The phonon dispersion curves at low temperatures in broader energy interval are desirable.
Although some ”extra” modes have analogs in the B or C points, folding of the A and X points to the zone center (the A + X folding scheme) is the simplest option to explain appearance of all ”extra” modes. It corresponds to the ”diagonal” doubling of the unit cell {a, b, c} $``$ {a+c, b, a-c} (see Fig. 14). This scenario looks attractive, because new crystal axes correspond to principal anisotropy directions for several physical properties (exchange constants, sound velocity, principal dielectric axes etc.), and the unit cell is the same as the magnetic unit cell in the AFM phase .
In the framework of this scenario the $`A_u^3`$ principal mode, the 485 cm<sup>-1</sup> and the 425 cm<sup>-1</sup> ”extra” modes should be the phonon modes from the points $`\mathrm{\Gamma }`$, A and X correspondingly belonging to the same dispersion branch, according to ”rigid-ion” modelling of dispersion curves by Reichardt et al . It is in a good agreement with some similarity of temperature dependences of parameters of these modes (see Figs. 6 and 7).
In summary, we propose, that in reality the crystal structure is more complicated, than was considered before; it is the case, to our mind, already at room temperature, because most of ”extra” IR-active modes are present in the whole temperature range. On the base of the fact that IR-active and Raman-active modes have different frequencies one can suggest that the crystal space-group is still centro-symmetric, although the copper atoms are not necessarily located in the $`C_i`$ position. However, the alternative space-group for the CuO, proposed by Asbrink and Waskowska, is not centro-symmetric (C$`{}_{}{}^{4}{}_{s}{}^{}`$); the solution of this mismatch is unclear.
One of the central issues is the mechanism of the unit cell multiplication and formation of superstructures. There exist a variety of examples of such effect, when ”extra” IR and/or Raman modes are emerging, which are undoubtedly the zone-boundary folded phonons. In the extensively studied compounds CuGeO<sub>3</sub> and $`\alpha ^{}`$-NaV<sub>2</sub>O<sub>5</sub> the unit cell doubles as a result of spin-phonon interaction, and new IR modes are observed . Another example present compounds with a charge disproportionation, out of those the BaBiO<sub>3</sub> is probably the most famous one. In this system the bismuth is disproportionated according to scheme: 2 Bi<sup>4+</sup>$``$ Bi<sup>3+</sup> \+ Bi<sup>5+</sup>, and atoms in different valence states form superstructure, which is a reason for ”extra” quite strong IR line. One cannot exclude a collective Jahn-Teller effect as an engine of superlattice formation. The question about particular type of spin-charge-lattice ordering in CuO is open. It could be closely connected with formation of inhomogeneity phases in cuprates.
### E High-frequency dielectric function
The reflectivity and the dielectric function in the mid-infrared range well above the maximum phonon energy ($``$ 0.08 eV) but below the optical gap ($``$ 1.3 eV) are determined by the electronic polarizability. We observe significant difference between values of mid-infrared reflectivity for different polarizations. It results in an appreciable anisotropy of the high-frequency dielectric tensor. The temperature dependence of the mid-IR reflectivity is small, therefore, we discuss only room-temperature data. All the components of the $`\widehat{ϵ}^{\mathrm{}}`$ can be found from the dispersion analysis of spectra, see formulas (1), (2).
The smallest value of $`ϵ^{\mathrm{}}`$ is observed along the b-axis: $`ϵ_b^{\mathrm{}}=5.9`$; it is apparently one of the prinicipal values of the dielectric tensor. The diagonalization of the tensor gives directions and values of maximal and minimal levels of $`ϵ^{\mathrm{}}`$ within the (ac)-plane. The maximal value $`ϵ_{\text{max}}^{\mathrm{}}=7.8`$ is observed along the direction, corresponding to the angle $`\varphi _{\text{max}}=36^{}`$, the minimal value $`ϵ_{\text{min}}^{\mathrm{}}=6.2`$ is along the orthogonal direction for $`\varphi _{\text{min}}=54^{}`$. One can see that direction $`\varphi _{\text{max}}`$ is very close to direction of the \[10$`\overline{1}`$\] chains ($`\varphi _{\text{[10}\overline{1}\text{]}}=42.5^{}`$) while $`\varphi _{\text{min}}`$ almost exactly corresponds to $`\varphi _{\text{[101]}}=52.8^{}`$. In other words, the high-frequency dielectric constant within the (ac)-plane is maximal along the \[10$`\overline{1}`$\] direction, and minimal along the direction. Authors have measured the high-frequency dielectric function of polycrystalline sample by the ellipsometric method. Their reported value $`ϵ^{\mathrm{}}=6.45`$ is in a good agreement with the average quantity $`(ϵ_{xx}^{\mathrm{}}+ϵ_{yy}^{\mathrm{}}+ϵ_{zz}^{\mathrm{}})/3=6.6`$ obtained here.
### F Effective charges
Another important value derived from infrared spectra is an effective ionic charge. There are several definitions of the effective charge. We would mention the Born, or ”transverse” charge $`e_\text{T}^{}`$, the Szigeti charge $`e_\text{s}^{}`$ and the Scott charge $`Ze`$ . The ”transverse” charges can be calculated using the sum-rule:
$$\underset{i}{}\omega _{\text{p},i}^2=\frac{4\pi }{v_c}\underset{k}{}\frac{(e_{\text{T},k}^{})^2}{m_k},$$
(14)
where $`v_c`$ is the primitive cell volume, the sum in the left side is over IR-active modes, the sum in the right side is over all atoms in the primitive cell. For the binary compound this relation in combination with the electric neutrality condition directly yields the ”transverse” charge of both atoms. In CuO due to anisotropy the transverse charge slightly differs for different directions. The values along the (b-axis) the $`[101]`$ and the $`[10\overline{1}]`$ directions at room temperature are 1.92, 2.04 and 2.10 correspondingly (the $`\omega _{\text{p},i}`$ are taken from the Figs. 6, 11).
The Scott charge and Szigeti charges are related to the ”transverse” charge:
$$Ze=\frac{e_\text{T}^{}}{\sqrt{ϵ^{\mathrm{}}}},e_s^{}=\frac{3e_\text{T}^{}}{ϵ^{\mathrm{}}+2}.$$
(15)
The value of $`Ze`$ divided to the nominal valence is often considered as the degree of ionicity . For the CuO it appears to be about 40 %. The Szigeti charge is useful in the context of the polarized point ions model . In literature all these charges are used, so that it is reasonable to report all of them.
In the Table III oxygen effective charges for the CuO along with other related copper oxides (Cu<sub>2</sub>O, La<sub>2</sub>CuO<sub>4</sub>, Nd<sub>2</sub>CuO<sub>4</sub> and YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6</sub>) are collected. It is seen that charge values in CuO reasonably agree to the corresponding values in other relevant oxides.
## VII Conclusions
We have measured far- and mid-infrared reflectivity spectra of monoclinic CuO from the (010) and (001) crystal faces in wide temperature range. We obtained for the first time characteristics of pure TO $`A_u`$ and $`B_u`$ modes separately. Our data finally confirm that there are 3$`A_u`$ \+ 3$`B_u`$ intense modes, in accordance with prediction of the FG-analysis for the $`C_{2h}^6`$ space-group.
We report existence of several ”extra” less intense IR-active modes in CuO. Analysis of the phonon dispersion curves leads to the conclusion that each ”extra” IR-active as well as reported earlier Raman-active mode could be folded phonon from either X (1, 0, 0) or A (1/2, 0, 1/2) symmetry points. Such folding is compatible with the ”diagonal” doubling of the unit cell with the new basis {a+c, b, a-c}. So the space-group in reality is lower than that was considered, but still centro-symmetric.
The 690 cm<sup>-1</sup> ”extra” IR-active mode exhibits anomalous hardening, similar to behaviour of the 240 cm<sup>-1</sup> Raman mode; the reason could be in additional rigidity of lattice due to magnetization, a special manifestation of the spin-phonon interaction. Another effect, which can be explained in a similar way, is a slope change of the phonon frequencies vs. temperature at the Néel point.
The anomalous softening and narrowing of the $`A_u^3`$ mode 410 cm<sup>-1</sup> we explain by its strong resonance coupling to the optical or acoustic bi-magnons. Reconstruction of magnetic excitations spectrum at the AFM transition strongly affects the phonon characteristics.
In summary, the CuO demonstrates a variety of anomalous properties, which show complex interplay of spin, charge, and phonon subsystems already in the simplest copper(II) oxide. A further insight into physics of CuO may contribute to elaboration of non-contradictory picture of antiferromagnetism and superconductivity in the high-$`T_c`$ cuprates.
###### Acknowledgements.
This investigation was supported by the Netherlands Foundation for Fundamental Research on Matter (FOM) with financial aid from the Netherlandse Organisatie voor Wetenschappelijk Onderzoek (NWO). The activity of A.B.K., E.A.T. and A.A.B. was also supported by the Russian Foundation for Basic Research (RFBR), grant No 99-02-17752. A lot of thanks we address to H. Bron and F. van der Horst (University of Groningen) for invaluable help in samples characterization. |
warning/0001/astro-ph0001180.html | ar5iv | text | # SIMULATING SUPERNOVAE REMNANTS IN GAS CLOUDS
## 1 INTRODUCTION
The mechanism by which energy and metals are fed back into the inter-stellar medium (ISM) by the supernovae of massive stars is an interesting and complex problem with applications to many areas of astrophysics. In this paper we present an adaptation of the Hydra $`N`$-body hydrodynamics code (Couchman, Thomas & Pearce 1995) which is capable of modelling the effects of multiple interacting supernovae over a large range of scales.
Hydrodynamical simulations of galaxy evolution and cosmological structure formation have recently been a topic of particular interest (eg. Katz 1992; Katz, Hernquist & Weinberg 1992, 1999; Navarro & White 1993; Steinmetz & Muller 1995; Frenk et al. 1996; Pearce et al. 1999). Feedback from supernovae appears to be required to reheat gas and prevent the ‘cooling catastrophe’ (White & Rees 1978; White & Frenk 1991), caused because cooling is very efficient in small halos where the cooling time is less than the dynamical time. As the particle mass in a simulation increases, more and more of these halos are resolved and a greater fraction of the available gas cools, in contradiction with the observation that only a small fraction of baryonic material is in the form of stars or cold gas. Presently the efficiency at which energy is returned to the ISM and beyond to the intergalactic medium (IGM) is a free parameter with little more than handwaving arguments to support the values adopted (eg. Katz 1992; Navarro & White 1993; Mihos & Hernquist 1994; Gerritsen & Icke 1997). It is to be expected that such efficiency parameters depend upon the density and metallicity of the surrounding gas at the very least. One of the main aims of this project is to better constrain the values of these feedback parameters (the amount of energy and material returned to the ISM as well as the spatial extent and form of the feedback) for a wide range of stellar systems.
An understanding of the early evolution of star clusters requires knowledge of how the residual gas, left within a cluster after star formation has occurred, is expelled (eg. Tenorio-Tagle et al. 1986; Goodwin 1997a, b). Such an understanding would also help provide constraints on the initial conditions and formation models of star clusters (Goodwin 1997b). Within a larger context, information on the effects of multiple supernovae is required to construct detailed stimulated star formation models.
Much detailed work (both hydrodynamic simulations and analytic calculations) has been carried out upon the dynamics of individual supernovae expanding into a uniform ISM (see Lozinskaya 1992 and references therein), however, little of this work has concentrated upon the late phases of that expansion (exceptions including Cioffi, McKee & Bertschinger 1988; Slavin & Cox 1992; Thornton et al. 1998). Even less has been published on the effects of supernovae in gas clouds; analytic calculations were made of the effects of multiple supernovae on gas clouds by Dopita & Smith (1986) and Morgan & Lake (1989), but due to the complexity of the problem the treatment was understandably simplistic. More recently Petruk (1999a, b) has published simulations of the very early evolution of a SNR in a density gradient.
In this paper we describe the extensions we have made to the Hydra code to model this problem and present a number of convergence tests that we have performed. We also compare the results from this code to previous analytic and simulation work on individual supernovae. In section 2 we describe the testing of the code and show that it replicates the behaviour of a supernova remnant expanding into a uniform medium as studied by other authors. In section 3 we study the effects of a central supernova in Plummer clouds of various masses and characteristic radii. In section 4 conclusions are drawn and the future applications of the code detailed.
## 2 TESTING
In this section we present the code we have used and test both its self-consistency and it’s ability to reproduce known results for the evolution of SNRs in a uniform density ISM. Firstly we outline the modifications to the Hydra code. Then we show that the evolution of the energy content of an SNR and its expansion law is in good agreement with other simulations. We show that these laws hold for different ambient densities and that the properties of an SNR scale with its initial radius.
We do this to show that we are confident of the code’s ability to model the evolution of a SNR when applied to the problems of SNR evolution in a cloud in Section 3 (and in future papers).
### 2.1 Hydra
The public release version of Hydra, upon which the code used in this work is based, has been extensively tested in a variety of astrophysical fields (Couchman 1991, Couchman, Thomas & Pearce 1995, Thacker et al. 1999, Kay et al. 1999, Benson et al. 1999). It is an adaptive particle-particle, particle-mesh code that includes smoothed particle hydrodynamics (SPH) (Gingold & Monaghan 1977; Lucy 1977) to follow the gaseous component.
For this work we have chosen to use the latest public release version of Hydra, (Hydra3.0). This version incorporates a standard pairwise Monaghan type $`𝐫.𝐯`$ viscosity rather than one based on the divergence of the local flow as used in earlier release versions of Hydra, because the $`𝐫.𝐯`$ viscosity provides better shock capturing. We also employ the improved neighbour counting estimator as described by Thacker et al. (1999). To follow the radiative cooling of the gas we use the cooling rates calculated by Sutherland & Dopita (1993) interpolated to the chosen gas metallicity.
In addition we have rewritten the SPH algorithm in order to perform complete neighbour counting even in low density regions. Previously, neighbour searching within Hydra was limited by the search length imposed by the size of the particle-particle grid, resulting in an uncomfortably high minimum resolved density. For our problem we wish to resolve the hot, low density gas within expanding supernovae shells and so have removed the maximum search-length restriction.
### 2.2 Individual supernovae
Here the mechanism for adding individual SNRs into a simulation is outlined.
The initial stages of a SNRs evolution are complex as an enriched, very high temperature gas ploughs into the surrounding ISM. We do not wish to simulate this phase of SNR evolution. In order to deal with multiple SNRs over a large range of length scales we ignore evolution prior to the onset of the late stages of evolution characterised by the pressure-driven snowplough (PDS) phase (cf Cioffi et al. 1988).
The choice of the beginning of the PDS phase as the starting point of the evolution within the code is due to a number of considerations. Firstly we wish to simulate the evolution of many SNRs in complex environments - the mass and time resolution (as well as the complex input physics) required to simulate the earlier stages of evolution are beyond our remit. Secondly, the evolution of a SNR to the PDS phase is reasonably well understood and provides a good basis for our simulations, allowing us reduce the mass resolution required to model its evolution. Lastly, the evolution of a SNR through the adiabatic Sedov phase even in a strong density contrast such as would be found in a molecular cloud is expected to be roughly spherical; significant non-spherical evolution would only be expected in the late PDS phase (Dopita & Smith 1986). In dense environments, such as molecular clouds, the radius at which the PDS phase starts is much less than 1 pc, thus for significant non-spherical evolution during the Sedov phase the density change must be on a similarly small scale.
The early evolution of SNRs comprise a short-lived initial free-expansion phase followed by an adiabatic expansion (the Sedov phase) once the mass of swept-up material exceeds the initial ejecta mass (Spitzer 1978). We will concern ourselves only with the later stages of SNR evolution after the Sedov phase has finished.
A few $`\times 10^4`$ yrs after the initial supernova, the temperature of the adiabatically expanding remnant falls to the point at which cooling is efficient in the outer regions: $`T(56)\times 10^5`$ K. At this point around half of the remnant mass forms a thin shell which expands into the ISM in a pressure-driven snowplough (McKee & Ostriker 1977; Lozinskaya 1992) . It is at the point at which the PDS phase begins that we will start to simulate the evolution of SNRs.
At the beginning of the PDS phase the thermal energy $`E_T`$ in the remnant is $`E_T0.36E_{51}`$. The temperature $`T_{\mathrm{PDS}}`$ of the interior gas is given by
$$T_{\mathrm{PDS}}=1.545\times 10^{10}\frac{E_{51}\mathrm{pc}^3}{n_0r_{\mathrm{PDS}}^3}K,$$
(1)
typically a few $`\times 10^6`$ K. This leads to a value of the interior pressure in very good agreement with that of Chevalier (1974). The thin shell (containing half the mass) is at around $`5\times 10^5`$ K.
Initially a supernova inputs energy into the ISM such that the thermal energy is $`E_T0.72E_{\mathrm{SN}}`$ and the remainder is in the form of kinetic energy where $`E_{\mathrm{SN}}`$ is the total supernova energy, where the usual value is $`E_{\mathrm{SN}}=10^{51}`$ ergs (Chevalier 1974). By the time a remnant enters the PDS phase, around 50 per cent of the thermal energy will have been lost (Lozinskaya 1992). The shell temperature is $`(56)\times 10^5`$ K and the interior pressure will have fallen to $`PE_{\mathrm{SN}}/(4\pi R_{\mathrm{PDS}}^3)`$ where $`r_{\mathrm{PDS}}`$ is the radius at which the transfer to the PDS phase occurs (McKee & Ostriker 1977).
The radius $`r_{\mathrm{PDS}}`$ at which the PDS phase starts is given by Cioffi et al. (1988) as
$$r_{\mathrm{PDS}}=14E_{\mathrm{SN}}^{0.29}n_0^{0.43}\zeta _m^{0.143}\mathrm{pc}$$
(2)
where $`n_0`$ is the ambient density in Hydrogen atoms cm<sup>-3</sup> and $`\zeta _m`$ is related to the metallicity (for Solar metallicity $`\zeta _m=1`$).
The initial velocity $`v_{\mathrm{PDS}}`$ of the shell is given by
$$v_{\mathrm{PDS}}=413n_0^{0.143}E_{51}^{0.071}\zeta _m^{0.214}\mathrm{km}\mathrm{s}^1$$
(3)
These parameters (radius and velocity) are similar to those given by other authors (eg. Chevalier 1974; Falle 1981; Blondin et al. 1998). Differences are mainly due to the use of different cooling functions although the shell velocity predicted at the same radius is very similar for all calculations.
The main reasons for preferring the calculations of Cioffi et al. (1988) over other calculations are that they include a correction factor for non-Solar metallicities and their good agreement with the simulations of Thornton et al. (1998).
### 2.3 SNRs in the code
In the remainder of this section we compare our results for a SNR in a uniform medium with those of other authors as well as testing the convergence and self-consistency of our results.
SNRs are set-up by taking all of the particles in a sphere of radius $`r_{\mathrm{PDS}}`$ centred on the position of the supernova. Half of the particles are unmoved but are heated up to $`T_{\mathrm{PDS}}`$ to represent the hot interior (uniform density and pressure, no bulk motions). The other half of the particles are distributed on the surface of a sphere of radius $`r_{\mathrm{PDS}}`$ and given a velocity outwards from the centre of the sphere of $`v_{\mathrm{PDS}}`$ at a temperature of $`5\times 10^5`$ K.
Low particle numbers within $`r_{\mathrm{PDS}}`$ have an effect on the evolution of the SNR due to shot noise causing significant non-sphericity. Generally more than $`60`$ particles are required within $`r_{\mathrm{PDS}}`$ to model the SNR well although as few as 20 do a reasonable job (we would not wish however to trust simulations where the included particle number is this low).
#### 2.3.1 SNR energy
Table 1 gives the parameters used in 10 example runs covering a factor of 8 in $`N`$, 2.5 in box size and over 30 in mass resolution for the same physical initial conditions. These runs are used in the next 2 subsections to illustrate the stability of the code at reproducing results over this range of parameters and to compare with other authors to show the ability of the code to well represent the post-Sedov phase evolution of a SNR.
Figure 1 shows the change in the total energy within the simulation volume, normalised to the value at 1 Myr to take into account the different amounts of total thermal energy present in simulation volumes of varying size (in large boxes the thermal energy of the undisturbed gas is by far the largest contributor to the total energy). The ‘lost’ energy has been radiated away by the hot gas, the energy conservation being good, a maximum of $`0.10.2`$ per cent. Good agreement is seen over a wide variety of particle mass and box size. The small differences are mainly due to the requirement that there be an integer number of particles in the interior and shell of the SNR. This leads to small differences between the initial thermal and kinetic energies of each model which amount to a few per cent, a difference which is amplified as time progresses. For instance, the initial kinetic energy of runs 1019 and 1020 differ by 4 per cent, by 6 Myr the difference in fig. 1 of $`0.02E_{51}`$ represents approximately 4 per cent of the total energy change.
A more detailed examination of the energy balance for three of the SNRs is shown in fig. 2, approximately 0.6 Myr after the supernova explosion. The solid lines show the sum of the thermal and kinetic energy within each radius with the dashed and dot-dashed lines showing the total thermal and kinetic energy respectively within each radius. The outwardly moving shell is clearly visible. The differences in interior thermal energy are due to shot noise and very low particle numbers in these regions - the high thermal energies of runs 1002 and 1017 within the inner 10 pc are due to 1 hot particle. The agreement of the energies over a factor of more than 30 in mass resolution is good and is representative of all runs. The shell has a peak density at approximately 60 pc at 0.6 Myr and is wider where the mass resolution is poor.
#### 2.3.2 SNR expansion law
The expansion of a SNR can be roughly characterised as a power law of the form $`rt^\nu `$ (with $`t`$ being measured from the time of the supernova). McKee & Ostriker (1977) find for negligible external pressure that $`\nu =2/7(0.286)`$. Chevalier (1974) finds $`\nu =0.305`$ for early times and for late evolution ($`t>0.75`$ Myr) $`\nu =0.32`$. This is close to the Cioffi et al. (1988) offset solution where $`r(tt_{\mathrm{offset}})^{0.3}`$. From inspection of fig. 3 in Thornton et al. (1998), they appear to obtain a value around $`\nu =0.3`$. All these values are below the power law for a purely momentum conserving snowplough where $`\nu =0.4`$ (Spitzer 1978).
Figure 3 shows the evolution of the shell radius with time for the ten convergence test runs (which model the same physical conditions at different mass resolutions and box sizes) detailed in table 1. In fig. 3 the shell radius is determined by the mean radial distance from the centre of the SNR of the densest particles. At low shell radii $`r`$ this is less reliable due to the smaller number of particles in the shock (especially when the mass resolution is poor).
The best fit to the first 1.5 Myr of evolution is shown by the dashed curve fitted with;
$$r(t)=r_{\mathrm{PDS}}\left(\frac{t}{t_{\mathrm{PDS}}}+1\right)^\nu $$
(4)
with $`r`$ in pc and $`t`$ in Myr ($`t`$ is the simulation time, ie. the time since the onset of the PDS phase), $`t_{\mathrm{PDS}}`$ is of order the duration of the PDS phase. Leaving $`r_{\mathrm{PDS}}`$, $`t_{\mathrm{PDS}}`$ and $`\nu `$ as free parameters the best fit is given by $`r_{\mathrm{PDS}}=18.34`$ pc, $`t_{\mathrm{PDS}}=0.012`$ Myr and $`\nu =0.317`$. For an $`n_0=0.5`$ cm<sup>-3</sup> ISM, $`r_{\mathrm{PDS}}`$ is actually 18.85 pc (from eqn. 2) and $`t_{\mathrm{PDS}}=r_{\mathrm{PDS}}/v_{\mathrm{PDS}}`$ is around 0.02 Myr (Cioffi et al. 1988).
The late stages ($`>1.5`$ Myr) are well fitted by the solid line, $`r=55.8+16.8t`$, a reasonably good fit to free expansion of a sound wave, the sound speed in this case of 15.6 pc Myr<sup>-1</sup>.
#### 2.3.3 Different ambient densities
We have investigated the differences in SNR evolution within ISMs of different ambient densities. Table 2 shows the range of ambient densities tested, covering 4 orders of magnitude. Previous studies (with the exception of Thornton et al. 1998) have concentrated on investigating low ambient densities representative of the ISM in the Solar neighbourhood.
Figure 4 shows the shell evolution for the range of densities presented in table 2. Also shown is the radius and time at which the shell velocity falls to the sound speed for each density.
From the data presented in table 2 the radius at which the sound speed is reached is fitted by a power-law in density of $`r_{\mathrm{sound}}=65n_0^{0.37}`$ pc and the time by a power-law of the form $`t_{\mathrm{sound}}=1.2n_0^{0.38}`$ Myr.
These power-law fits to $`r_{\mathrm{sound}}`$ and $`t_{\mathrm{sound}}`$ are almost exactly what are expected. Using eqn. 4 and the known dependencies on $`n_0`$ of $`r_{\mathrm{PDS}}(n_0^{0.43}`$ \- eqn. 2) and $`t_{\mathrm{PDS}}(r_{\mathrm{PDS}}/v_{\mathrm{PDS}}n_0^{0.57}`$ -eqns. 23) gives for $`t/t_{\mathrm{PDS}}>>1`$
$$rn_0^{0.25}t^{0.32}$$
(5)
as $`v_{\mathrm{sound}}=r_{\mathrm{sound}}/t_{\mathrm{sound}}=constant`$ leads to the predictions that $`t_{\mathrm{sound}}n_0^{0.36}`$ and $`r_{\mathrm{sound}}n_0^{0.37}`$, very close to the fitted relations from the simulations.
The initial energy (thermal and kinetic) of the SNR and the evolution of the total energy is approximately the same for each density if the evolution of a SNR through the PDS phase is scaled so that the radial properties are in units of $`r_{\mathrm{PDS}}`$ as illustrated in fig. 5. The radial total, thermal and kinetic energies are plotted against $`r/r_{\mathrm{PDS}}`$ for 3 models spanning a range of 4 orders of magnitude in density for the time when $`r_{\mathrm{shell}}2\times r_{\mathrm{PDS}}`$. For $`n_0=0.01`$ cm<sup>-3</sup>, $`2r_{\mathrm{PDS}}200`$ pc, for $`n_0=0.5`$ cm<sup>-3</sup>, $`2r_{\mathrm{PDS}}38`$ pc while for $`n_0=100`$ cm<sup>-3</sup>, $`2r_{\mathrm{PDS}}3.9`$ pc.
As is clear in fig. 5 the energy profiles of the SNRs in terms of $`r/r_{\mathrm{PDS}}`$ are very similar, indeed the main differences are due to the fact that the output radii shown are not exactly $`r_{\mathrm{shell}}=2\times r_{\mathrm{PDS}}`$ but differ by a few per cent (due to the finite number of outputs).
When $`r=2\times r_{\mathrm{PDS}}`$ the mass of the hot interior is around twice the interior mass at the onset of the PDS phase as most of the swept-up material remains in the dense shell. The interior temperature falls by a factor of around 2.7 while the shell temperature has fallen to $`10^4`$ K. The shell velocity has fallen by a factor of 5 (as expected if $`v`$ scales as $`v_{\mathrm{PDS}}(t/t_{\mathrm{PDS}}+1)^{0.68}`$, cf. eqn. 4). The scaling of the SNR properties with $`r_{\mathrm{PDS}}`$ across all densities is because the energy content of SNRs is density independent - a doubling of the volume of the SNR will result in the same energy changes throughout the SNR independent of the ambient density (as swept-up mass dominates).
#### 2.3.4 Late-starting SNRs
Figure 6 shows the late evolution of three SNRs, two started at the onset of the PDS phase and the other started at $`r=2\times r_{\mathrm{PDS}}`$ with $`T_{\mathrm{shell}}=10^4`$ K and the interior temperature given by $`T_{\mathrm{SNR}}=0.37\times T_{\mathrm{PDS}}`$ (where $`T_{\mathrm{PDS}}`$ is given by eqn. 1) and $`v_{\mathrm{shell}}=0.2v_{\mathrm{PDS}}`$. The total energy within the late-starting SNR is very similar to that in those that started at $`r_{\mathrm{PDS}}`$ but the relative contributions of the kinetic and thermal energies within the bubble are slightly different. This is due to the shorter amount of time which the $`2r_{\mathrm{PDS}}`$ simulation has had to convert bulk motions into thermal energy within the bubble. Again the main differences between the simulations are due to outputs not being at exactly the right times. The ability to start modelling each SNR at $`2\times r_{\mathrm{PDS}}`$ will become very useful when we wish to add SNRs to massive gas clouds reducing the required mass resolutions in simulations by a factor of 8 (see section 3).
#### 2.3.5 Metallicity effects
The effect of different metallicity environments can be significant during the early evolution of SNRs, during the PDS phase however the differences due to metallicity are found to be negligible (see also Thornton et al. 1998). Initially a metallicity of $`\zeta _m=0.01`$ will make a factor of 1.93 difference in $`r_{\mathrm{PDS}}`$ and 0.37 in $`v_{\mathrm{PDS}}`$ as compared to a $`\zeta _m=1`$ SNR (from eqns. 2 & 5 respectively). However, the effect of these different initial SNR conditions is not as significant as it might first appear, because the velocity of a $`\zeta _m=1`$ SNR at $`r=1.93\times r_{\mathrm{PDS}}(\zeta _m=1)`$ is $`0.22v_{\mathrm{PDS}}(\zeta _m=1)`$ (see above) resulting in only slightly different expansions.
The subsequent evolution of SNRs with different metallicities through the PDS phase are virtually indistinguishable (as also found by Thornton et al. 1998). The evolution of the shell radii is slightly different at early times due to different values of $`r_{\mathrm{PDS}}`$ but as mentioned above this difference is soon nearly cancelled out and at late times the shell radii differ by only a few pc. When SNRs are started with the same $`r_{\mathrm{PDS}}`$ but different metallicities the difference is negligible. The evolution of the total energies within each SNR are also remarkably similar. Differences in metallicity affect the cooling of the hot interior and hence the pressure and the extra driving force on the shell. However, during expansion into a uniform ambient medium the shock acceleration is dominated by the deceleration due to the sweeping-up of material, making metallicity (which effects the pressure-driven acceleration) a minor factor in the late-time evolution of the SNR.
The differences in energy between different metallicities are accounted for by the different initial conditions at the onset of the PDS phase however it is not clear that even the early stages of SNR evolution would be affected by low-metallicity effects. A type II supernova with a progenitor mass of $`25M_{}`$ will produce around $`2.4M_{}`$ of heavy metals (Tsujimoto et al. 1995), the majority ($`1.8M_{}`$) being O<sup>16</sup>. In the case of a SNR expanding into an ISM of density $`n_0=100`$ cm<sup>-3</sup> (where the mass within $`r_{\mathrm{PDS}}`$ is $`75.5M_{}`$) this will enrich the swept-up material from $`\zeta _m=1`$ (Solar) to $`\zeta _m=2.6`$ and from $`\zeta _m=0.01`$ to $`\zeta _m=1.6`$. Such enrichment is very significant and will alter the evolution of the pre-PDS phases. However, by $`5r_{\mathrm{PDS}}`$ the amount of swept-up material is $`9440M_{}`$ and has been enriched from $`\zeta _m=0.01`$ to $`\zeta _m=0.023`$, a fairly insignificant (from the point of view of cooling) amount. Thornton et al. (1998) did not find a significant effect when including the metals ejected from the supernova (Thornton 1999). Obviously this enrichment and its history is of vital importance for understanding chemical evolution, and the enrichment history can be followed by this code (once a number of assumptions about mixing have been included), but it has little bearing on the late-time evolution of an individual SNR in a uniform medium.
#### 2.3.6 Computational aspects
The softening included in the simulations does not greatly affect the results of the shell expansion and energy transfer. Similar runs differing in softening are virtually indistinguishable in their results (two such runs, differing by a factor of over 4 in softening - 1002 & 1005 - are included in fig. 1). The softening selected for simulations should be as small as is reasonably practicable given that the smallest region in which results can be believed in detail is no less than the softening length, which also determines the width of any shock fronts that are present. Unfortunately the number of timesteps required rises as the softening length is decreased.
The width of the shell that forms the snowplough is difficult to determine accurately as it becomes slightly aspherical as it evolves due to initial Poisson fluctuations in the particle positions. This effect is worse for the models which began with only a few particles in the SNR. When spherically averaged profiles are produced this leads to an artificial broadening of the snowplough and general smoothing of the steep shock fronts. The snowplough width is expected to be a few pc (Thornton et al. 1998), in these simulations the shock width is approximately the softening.
### 2.4 Summary
In this section we have shown that the code is stable and converges for reasonable selections of box size and particle mass. The code is able to reproduce the results of other authors for the situation of a SNR expanding into a uniform ambient medium during the PDS phase of evolution. In the next section we will place SNRs in the centres of Plummer model gas clouds.
## 3 CENTRAL SUPERNOVAE IN A GAS CLOUD
While the evolution of a SNR in a uniform density ISM provides a good test of the ability of the code to reproduce the characteristics of a SNR as found in more detailed hydrodynamic calculations, it is not physically realistic. The effects of a supernovae in a stratified medium (eg. a gas cloud) is a more practical problem. Massive stars are found in young clusters, their lifetimes being so short that they are expected to still be embedded in gas remaining from star formation. The significant evolution of a young SNR is therefore expected to occur in the confines of its parent GMC. Massive stars are also found to be centrally-concentrated in clusters (Hillenbrand 1997; Carpenter et al. 1997) so that an investigation of a central SNR is a good approximation to the evolution of the first SNR in a cluster.
Some analytic calculations of the effects of supernovae in a gas cloud have been made by Dopita & Smith (1986) and Morgan & Lake (1989) to estimate the number of supernovae required to totally disrupt a gas cloud of a given mass. Morgan & Lake (1989) using more detailed cooling functions than Dopita & Smith (1986) found that the minimum mass of a ($`1/r^2`$) cloud for it to confine a single central supernova was $`4\times 10^4M_{}`$.
As noted in section 2.3, Dopita & Smith (1986) find that a SNR within the high-density environment of a gas cloud will remain roughly spherical during the adiabatic phase of its evolution. Using this result we may place SNRs in the cloud at the start of the PDS phase in the same way as they are placed in the uniform medium simulations. The high densities in gas clouds ($`n_0>10^4`$ cm<sup>-3</sup>) means that the adiabatic phase will end whilst the shell radius is very small ($``$ pc).
Cloud initial conditions are based on a Plummer model ($`n=5`$ polytrope) with mass distribution;
$$M(r)=\frac{Mr^3}{R^3}\frac{1}{[1+r^2/R^2]^{3/2}}$$
(6)
where $`M`$ is the total mass and $`R`$ is the scale length. The half-mass radius of a cloud is then $`1.3R`$. Clouds are constructed out to a maximum radius $`r_{\mathrm{max}}`$ of a few $`R`$ (normally $`r_{\mathrm{max}}=20`$ pc), resulting in the actual mass being slightly less than $`M`$.
In this paper we will only deal with Plummer model clouds. Obviously the effects of a SNR will depend, at least to some extent and possibly very significantly, upon the density distribution of the parent cloud. These will be dealt with in detail in a paper to follow.
Observationally, Giant Molecular Clouds (GMCs) can approximated very roughly by clouds with $`M=`$ a few $`\times 10^4`$ to a few $`\times 10^5M_{}`$ and $`R=`$ a few pc (eg. Harris & Pudritz 1994) which we used as initial conditions for our simulations (full details in table 3). Before inserting any supernovae, our clouds are allowed to relax to a stable state in which they are virialised. Support is provided by the thermal energy of the gas and the bulk kinetic energy is negligible (although the temperature of the gas represents the turbulent velocity which supports the cloud rather than the molecular temperature). The clouds are self-gravitating as are real giant molecular clouds (Rivolo & Solomon 1988) and so require no pressure confinement from an external hot diffuse gas.
The insertion of a supernova has the effect of adding considerable amounts of kinetic and thermal energy to the cloud. By the end of the Sedov phase, the extra energy amounts to around $`6.4\times 10^{50}`$ ergs, comparable to or greater than the potential energy of the system (roughly $`8.6\times 10^{40}(M/M_{})^2/(R/\mathrm{pc})`$ ergs). The shock from the supernova then passes through the the cloud.
### 3.1 Some numerical considerations
The number of particles within $`r_{\mathrm{PDS}}`$ required for the SNR evolution to be a good approximation to more complex simulations is $`>40`$ with $`>60`$ giving the best results (see section 2). For a Plummer model cloud the central density is
$$\rho _0=\frac{3M}{4\pi R^3}$$
(7)
where $`M`$ is the mass of the cloud and $`R`$ is the Plummer scale length. Using equation 2 the mass interior to $`r_{\mathrm{PDS}}`$ at the cloud core will be
$$M_{\mathrm{PDS}}5.48\frac{R_{\mathrm{pc}}^{6/7}}{M_5^{2/7}}M_{}$$
(8)
where $`R_{\mathrm{pc}}=R/`$pc and $`M_5=M/10^5M_{}`$ (we shall be using $`R_{\mathrm{pc}}`$ and $`M_5`$ throughout the rest of this section). We require a mass resolution at least 60 times smaller than this giving, for $`R_{\mathrm{pc}}=3.5`$ and $`M_5=1`$ a minimum particle number of $`N>3.74\times 10^5`$.
Such resolutions, while attainable, require significant computing time to run. Our aim here is to introduce SNRs in such a way as to allow many simulations covering a variety of initial conditions to be completed, so as to explore the parameter space. Starting SNR evolution at $`r=2\times r_{\mathrm{PDS}}`$ (see Section 2.3.4) allows a factor of 8 reduction in the mass resolution and therefore the number of particles required to model the cloud. In the case of a $`10^5M_{}`$ cloud $`2\times 32^3`$ particles are required, on a workstation each timestep takes around 3 minutes resulting in a 20 Myr simulation in around 30 hours.
### 3.2 Feedback as a function of cloud structure
Here we present the analysis of the energy and mass feedback as a function of $`M_5`$ and $`R_{\mathrm{pc}}`$, the cloud models are presented in table 3. Initially all clouds are virialised and particles are said to be lost once they cross the box edges which are a distance $`2r_{\mathrm{max}}`$ from the centre of the cloud (at a distance roughly equal to the tidal radius of a GMC). It is these lost particles that comprise the feedback from the cloud into the larger ISM. In most cases the SNRs are started at $`2r_{\mathrm{PDS}}`$.
As a ’standard’ model we use the run with parameters $`M_5=1`$, $`R_{\mathrm{pc}}=3.5`$ to illustrate the general features of SNR evolution in a Plummer cloud. Initially the cloud has $`0.0356\times 10^{51}`$ ergs of thermal energy and a potential energy of $`0.0711\times 10^{51}`$ ergs. The SNR starting at $`2r_{\mathrm{PDS}}`$ adds a total of $`0.123\times 10^{51}`$ ergs of kinetic energy and $`0.289\times 10^{51}`$ ergs of thermal energy (calculated from section 2) creating a net positive energy for the system of $`0.38\times 10^{51}`$ ergs.
#### 3.2.1 The evolution of the SNR
The evolution of the shell velocity with radius for the standard model is shown in fig. 7. It is typical of all runs that the velocity drops from its initial value ($`v_{\mathrm{PDS}}`$ or $`0.2v_{\mathrm{PDS}}`$ of started at $`2r_{\mathrm{PDS}}`$) reaching a minimum value $`v_{\mathrm{min}}`$ at around 2 to 3 $`R_{\mathrm{pc}}`$ before accelerating to around $`1.5\times v_{\mathrm{min}}`$ on leaving the box.
It is found later that $`v_{\mathrm{min}}`$ and associated quantities may be used to parameterise all of the effects of a SNR on a gas cloud. In this subsection we sketch a model of the SNR evolution that predicts $`v_{\mathrm{min}}`$ in terms of the two cloud parameters $`M_5`$ and $`R_{\mathrm{pc}}`$.
The velocity reaches $`v_{\mathrm{min}}`$ when the acceleration due to the interior pressure matches the deceleration due to the sweeping-up of material. The force on the shell due to pressure $`F_{\mathrm{press}}(r)`$ goes as
$$F_{\mathrm{press}}=(4\pi r^2)P$$
(9)
where $`P`$ is the pressure when the shell radius is $`r(t)`$.
The initial radius of the shell ($`r_{\mathrm{PDS}}`$) is very small compared to $`R`$ and as adiabatic cooling goes as $`(r_{\mathrm{PDS}}/r)^5`$ the majority of the cooling will be adiabatic, especially at late times so
$$PP_{\mathrm{PDS}}\frac{r_{\mathrm{PDS}}^5}{r^5}$$
(10)
with $`P_0`$ given by $`\rho _0kT_{\mathrm{PDS}}`$ (where $`\rho _0`$ is the central density $`M_5/R_{\mathrm{pc}}^3`$). The initial temperature $`T_{\mathrm{PDS}}`$ is given by eqn. 1 as
$$T_{\mathrm{PDS}}\frac{M_5^{0.29}}{R_{\mathrm{pc}}^{0.87}}$$
(11)
The force due to the sweeping-up of material (momentum conservation) $`F_{\mathrm{mom}}(r)`$ is
$$F_{\mathrm{mom}}(r)=(4\pi r^2)\rho (r)v(r)^2$$
(12)
equating eqns. 9 and 12 using eqn. 10 gives
$$\rho (r_{\mathrm{min}})v_{\mathrm{min}}^2=P_0\frac{r_{\mathrm{PDS}}^5}{r_{\mathrm{min}}^5}$$
(13)
using eqn. 11 and assuming that $`r_{\mathrm{min}}R`$ gives
$$v_{\mathrm{min}}\frac{R_{\mathrm{pc}}^{0.29}}{M_5^{0.93}}$$
(14)
Fitting the powers of $`M_5`$ and $`R_{\mathrm{pc}}`$ simultaneously gives a best linear fit to $`v_{\mathrm{min}}`$ of
$$v_{\mathrm{min}}=6.58\frac{R_{\mathrm{pc}}^{0.32}}{M_5^{0.77}}\mathrm{km}\mathrm{s}^1$$
(15)
which is compared to the simulations in fig. 9. The similarities between eqns. 14 and 15 are remarkable considering the very simple assumptions that went into the formulation of eqn. 14.
The accelerating force due to the pressure goes as $`r^5`$ and so rapidly becomes negligible at high $`r`$. It is still able, however, to increase the velocity of the shell by the time it leaves the cloud to $`1.5v_{\mathrm{min}}`$.
The above rather simplistic treatment avoids (as is necessary) many details of the actual SNR evolution in a cloud. In the simulations the main shell of the SNR is followed by a weaker shock which is created by the shocking of infalling material filling the interior when it reaches the centre (and meets other infalling material). This creates a very complex density/pressure structure in the cloud as the SNR evolves.
We know go on to show how the effects of a SNR on a cloud can be parameterised very simply in terms of $`v_{\mathrm{min}}`$ and closely related quantities.
#### 3.2.2 Metallicity effects
As the temperature factor $`T_0`$ is so important in setting the value of $`v_{\mathrm{min}}`$, then metallicity, $`Z`$, is an important factor in that in clouds of lower $`Z`$ than the fiducial Solar $`Z=1`$, cooling will be less efficient and $`T_0`$ and $`v_{\mathrm{min}}`$ higher. In the case of a $`Z=0.01`$ metallicity cloud the kinetic energy lost increases to $`2.5\times 10^{49}`$ ergs and the mass lost to $`5.4\times 10^4M_{}`$, a 67 per cent increase on the $`Z=1`$ cloud in both cases. For a $`Z=0.1`$ cloud the increase in both values is 46 per cent. This would indicate that the effects of the first generation of supernovae would have been much more dramatic than those which occurred later once metal enrichment had taken place.
#### 3.2.3 Mass loss and disruption
The mass loss from the cloud would be expected to be related the escape velocity from the cloud at the point where the minimum velocity is reached. The escape velocity at $`2R_{\mathrm{pc}}`$ is given by $`18.6\sqrt{M_5/R_{\mathrm{pc}}}`$ km s<sup>-1</sup> and it is found that if $`v_{\mathrm{min}}>v_{\mathrm{esc}}(r_{\mathrm{min}})`$ then the cloud completely disrupts, ie. all of the mass is lost. At the other extreme if $`v_{\mathrm{min}}<0.2v_{\mathrm{esc}}(r_{\mathrm{min}})`$ then no mass loss (and hence no feedback) occurs. In the intermediate regime the mass loss is related to $`v_{\mathrm{esc}}(r_{\mathrm{min}})`$ by
$$\frac{M_{\mathrm{lost}}}{M_5}=0.25\times \left(\frac{v_{\mathrm{min}}}{v_{\mathrm{esc}}}\right)^{2.5}$$
(16)
as illustrated in fig. 10.
#### 3.2.4 Energy feedback
The amount of energy returned to the ISM from a cloud would be expected to be related to the kinetic energy of the cloud at $`v_{\mathrm{min}}`$ given by $`T_{\mathrm{min}}=1/2M_{\mathrm{shell}}v_{\mathrm{min}}^2`$ where $`M_{\mathrm{shell}}`$ is the mass of the shell at $`r_{\mathrm{min}}`$, most simply we might expect $`M_{\mathrm{shell}}M_5`$.
Illustrated in fig. 11 is the feedback energy from table 3 against $`T_{\mathrm{min}}/10^{51}`$ ergs $`=0.001M_5(v_{\mathrm{min}}/`$km s$`{}_{}{}^{1})^2`$. Two linear regimes are obvious in the behaviour of the feedback energy. In the upper limit where the cloud has totally disrupted the feedback energy is fitted by
$$E_{\mathrm{lost}}=0.0012+0.25M_5v_{\mathrm{min}}^2\times 10^{51}\mathrm{ergs}$$
(17)
in the lower limit where the mass loss is negligible the feedback energy is
$$E_{\mathrm{elost}}=0.012+0.25M_5v_{\mathrm{min}}^2\times 10^{51}\mathrm{ergs}$$
(18)
as can be easily seen the slope of these two relationships is the same and the constant factor, different by exactly one order of magnitude, is the only difference. The intermediate stage marks the change from the total cloud destruction and high feedback to the low mass loss and low feedback regime. This region is only a very small region of the total $`M_5R_{\mathrm{pc}}`$ parameter space.
#### 3.2.5 The final state of clouds
As stated above when $`v_{\mathrm{min}}>v_{\mathrm{esc}}`$ a single SNR is enough to disrupt a cloud completely and when $`v_{\mathrm{min}}<0.2v_{\mathrm{esc}}`$ no mass loss occurs at all. In the intermediate regime some gas is retained in a bound object - a new cloud with lower mass than the original.
The ratio of the final to initial potential energy, $`\mathrm{\Omega }_f/\mathrm{\Omega }_i`$ is related again to the ratio of the minimum velocity of the shell and the escape velocity $`v_{\mathrm{min}}/v_{\mathrm{esc}}`$
$$\frac{\mathrm{\Omega }_f}{\mathrm{\Omega }_i}=1.261.05\frac{v_{\mathrm{min}}}{v_{\mathrm{esc}}}$$
(19)
as illustrated in fig. 12. This relationship is only valid for $`0.3<v_{\mathrm{min}}/v_{\mathrm{esc}}<1`$, beyond 1 the cloud is totally destroyed and $`\mathrm{\Omega }_f/\mathrm{\Omega }_i=0`$ and less than 0.3 $`\mathrm{\Omega }_f/\mathrm{\Omega }_i`$ asymptotes to 0.
The change in potential energy of a gas cloud due to a central supernova can sometimes be significantly more important than the feedback of kinetic energy as illustrated in table 3. Clouds lose all of their initial potential energy if they are completely disrupted. If the cloud loses only part of its mass then the significant energy change occurs as a loss of potential energy from the cloud (a net gain of energy). Even when no mass loss or feedback occurs a small change (of the order of a few percent) in the potential energy is observed as the cloud expands slightly.
As fig. 13 shows, the remaining material settles back into a new configuration that can also be described as a Plummer model with a lower central density and larger scale length. Thus knowing the initial mass $`M`$ and characteristic radius $`R`$ gives the final characteristic radius $`R_f`$ via;
$$R_f=\left(1\frac{M_{\mathrm{lost}}}{M}\right)^2R\left(\frac{\mathrm{\Omega }_f}{\mathrm{\Omega }_i}\right)^1$$
(20)
Knowing the final mass and characteristic radius of the cloud allows the effects of further central supernovae on that cloud to be calculated - as long as the interval between supernovae is greater than the time required for the cloud to relax to a new equilibrium.
The time taken for a $`5\times 10^5M_{}`$ cloud to recover from a central supernova is very short, the majority of the SNR’s kinetic and thermal energy is radiated away in $`<1`$Myr. Assuming a Salpeter IMF, the number of massive stars within a cloud that will go supernovae is $`N_{\mathrm{SNe}}0.006Mϵ`$ where $`M`$ is the mass of the cloud and $`ϵ`$ is the star formation efficiency. Taking $`ϵ=0.01`$ and $`M=5\times 10^5M_{}`$ gives $`N_{\mathrm{SNe}}=30`$. If these supernovae occur centrally and evenly spaced over a 30 Myr period then no ejecta will escape the cloud and the cloud will not be disrupted by the SNRs.
These results hold for Plummer model clouds that can be characterised solely by $`M`$ and $`R`$. This study will be extended to other density distributions in a future paper which we are preparing.
## 4 CONCLUSIONS
The Hydra $`N`$-body SPH code has been extended to allow the simulation of the evolution of supernova remnants (SNRs) from the onset of the pressure-driven snowplough (PDS) phase. In section 2 this code was seen to be able to produce convergent results over a wide range of parameter space and reproduce the results on the evolution of SNRs from a variety of previous authors. The power of this code is the ability to simulate the evolution of SNRs in a variety of environments using a workstation in a reasonable time (of the order of days).
This code represents the first time that we are able to model multiple, interacting SNRs in gas clouds and complexes. This paper is the first detailed analysis of SNR evolution in a gas cloud to examine the feedback parameters so important in galaxy formation and evolution calculations.
In section 3 we presented new results on the effect of a single, central supernova on Plummer model gas clouds of various masses and characteristic radii. The results of Dopita & Smith (1986) and Morgan & Lake (1989) were found to be too simplistic. The evolution of a central SNR is very complex, the late stages of evolution are governed by pressure-driving from a hot interior accelerating the SNR shell down the density gradient.
Feedback from a cloud is defined to be the mass and energy which pass out of our simulation box (whose size is approximately the tidal radius of the cloud). The results may be summarised as:
* The efficiency of energy feedback, mass loss and cloud destruction for a central supernova in a Plummer cloud of mass $`M`$ and characteristic radius $`R`$ is related to the minimum velocity $`v_{\mathrm{min}}`$ that the shock reaches during it’s evolution. This minimum velocity is given by $`v_{\mathrm{min}}=R_{\mathrm{pc}}^{0.32}/M_5^{0.77}`$.
* For $`v_{\mathrm{min}}>v_{\mathrm{esc}}`$ the cloud is totally destroyed, while for $`v_{\mathrm{min}}<0.2v_{\mathrm{esc}}`$ the SNR is completely contained and no feedback occurs.
* The mass lost from a cloud is related to the ratio of the minimum velocity to the escape velocity $`v_{\mathrm{min}}/v_{\mathrm{esc}}`$ as $`M_{\mathrm{lost}}/M=0.25(v_{\mathrm{min}}/v_{\mathrm{esc}})^{2.5}`$.
* The energy feedback $`E_{\mathrm{lost}}`$ has two main regimes with equal slopes where $`E_{\mathrm{lost}}=C_E+0.25M_5V_{\mathrm{min}}^2\times 10^{51}`$ ergs where $`C_E`$ = 0.0012 when $`v_{\mathrm{min}}/v_{\mathrm{esc}}>1.1`$ and $`C_E`$ = 0.012 when $`v_{\mathrm{min}}/v_{\mathrm{esc}}<0.9`$.
* The loss of (negative) potential energy is often the largest (positive) increase of energy in the system and at least of order the feedback of kinetic energy, the feedback of thermal energy being negligible in comparison to both.
* The final state of a cloud that is not destroyed is close to a Plummer model with final characteristic radius $`R_f`$ related to the initial parameters by $`R_f=(1M_{\mathrm{lost}}/M)^2R(\mathrm{\Omega }_f/\mathrm{\Omega }_i)^1`$.
* The efficiency of feedback increases rapidly with decreasing metallicity suggesting that feedback at early epochs was far more efficient.
Our simulations are obviously not perfect. The code takes no account of magnetic fields (following standard astrophysical practice) which may well be important, especially at late times. In addition the feedback from massive stars into the cloud before they become supernovae is neglected, although in the dense environments of cloud cores we may be justified in ignoring this effect (Franco, García-Segura & Plewa 1996). A simple method of including such effects in the code is being developed. Despite this we believe that this code represents a significant step in modelling the effects of multiple SNRs on clouds and the larger ISM.
### 4.1 Future directions
As noted in the introduction, a code such as this has many interesting applications. It can be used to investigate how gas clouds of different sizes and shapes are affected by internal supernovae, calculating the feedback of energy and mass into the ICM. We are in the process of preparing papers that will expand the current work into an investigation of off-centre SNRs in gas clouds and of the effect of different density distributions on feedback parameters.
In addition we will be able to investigate how the cloud is disrupted and in what way gas is expelled on a small scale. This has important consequences for investigating the early evolution of star clusters. A paper is in preparation on the effects of gas expulsion on the dynamics of the stellar content of a cluster and what star formation efficiency is required for a bound object to remain.
Using supercomputers $`N=10^7`$ or more is possible, giving the resolution to model fully the dynamical and chemical evolution of a dwarf galaxy. In addition feedback can be placed into simulations of the formation of the first cosmological objects.
## ACKNOWLEDGMENTS
PAT is a PPARC Lecturer Fellow. We would like to thank Hugh Couchman for useful discussions about this paper and acknowledge NATO CRG 970081 which facilitated this interaction. We would also like to thank Anne Green and Roger Hutchings for their help. This project was completed using the computer facilities of the University of Sussex Astronomy Centre and University of Sussex BFG computer using the Hydra code (public release version available from the Hydra Consortium at http://phobos.astro.uwo.ca/hydra\_consort/).
In memory of Jenny. |
warning/0001/astro-ph0001354.html | ar5iv | text | # New Constraints from High Redshift Supernovae and Lensing Statistics upon Scalar Field Cosmologies
## I Introduction
Recent observations of type Ia supernovae (SNe Ia) at high redshift suggest that the expansion of the Universe is accelerating : these calibrated ‘standard’ candles appear fainter than would be expected if the expansion were slowing due to gravity. While concerns about systematic errors (such as possible evolution of the source population and grey dust) remain, the current evidence indicates that the high-redshift supernovae appear fainter because, at fixed redshift, they are at larger distances. According to the Friedmann equation, $`\ddot{a}/a=(4\pi G/3)(\rho +3p)`$, accelerated expansion requires a dominant component with either negative energy density, which is physically inadmissible, or effective negative pressure. Dark energy, dynamical-$`\mathrm{\Lambda }`$ (dynamical vacuum energy), or quintessence are different names that have been used to denote this component. A cosmological constant, with $`p_\mathrm{\Lambda }=\rho _\mathrm{\Lambda }`$, is the simplest possibility.
Recent studies incorporating new CMB data confirm previous analyses suggesting a large value for the total density parameter, $`\mathrm{\Omega }_{total}>0.4`$, and favor a nearly flat Universe ($`\mathrm{\Omega }_{total}=1`$). A different set of observations now unambiguously point to low values for the matter density parameter, $`\mathrm{\Omega }_{m0}=0.3\pm 0.1`$. In combination, these two results provide independent evidence for the conventional interpretation of the SNe Ia results and strongly support a spatially flat cosmology with $`\mathrm{\Omega }_{m0}0.3`$ and a dark energy component with $`\mathrm{\Omega }_X0.7`$. These models are also theoretically appealing since a dark energy component that is homogeneous on small scales (20–30 $`h^{1\text{ }}`$ Mpc) reconciles the spatial flatness predicted by inflation with the sub-critical value of $`\mathrm{\Omega }_{m0}`$ .
The cosmological constant has been introduced several times in modern cosmology to reconcile theory with observations and subsequently discarded when improved data or interpretation showed it was not needed. However, it may be that the “genie” will now remain forever out of the bottle . Although current cosmological observations favor a cosmological constant, there is as yet no explanation why its value is 50 to 120 orders of magnitude below the naive estimates of quantum field theory. One of the original motivations for introducing the idea of a dynamical $`\mathrm{\Lambda }`$-term was to alleviate this problem. There are also observational motivations for considering dynamical-$`\mathrm{\Lambda }`$ as opposed to constant-$`\mathrm{\Lambda }`$ models. For instance, the COBE-normalized amplitude of the mass power spectrum is in general lower in a dynamical-$`\mathrm{\Lambda }`$ model than in a constant-$`\mathrm{\Lambda }`$ one, in accordance with observations . Further, since distances are smaller (for fixed $`z`$ and $`\mathrm{\Omega }_{m0}`$), constraints from the statistics of lensed QSOs are weaker in dynamical-$`\mathrm{\Lambda }`$ models.
## II Scalar Field Models
A number of models with a dynamical $`\mathrm{\Lambda }`$ have been discussed in the literature . We report here new constraints from gravitational lensing statistics and high-z SNe Ia on two representative scalar field potentials that give rise to effective decaying $`\mathrm{\Lambda }`$ models: pseudo-Nambu-Goldstone bosons (PNGB), with potential of the form $`V(\varphi )=M^4\left(1+\mathrm{cos}\left(\varphi /f\right)\right)`$, and inverse power-law models, $`V(\varphi )=M^{4+\alpha }\varphi ^\alpha `$. These two models are chosen to be representative of the range of dynamical behavior of scalar field ‘quintessence’ models. In the PNGB model, the scalar field at early times is frozen and therefore acts as a cosmological constant; at late times, the field becomes dynamical, eventually oscillating about the potential minimum, and the large-scale equation of state approaches that of non-relativistic matter ($`p=0`$). The power-law model, on the other hand, exhibits “tracker” solutions : at high redshift, the scalar field equation of state is close to that of non-relativistic matter, and at late times it approaches that of the cosmological constant.
Let us consider first the motivation for the PNGB model. All “quintessence” models involve a scalar field with ultra-low effective mass. In quantum field theory, such ultra-low-mass scalars are not generically natural: radiative corrections generate large mass renormalizations at each order of perturbation theory. To incorporate ultra-light scalars into particle physics, their small masses should be at least ‘technically’ natural, that is, protected by symmetries, such that when the small masses are set to zero, they cannot be generated in any order of perturbation theory, owing to the restrictive symmetry. Pseudo-Nambu-Goldstone bosons (PNGBs) are the simplest way to have naturally ultra–low mass, spin–$`0`$ particles. These models are characterized by two mass scales, a spontaneous symmetry breaking scale $`f`$ (at which the effective Lagrangian still retains the symmetry) and an explicit breaking scale $`M`$ (at which the effective Lagrangian contains the explicit symmetry breaking term). In order to act approximately like a cosmological constant at recent epochs with $`\mathrm{\Omega }_\varphi 1`$, the potential energy density should be of order the critical density, $`M^43H_0^2m_{Pl}^2/8\pi `$, or $`M3\times 10^3h^{1/2}`$ eV. As usual we set $`V=0`$ at the minimum of the potential by the assumption that the fundamental vacuum energy of the Universe is zero – for reasons not yet understood. Further, since observations indicate an accelerated expansion, at present the field kinetic energy must be relatively small compared to its potential energy. This implies that the motion of the field is still (nearly) overdamped, that is, $`\sqrt{|V^{\prime \prime }(\varphi _0)|}3H_0=5\times 10^{33}h`$ eV, i.e., that the PNGB is ultra-light. The two conditions above imply that $`fm_{Pl}10^{19}`$ GeV. Note that $`M10^3`$ eV is close to the neutrino mass scale for the MSW solution to the solar neutrino problem, and $`fm_{Pl}10^{19}`$ GeV, the Planck scale. Since these scales have a plausible origin in particle physics models, we may have an explanation for the ‘coincidence’ that the vacuum energy is dynamically important at the present epoch . Moreover, the small mass $`m_\varphi M^2/f`$ is technically natural.
Next consider the inverse power-law model: this potential gives rise to attractor (tracking) solutions. If $`\rho _\varphi `$ and $`\rho _B`$ denote the mean scalar and dominant background (radiation or matter) densities, then if $`\rho _\varphi \rho _B`$, the following ‘tracker’ relationship is satisfied: $`\rho _\varphi ^{TR}a^{3(\gamma _B\gamma _\varphi ^{TR})}\rho _B`$, where $`\gamma _\varphi ^{TR}=\gamma _B\alpha /(2+\alpha )<\gamma _B`$ . Here, $`a(t)`$ is the cosmic scale factor, and $`\gamma _B=(p_B+\rho _B)/\rho _B`$ denotes the adiabatic index of the background ($`\gamma _B=4/3`$ during the radiation-dominated era and $`\gamma _B=1`$ during the matter-dominated epoch (MDE)). If the scalar field is in the tracker solution, its energy density decreases more slowly than the background energy density, and the field eventually begins to dominate the dynamics of the expansion. If the field is on track during the MDE, its effective adiabatic index is less than unity—its effective pressure $`p_\varphi =(\dot{\varphi }^2/2)V(\varphi )`$ is negative. This condition by itself does not guarantee accelerated expansion: the field must have sufficiently negative pressure and a sufficiently large energy density such that the total effective adiabatic index (of the field plus the matter) is less than 2/3. Moreoever, for inverse power-law potentials, at late times $`\mathrm{\Omega }_\varphi 1`$, such that when the growing $`\mathrm{\Omega }_\varphi `$ starts to become appreciable, $`\gamma _\varphi `$ deviates from the above tracking value, decreasing toward the value $`\gamma _\varphi 0`$. Thus, even if $`\alpha >4`$, such that initially $`\gamma _\varphi =\gamma _\varphi ^{TR}>2/3`$ in the MDE, when the field begins to dominate the energy density and $`\gamma _\varphi `$ decreases, the Universe will enter a phase of accelerated expansion. If $`\mathrm{\Omega }_{m0}`$ and $`\alpha `$ are sufficiently small, this will happen before the present time. For inverse power-law potentials, the two conditions $`\mathrm{\Omega }_{\varphi 0}1`$ and the preponderance of the field potential energy over its kinetic energy (the condition for negative pressure) imply $`M10^{\frac{27\alpha 12}{\alpha +4}}`$ eV and $`\varphi _0m_{Pl}`$. Since $`\varphi _0m_{Pl}`$, quantum gravitational corrections to the potential may be important and could invalidate this picture .
In the very early Universe, in order to successfully achieve tracking, the scalar field energy density must be smaller than the radiation energy density. If, in addition, $`\rho _\varphi `$ is smaller than the initial value of the tracking energy density, the field will remain frozen until they have comparable magnitude; at that point, the field starts to follow the tracking solution. On the other hand, if $`\rho _\varphi `$ is larger than the initial value of the tracking energy density, the field will enter a phase of kinetic energy domination ($`\gamma _\varphi 2`$); this causes $`\rho _\varphi `$ to decrease rapidly ($`\rho _\varphi a^6`$), overshooting the tracker solution . Subsequently, as in the case above, the field is frozen and later begins to follow the tracking solution when its energy density becomes comparable to the tracking energy density. In either case, there is always a phase before tracking during which the field is frozen. Consequently, an important variable is the value of the field energy density when it freezes. For instance, is it smaller or larger than $`\rho _{eq}`$, the mean energy density at the epoch of radiation-matter equality? Did the field have time to completely achieve tracking or not? In fact, the exact constraints imposed by cosmological tests on the parameter space of this model depend upon this condition.
In a previous study , we numerically evolved the scalar field equations of motion forward from the epoch of matter-radiation equality, assuming the field is initially frozen, $`\dot{\varphi }(t_{eq})=0`$. In this case, depending on the values of $`\alpha `$ and $`\mathrm{\Omega }_{m0}`$, it may happen that the field does not have time to reach the tracking solution before the present. In general, if $`\mathrm{\Omega }_{m0}`$ is large, we observe that at the present $`\gamma _\varphi `$ is still growing away from its initial value $`\gamma _\varphi =0`$. On the other hand, if $`\mathrm{\Omega }_{m0}`$ is sufficiently low, $`\gamma _\varphi `$ will reach a maximum value (not necessarily the tracking value) at some point in the past and at the present time will be decreasing to the value $`\gamma _\varphi 0`$. Here we follow a different approach. In our numerical computation we now start the evolution of the scalar field during the radiation dominated epoch and assume that it is on track early in the evolution of the Universe.<sup>*</sup><sup>*</sup>* In fact this is true only if $`\alpha `$ is not close to zero. The case $`\alpha =0`$ is equivalent to a cosmological constant, and the field remains frozen always. When $`\rho _\varphi `$ becomes non-negligible compared to the matter density, $`\gamma _\varphi `$ starts to decrease toward zero. Recently, constraints from high-z SNe Ia on power-law potentials with the field rolling with this set of initial conditions were obtained by Podariu and Ratra. We complement their analysis by including the lensing constraints as well. In the next section we show using the scalar field equations that present data prefer low values of $`\alpha `$. We also update and expand the observational constraints on the PNGB models .
## III Observational Constraints
In the following we briefly outline our main assumptions for lensing and supernovae analysis. Our approach for lensing statistics is based on Refs: and is described in more detail in . To perform the statistical analysis we consider data from the HST Snapshot survey (498 highly luminous quasars (HLQ)), the Crampton survey (43 HLQ), the Yee survey (37 HLQ), the ESO/Liege survey (61 HLQ), the HST GO observations (17 HLQ), the CFA survey (102 HLQ) , and the NOT survey (104 HLQ) . We consider a total of 862 ($`z>1`$) highly luminous optical quasars plus 5 lenses. The lens galaxies are modeled as singular isothermal spheres (SIS), and we consider lensing only by early-type galaxies, since they are expected to dominate the lens population. We assume a conserved comoving number density of lenses, $`n=n_0(1+z)^3`$, and a Schechter form for the early type galaxy population, $`n_0=_0^{\mathrm{}}n_{}\left(\frac{L}{L^{}}\right)^\alpha \mathrm{exp}\left(\frac{L}{L^{}}\right)\frac{dL}{L^{}}`$ , with $`n_{}=0.61\times 10^2h^3\text{Mpc}^3`$ and $`\alpha =1.0`$ . We assume that the luminosity satisfies the Faber-Jackson relation , $`L/L^{}=(\sigma _{||}/\sigma _{||}^{})^\gamma `$, with $`\gamma =4`$. Since the lensing optical depth depends upon the fourth power of the velocity dispersion of an $`L^{}`$ galaxy, a correct estimate of this quantity is crucial for strong lensing calculations. The image angular separation is also very sensitive to $`\sigma _{||}^{}`$: larger velocities give rise to larger image separations. In our likelihood analysis we take into account the observed image separation of the lensed quasars and adopt the value $`\sigma _{||}^{}=225`$ km/s, which gives the best fit to the observed image separations .
For SIS, the total lensing optical depth can be expressed analytically, $`\tau (z_S)=\frac{F}{30}(d_A(0,z_S)(1+z_S))^3(cH_0^1)^3`$, where $`z_S`$ is the source redshift, $`d_A(0,z_S)`$ is its angular diameter distance, and $`F=16\pi ^3n(cH_0^1)^3(\sigma _{||}^{}/c)^4\mathrm{\Gamma }(1+\alpha +4/\gamma )0.026`$ measures the effectiveness of the lens in producing multiple images . We correct the optical depth for the effects of magnification bias and include the selection function due to finite angular resolution and dynamic range . We assume a mean optical extinction of $`\mathrm{\Delta }m`$ =$`0.5`$ mag, as suggested by Falco et al. : this makes the lensing statistics for optically selected quasars consistent with the results for radio sources, for which there is no extinction. When applied to spatially flat cosmological constant models, our approach yields the upper bounds $`\mathrm{\Omega }_\mathrm{\Lambda }0.76`$ (at $`2\sigma `$) and $`\mathrm{\Omega }_\mathrm{\Lambda }0.61`$ (at $`1\sigma `$), with a best-fit value of $`\mathrm{\Omega }_\mathrm{\Lambda }0.39`$. Recent statistical analyses using both HLQ and radio sources slightly tighten these constraints on a cosmological constant . A combined (optical+radio) lensing analysis for dynamical-$`\mathrm{\Lambda }`$ models is still in progress; qualitatively, we expect this to tighten the lensing constraints below by approximately $`1\sigma `$.
For the SNe Ia analysis , we consider the latest published data from the High-z Supernovae Search Team . We use the 27 low-redshift and 10 high-redshift SNe Ia (including SN97ck) reported in Riess et al. and consider data with the MLCS method applied to the supernovae light curves. Following a procedure similar to that described in Riess et al., we determine the cosmological parameters through a $`\chi ^2`$ minimization, neglecting the unphysical region $`\mathrm{\Omega }_{m0}<0`$.
In Fig. 1 we show the $`95.4\%`$ and $`68.3\%`$ C. L. limits from lensing (short dashed contours) and the SNe Ia data (solid curves) on the parameters $`f`$ and $`M`$ of the PNGB potential. As in , these limits apply to models with the initial condition $`\frac{4\sqrt{\pi }\varphi (t_i)}{m_{Pl}}=1.5`$ and $`\frac{d\varphi }{dt}(t_i)=0`$, with $`t_i=10^5t_0`$ ; for other choices, the bounding contours would shift by small amounts in the $`fM`$ plane. We also plot some contours of constant $`\mathrm{\Omega }_{m0}`$ (dashed) and the curve $`q_0=0`$ (long dashed contour) as a function of the parameters $`f`$ and $`M`$. The allowed region (shown by the shaded area in Fig. 1) is limited by the lensing and SNe Ia $`95.4\%`$ C. L. contours and also by the constraint $`\mathrm{\Omega }_{m0}>0.15`$, which we interpret as $`2\sigma `$ lower bound from observations of galaxy clusters. The data clearly favors accelerated expansion (the region above the $`q_0=0`$ curve) but curiously there is a small region in the parameter space, close to the point where the $`\mathrm{\Omega }_{m0}=0.15`$ and the Sne Ia $`2\sigma `$ curves cross, where the Universe is not in accelerated expansion by the present time. This small area disappears if we adopt the tighter constraint $`\mathrm{\Omega }_{m0}>0.3`$. We note that the bulk of the $`2\sigma `$-allowed parameter space, where the lensing and SNe contours are nearly vertical, corresponds to the scalar field being nearly frozen, i.e., in this region the model is degenerate with a cosmological constant.
In Fig. 2 we show the $`95.4\%`$ and $`68.3\%`$ C. L. limits from lensing (thick dashed contours) and the SNe Ia data (solid curves) on the parameters $`\alpha `$ and $`\mathrm{\Omega }_{m0}`$ of the inverse power-law potential. The horizontal dotted line shows a lower bound on the matter density inferred from the dynamics of galaxy clusters, $`\mathrm{\Omega }_{m0}=0.15`$. We also show contours of the present equation of state $`w_0=\gamma _01`$ (thin dotted curves) and the curve $`q_0=0`$ (long dashed curve). At $`95.4\%`$ confidence, the SNe Ia and $`\mathrm{\Omega }_{m0}`$ constraints require $`\alpha <5`$ and $`w_0<0.5`$; the latter bound agrees roughly with the constraint obtained by assuming a time-independent equation of state , an approximation sometimes used for the inverse power-law model. We also observe that the lensing constraints on the model parameters are weak, constraining only low values of $`\mathrm{\Omega }_{m0}`$ and $`\alpha `$. We remark, however, that they are consistent with the SNe Ia constraints. We can tighten the constraints on the equation of state if we consider a higher value for the $`\mathrm{\Omega }_{m0}`$ lower bound. For instance, if we adopt $`\mathrm{\Omega }_{m0}>0.3`$, as suggested in , we obtain $`w_0<0.67`$ and $`\alpha <1.8`$. In both models, a larger lower bound on $`\mathrm{\Omega }_{m0}`$ pushes the scalar field behavior toward that of the cosmological constant ($`w=1`$).
## IV Conclusion
A consensus is beginning to emerge that we live in a nearly flat, low-matter-density Universe with $`\mathrm{\Omega }_{m0}0.3`$ and a dark energy, negative-pressure component with $`\mathrm{\Omega }_X0.7`$. The nature of this dark energy component is still not well understood; further developments will require deeper understanding of fundamental physics as well as improved observational tests to measure the equation of state at recent epochs, $`w(t)`$, and determine if it is distinguishable from that of the cosmological constant . Classical scalar field models provide a simple dynamical framework for posing these questions. In this paper we analyzed two representative scalar field models, the PNGB and power-law potentials, which span the range of expected dynamical behavior. The inverse power-law model displays tracking solutions which allow the scalar field to start from a wide set of initial conditions. We showed that current data favors a small value of the parameter, $`\alpha <5`$. This may be a problem for these models: in Refs: it was shown that, starting from the equipartition condition after inflation, it is necessary to have $`\alpha >5`$ for the field to begin tracking before matter-radiation equality. Since the observational constraints indicate that tracking could only be achieved (if at all) at more recent times, it is not clear what theoretical advantage, in terms of alleviating the ‘cosmic coincidence’ problem, is gained by the tracking solution. Although well motivated from the particle physics viewpoint, the PNGB model is strongly constrained by the SNe Ia and lensing data. Finally, as noted above, these two models predict radically different futures for the Universe. In the inverse power law model, the expansion will continue accelerating and approach de Sitter space. In the PNGB model, the present epoch of acceleration may be brief, followed by a return to what is effectively matter-dominated evolution.
###### Acknowledgements.
We would like to deeply thank Luca Amendola, Robert Caldwell, Cindy Ng, Franco Ochionero, Silviu Podariu and Bharat Ratra for several useful discussions that helped us to improve this work. This work was supported by the Brazilian agencies CNPq and FAPERJ and by the DOE and NASA Grant NAG5-7092 at Fermilab. |
warning/0001/cond-mat0001254.html | ar5iv | text | # d_{𝑥²-𝑦²} Pairing of Composite Excitations in the 2D Hubbard Model
## I Introduction
After a protracted quest for superconductivity in the Hubbard model, it is surprising that the question is still open. This state of affairs has arisen primarily because the parameter space in which superconductivity is anticipated to reside is precisely the strong coupling limit, $`U|t|`$, in which traditional perturbative schemes fail. As a result, progress on this question has relied predominantly on exact diagonalization or Quantum Monte Carlo (QMC) simulations on finite systems. While early work showed promising results on the possible onset of superconducting pair correlations in the 2D Hubbard model, the most recent numerical work shows that true superconductivity with at least algebraic decay of anomalous correlations is absent. In light of such results, it has also been suggested that some combination of phonons and or next nearest-neighbour hopping are needed for superconductivity to survive. Further, Kuroki and Aoki have suggested that QMC simulations have yet to access the energy scale of order $`0.01t`$ where strong-coupling superconductivity is expected to occur. In addition, Beenen and Edwards have analyzed the 2D Hubbard model with an equations of motion technique considerably enhanced by Mancini, Matsumoto, and co-workers and shown that pairs with $`d_{x^2y^2}`$ symmetry emerge. However, this work has key weaknesses: the onset temperature and the window of doping over which d$`_{x^2y^2}`$ pairing occurs are highly sensitive to the decoupling scheme of the anomalous Green functions. For example, at least order of magnitude variations in the onset temperature were observed for $`U=8|t|`$ and a decoupling scheme which is supposedly exact in the limit of $`U\mathrm{}`$ resulted in a near vanishing of pairing for $`U>10|t|`$.
Despite these difficulties, the Hubbard model remains central to such strongly correlated problems as superconductivity in the cuprates as well as the organic conductors. In fact, in so far as the motion of holes in the CuO<sub>2</sub> planes in the cuprates can be modeled realistically by the Hubbard model in the vicinity of half-filling, the microscopic underpinnings of the pairing mechanism in these materials should arise from this model. Because D=2 is the marginal dimension for superconductivity, it might be that while the pairing mechanism originates within a single layer, interlayer tunneling is needed to maintain long-range phase coherence. Hence, the questions that face the Hubbard model are three-fold: 1) does local pairing of the right sort obtain to describe the cuprates, 2) what entities are involved in the pairing, and 3) is long-range phase coherence present? In this work we develop a strong-coupling approach that is sufficient to answer the first two questions. Our approach is based on a simple observation first made by Hubbard. Namely, that the bare electron annihilation operator can be written as a sum of two composite operators
$`c_{i\sigma }=c_{i\sigma }n_{i\sigma }+c_{i\sigma }(1n_{i\sigma })=\eta _{i\sigma }+\xi _{i\sigma }.`$ (1)
The quantities described by $`\eta _{i\sigma }`$ and $`\xi _{i\sigma }`$ are composite excitations. The $`\eta `$ excitation describes an electron restricted to move on sites already occupied with an electron of opposite spin whereas $`\xi `$ demands that there be no prior occupancy on the site. A key feature these operators possess is that they exactly solve the $`t=0`$ Hubbard model. Consequently, an equation of motion approach based on these operators will lead naturally to an expansion in $`t/U`$. Such an expansion is ideally suited for the strong coupling regime $`U|t|`$ in which both the cuprates as well as the organic conductors reside. It is this approach that we pursue here.
When the kinetic energy is treated as a perturbation, hybridization is introduced into the Hubbard atomic orbital basis. As a result, $`\xi _i`$ and $`\eta _i`$ are no longer the eigen-operators. The new type of excitations can be thought of heuristically as resonant valence-bond hybrids. To determine what operators create such excitations, it is sufficient to construct the equations of motion for the Hubbard operators. The relevant operators have the kinetic energy, $`t`$ as a coefficient. As we will see, the simplest terms that arise are of the form, $`\eta _{ij\sigma }=c_{i\sigma }n_{j\tau }`$ with $`(ij)`$ nearest neighbours. This operator is the generator of a composite excitation (or cexon) that is restricted to live on sites $`i`$ with a nearest neighbour site $`j`$ occupied. Such composites have unit charge as can be seen directly from the commutator,
$`[\eta _{ij\sigma },n_{l\sigma ^{}}]=\delta _{\sigma ,\sigma ^{}}\delta _{il}\eta _{ij\sigma }.`$ (2)
Also, quite obviously, its spin is 1/2. In fact, all the composite excitations generated in our scheme have unit charge. Hence, spin and charge are coupled. However, while the $`\eta _{ij}`$ excitation is fermionic in character, it is not confined to a single site. Rather, it is an excitation that lives on neighbouring sites. In traditional perturbative schemes, such as Fermi liquid theory, similar composite operators naturally appear. However, in such approaches, the operators that are generated do not describe new excitations but rather dress the non-interacting quasiparticles. In this approach, the new operators do in fact describe fundamentally new excitations as they give rise to new peaks in the spectral function. The appearence of new peaks in the spectral function is incompatible with a Fermi liquid description of the composite excitations.
It might be argued that the emergence of cexons in this model and the unrecoverability of Fermi liquid theory are inevitable in an approach based fundamentally on the atomic limit of the Hubbard model. However, any theory of high $`T_c`$ must be based on a starting point that captures the essence of strong-correlation physics. While the atomic limit might seem extreme, it is justified in this context as $`Ut`$ in high $`T_c`$ systems. Consequently, it stands to reason that the physically-relevant excitations should bear a greater family resemblance to the atomic limit eigentates than they would to the plane-wave excitations which populate the non-interacting limit. Cexons are the simplest local excitations that arise once the Hubbard operators are hybridized on nearest-neighbour sites. To determine if such excitations really do constitute new particles in the strongly-coupled limit, dynamical corrections must be included to the theory we present below.
As in earlier composite operator approaches , we find that the traditional correlation function, $`c_ic_j`$, does not determine the pairing gap in the 2D Hubbard model. Rather
$`\theta _{ij}=c_ic_in_{j\tau }=\eta _{ij}\eta _{ij}`$ (3)
is the relevant anomalous correlation function that determines pairing in the d-wave channel. The emergence of this correlation function as the d-wave order parameter in the Hubbard model and the seemingly innocuous rewriting (overlooked previously in similar Green function analyses) in terms of the $`\eta _{ij\sigma }`$ operators is particularly illuminating. Pairing based on the $`\eta _{ij\sigma }`$ excitations is the new feature which we develop explicitly in this work. This rewriting signifies that composite excitations rather than electron-like particles, as would be the case in a traditional Fermi liquid, are involved in the pairing process in a strongly-correlated system with repulsive interactions. In addition, the fact that the new composite operators describe the pairing lends further credence to their correspondence with new excitations in the strongly-correlated regime. While it is well known that quasiparticles in strongly-correlated systems do not resemble well electron-like excitations, as shown, for example, by the absence of a sharp peak at the Fermi surface in angle-resolved photoemission experiments on the under-doped normal state of the cuprates, the traditional mechanism invoked to explain this fact is spin-charge separation. In 1-dimension spin-charge separation is on firm footing. In addition, even in a singlet BCS superconductor, spin-charge separation occurs naturally as the Cooper pairs carry charge but no spin, whereas the quasiparticles are spin 1/2 but have no well-defined charge. In D=2, fractionalization of the electron also occurs in the fractional quantum Hall effect. For strongly-correlated systems, spin-charge separation has been proposed based on a slave-boson picture. However, recently Nayak has shown that the slave particles always remain confined. Our work suggests that rather than falling apart, electrons form clusters or composites when the Coulomb interaction is large and repulsive. Such entities form pairs in the 2D Hubbard model. We propose that composite excitations of the kind we describe here offer a natural explanation for the absence of electron-like excitations in the ARPES experiments on the underdoped cuprates. While it is tempting to draw a connection between pairing of composite excitations in the 2D Hubbard model and the pairing of composite fermions in quantum Hall systems, this connection would be tenuous at best.
Our approach then accomplishes two things: 1) First we are able to access the strong-coupling limit of the Hubbard model and identify the relevant composite excitations. 2) We are able to show that such composites pair in the $`d_{x^2y^2}`$ channel. Our work then represents a formulation of d-wave pairing from a microscopic model out of which emerges a hierarchy of new composite fermionic excitations. We find that the magnitude of the anomalous correlation function for composite excitation pairing increases as the compressibility in the normal state becomes more negative. In the next section, we outline the essential details of the strong-coupling approach. In Sec. III, we derive the equations of motion for the composite excitations and all the relevant correlation functions. The results for the anomalous pairing correlation functions are presented in Sec. IV.
## II Formalism
The starting point of our analysis is the on-site Hubbard model
$`H={\displaystyle \underset{i,j,\sigma }{}}t_{ij}c_{i\sigma }^{}c_{\sigma j}+U{\displaystyle \underset{i}{}}n_in_{}\mu {\displaystyle \underset{i}{}}n_i`$ (4)
where $`t_{ij}=t`$ if $`(i,j)`$ are nearest-neighbour sites and zero otherwise, $`\mu `$ is the chemical potential and $`U`$ the on-site Coulomb repulsion. Consider the retarded Green function
$`c_i;c_j=\theta (tt^{})\{c_i(t),c_j(t^{})\}`$ (5)
where $`\{A,B\}`$ is the anticommutator. As a result of the interaction term in Eq. (4), the equation of motion for this Green function will generate a new Green function that is proportional to $`[c_i,H];c_j`$. This Green function will contain the composite excitation $`\eta _{i\sigma }`$. All higher-order Green functions will emerge from time derivatives of composite operators. This suggests that we should start with the Hubbard operators, $`\eta _{\sigma i}=c_{i\sigma }n_{i\sigma }`$ and $`\xi _\sigma (i)=c_{i\sigma }(1n_{i\sigma })`$ introduced previously. The utility of these operators is immediate. Their commutator with the interaction part of the Hubbard model yields $`(\mu U)\eta _{i\sigma }`$ and $`\mu \xi _{i\sigma }`$, respectively. Hence, they can be used to diagonalise the $`t=0`$ Hubbard model. Consequently, if this basis is used to compute correlation functions, all higher-order correlation functions that will be generated will be multiplied by the hopping term. That is, these operators form the basis for a $`t/U1`$ or equivalently a strong coupling expansion. The essence of this approach was first proposed by Linderberg and Öhrn and has been applied to the Hubbard model as well as the p-d model for the cuprates.
To illustrate the utility of the composite operator technique, we switch to the 4-component basis,
$`\psi (i)=\left(\begin{array}{c}\xi _i\hfill \\ \eta _i\hfill \\ \xi _i^{}\hfill \\ \eta _i^{}\hfill \end{array}\right).`$ (10)
Let us define
$`j(i)=i_t\psi (i)=[\psi (i),H]`$ (11)
as the ‘current’ operator. Formally, we can write the n<sup>th</sup> element of the current as
$`j_n(i)={\displaystyle \underset{m}{}}K_{nm}\psi _m(i)+\delta j_n(i).`$ (12)
We can project out the part of the correction, $`\delta j`$, that is ‘orthogonal’ to the composite operator basis by rewriting
$`\delta j_n={\displaystyle \underset{ml}{}}\{\delta j_n,\psi _m^{}\}I_{ml}^1\psi _l+\delta \varphi _n,`$ (13)
where $`\delta \varphi `$ satisfies the equation, $`\{\delta \varphi ,\mathrm{\Psi }^{}\}=0`$. Let us introduce the normalization matrix
$`I_{il}=\{\psi _i,\psi _l^{}\}={\displaystyle \frac{\mathrm{\Omega }}{(2\pi )^2}}{\displaystyle d^2ke^{i𝐤(𝐫_i𝐫_𝐥)}I(𝐤)}`$ (14)
and the overlap matrix
$`M_{il}=\{j_i,\psi _l^{}\}={\displaystyle \frac{\mathrm{\Omega }}{(2\pi )^2}}{\displaystyle d^2ke^{i𝐤(𝐫_i𝐫_l)}M_{il}(𝐤)}`$ (15)
where k-integration is over the Brillouin zone and $`\mathrm{\Omega }`$ the inverse area of the Brillouin zone. The elements of the $`𝐈`$ and $`𝐊`$ matrices contain only the mutual correlations among the constituents of the composite operator basis. Because $`\{\delta \varphi ,\mathrm{\Psi }^{}\}=0`$, we can write the overlap matrix as
$`𝐌=\mathrm{𝐄𝐈}`$ (16)
where
$`E_{nm}=K_{nm}+{\displaystyle \underset{l}{}}\{\delta j_n,\psi _l^{}\}I_{lm}^1.`$ (17)
This matrix (the energy matrix) will play a crucial role in our approach. Because of the $`\delta j`$ contribution, the energy matrix contains higher order correlations which are not expressible as simple correlations of the $`\psi `$-fields. It is from the $`\delta j`$ terms that the composite excitations arise. Consequently, the higher-order correlation functions must either be decoupled or used as self-consistent parameters in equations obtained by imposing various symmetries, such as the Pauli principle. As a result of the orthogonality of $`\delta \varphi `$ to the composite operator basis, we will neglect the contribution from $`\delta \varphi `$ in all subsequent calculations. This approximation is equivalent to ignoring the dynamical corrections to the self-energy. Such contributions are expected to be small because the Hubbard operator basis solves exactly the $`U=\mathrm{}`$ limit. The principal role of $`\delta \varphi `$ is to broaden the energy levels associated with the composite operator spectrum. The role of such dynamical corrections will be the focus of future study. Within this approximation, we can write any Green function
$`S(i,j)=\psi (i);\psi ^{}(j)`$ (18)
in Fourier space as
$`S(𝐤,\omega )=(\omega 𝐄(𝐤))^1𝐈(𝐤)`$ (19)
by using the equations of motion. This equation and the expression, $`𝐌=\mathrm{𝐄𝐈}`$, are the central equations of this approach. The M and I matrices involve correlation functions which can be obtained self-consistently and by introducing symmetry properties such as the Pauli principle to close the equations of motion. It is the off-diagonal blocks of the energy matrix that contain information about anomalous pairing correlations.
There are three types of self-consistent equations that enter this approach. The simplest of these involves the correlations between the $`\psi `$ fields that appear in $`𝐌`$ and $`𝐈`$. The general expression for such correlation functions in terms of the corresponding Green function is given by
$`\psi _m(i)\psi _n^{}(j)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{(2\pi )^2}}{\displaystyle d^2k𝑑\omega e^{i𝐤(𝐫_i𝐫_j)}(1f(\omega ))}`$ (21)
$`\times \left({\displaystyle \frac{1}{\pi }}\right)\mathrm{Im}S_{mn}(𝐤,\omega )`$
$``$ $`C_{m,n}(i,j)`$ (22)
with $`f(\omega )`$ the Fermi-Dirac distribution function. The self-consistency of Eq. (21) follows from the dependence of $`𝐄`$ and $`𝐈`$ on correlations in the composite operator basis. To obtain explicitly the self-consistent equations, we rewrite Eq. (19) as
$`S(𝐤,\omega )={\displaystyle \underset{i=1,2}{}}\left[{\displaystyle \frac{\kappa _i^+}{\omega ϵ_i+i\eta }}+{\displaystyle \frac{\kappa _i^{}}{\omega +ϵ_ii\eta }}\right]`$ (23)
where
$`\kappa _i^\pm ={\displaystyle \frac{\lambda (\pm ϵ_i)}{\pm 2ϵ_i(ϵ_i^2ϵ_j^2)}}`$ (24)
with $`(i,j)`$ chosen from $`(1,2)`$ but $`ij`$, $`ϵ_i`$ the eigenvalues of the energy matrix and the matrix $`\lambda `$ is given by
$`\lambda (\omega )=\mathrm{Det}(\omega \mathrm{𝟏}𝐄)(\omega \mathrm{𝟏}𝐄)^1𝐈.`$ (25)
The four eigenvalues correspond to the four composite operator bands: i=1,2 refer to the $`\xi `$ and $`\eta `$ bands, respectively, whereas $`+,`$ index the particle and hole states, respectively. If we now use Eq. (23) in Eq.(21), we obtain the general self-consistent equation,
$`C_{m,n}(i,l)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{2(2\pi )^2}}{\displaystyle }d^2ke^{i𝐤(𝐫_i𝐫_l)}[I_{mn}+`$ (27)
$`{\displaystyle \underset{i=1,2}{}}((\kappa _i^+)_{mn}(\kappa _i^{})_{mn})]\mathrm{tanh}{\displaystyle \frac{\beta ϵ_i}{2}}`$
for correlation functions within the composite operator basis where $`T=1/k_B\beta `$.
It is well-known that in approximation schemes of this sort, certain correlation functions which vanish explicitly as a result of the Pauli principle, self-consistently iterate to a non-zero value. To avoid this shortcoming, we explicitly maintained the symmetries imposed by the Pauli principle, namely
$`C_{1,2}(i,i)=\xi _\sigma (i)\eta _\sigma ^{}(i)=0.`$ (28)
This is the second type of self-consistent integral equation that will be used in this method. Likewise, $`C_{2,1}(i,i)=0`$. As a result of imposing these symmetry relations, we will find that our theory is completely particle-hole symmetric about half-filling.
The final self-consistent equations arise from decoupling correlations of composite fields that are not contained in the basis described by Eq. (10). This procedure will be described in detail in the next section.
## III Computational Machinery
### A Equations of Motion
We start first by constructing the algebra
$`\{\xi _i,\xi _j^{}\}=\delta _{ij}\left(\mathrm{𝟏}+{\displaystyle \frac{1}{2}}\sigma ^\mu n_\mu \right)`$ (29)
$`\{\eta _i,\eta _j^{}\}=\delta _{ij}{\displaystyle \frac{1}{2}}\sigma ^\mu n_\mu `$ (30)
$`\{\xi _i,\xi _j\}=\{\eta _i,\eta _j\}=\{\xi _i,\eta _j^{}\}=0`$ (31)
of the composite operators, with $`\sigma ^\mu `$ for $`\mu =1,2,3`$ are the Pauli matrices, $`\sigma ^0`$ is the identity matrix, and
$`{\displaystyle \frac{1}{2}}\sigma ^\mu n_\mu =\left(\begin{array}{cc}n_i& c_i^{}c_i\\ c_i^{}c_i& n_i\end{array}\right).`$ (34)
From these equations, it follows that the I-matrix is given by
$`I=\left(\begin{array}{cccc}1\frac{n}{2}& 0& 0& \mathrm{\Delta }_0\\ 0& \frac{n}{2}& \mathrm{\Delta }_0& 0\\ 0& \mathrm{\Delta }_0^{}& 1\frac{n}{2}& 0\\ \mathrm{\Delta }_0^{}& 0& 0& \frac{n}{2}\end{array}\right)`$ (39)
where $`\mathrm{\Delta }_0=c_ic_i`$ and we have assumed that $`n_{}=n_{}=n/2`$. We now turn to the “current”. We can construct the “current” from the equations of motion
$`j(i)=i{\displaystyle \frac{}{t}}\psi (i)=[\psi (i),H]`$ (40)
for the $`\psi `$-fields with
$`j_1(i)`$ $`=`$ $`\mu \xi _i{\displaystyle \underset{j}{}}t_{ij}c_j4t\pi _i`$ (41)
$`j_2(i)`$ $`=`$ $`(\mu U)\eta _i+4t\pi _i`$ (42)
$`j_3(i)`$ $`=`$ $`\mu \xi _i^{}+{\displaystyle \underset{j}{}}t_{ij}c_j^{}+4t\pi _i^{}`$ (43)
$`j_4(i)`$ $`=`$ $`(\mu U)\eta _i^{}4t\pi _i^{}`$ (44)
and
$`4t\left(\begin{array}{c}\pi _i\\ \pi _i\end{array}\right)={\displaystyle \underset{j}{}}t_{ij}\left(\begin{array}{c}n_ic_j+c_i^{}c_ic_jc_ic_ic_j^{}\\ n_ic_j+c_i^{}c_ic_j+c_ic_ic_j^{}\end{array}\right).`$ (49)
We also need the M-matrix. The distinct elements of this matrix are
$`M_{11}(𝐤)=FT\{j_1,\xi _{}^{}\}`$ $`=`$ $`\mu (1{\displaystyle \frac{n}{2}})4te`$ (51)
$`4t\alpha (𝐤)(1n+p)`$
$`M_{12}(𝐤)=FT\{j_1,\eta _{}^{}\}`$ $`=`$ $`4te4t\alpha (𝐤)({\displaystyle \frac{n}{2}}p)`$ (52)
$`M_{22}(𝐤)=FT\{j_2,\eta _{}^{}\}`$ $`=`$ $`(\mu U){\displaystyle \frac{n}{2}}4te4t\alpha p`$ (53)
$`M_{13}(𝐤)=FT\{j_1,\xi _{}\}`$ $`=`$ $`4t\gamma (𝐤)\theta +8t\alpha \mathrm{\Delta }_04t\mathrm{\Delta }_{c\xi }`$ (54)
$`M_{14}(𝐤)=FT\{j_1,\eta _{}\}`$ $`=`$ $`(\mu +4t\alpha (𝐤))\mathrm{\Delta }_0`$ (56)
$`+4t\gamma (𝐤))\theta 4t\mathrm{\Delta }_{c\eta }`$
$`M_{23}(𝐤)=FT\{j_2,\xi _{}\}`$ $`=`$ $`(\mu U4t\alpha (𝐤))\mathrm{\Delta }_0`$ (58)
$`+4t\gamma (𝐤))\theta +4t\mathrm{\Delta }_{c\xi }`$
$`M_{24}(𝐤)=FT\{j_2,\eta _{}\}`$ $`=`$ $`4t\gamma (𝐤)\theta +4t\mathrm{\Delta }_{c\eta }`$ (59)
with the normal correlations
$`e=\xi _i^\alpha \xi _i^{}\eta _i^\alpha \eta _i^{},`$ (60)
and
$`p=n_{i\sigma }n_{i\sigma }^\alpha +c_i^{}c_i(c_i^{}c_i)^\alpha c_ic_i(c_i^{}c_i^{})^\alpha ,`$ (61)
and anomalous correlations
$`\mathrm{\Delta }_{c\xi }=c_i\xi _i^\alpha c_i\xi _i^\alpha ,`$ (62)
$`\mathrm{\Delta }_{c\eta }=c_i\eta _i^\alpha c_i\eta _i^\alpha ,`$ (63)
and
$`\theta =c_ic_i[n_{}(𝐫_i+\widehat{x})+n_{}(𝐫_i+\widehat{x})].`$ (64)
In these equations, $`FT`$ signifies the Fourier transform as defined in Eq. (15), $`\alpha `$ represents an average over nearest-neighbour sites, $`\widehat{x}`$ indexes nearest-neighbour sites, and $`\alpha (𝐤)=\mathrm{cos}k_x+\mathrm{cos}k_y`$. Unlike $`\alpha (𝐤)`$ which is the coefficient of the normal correlation functions, the Fourier coefficient $`\gamma (𝐤)`$ is a coefficient of an anomalous correlation function. Hence, it is sensitive to the sign change of the anomalous correlation functions as $`\mathrm{exp}(i𝐤𝐫)`$ is summed over nearest neighbour sites. For s-wave symmetry, $`\alpha (𝐤)=\gamma (𝐤)`$ while for d-wave symmetry $`\gamma (𝐤)=\mathrm{cos}k_x\mathrm{cos}k_y`$. The symmetry relationships among the elements of $`𝐌`$ are as follows:
$`M_{11}`$ $`=`$ $`M_{33}`$ (65)
$`M_{12}`$ $`=`$ $`M_{34}`$ (66)
$`M_{22}`$ $`=`$ $`M_{44}`$ (67)
$`M_{14}`$ $`=`$ $`M_{32}^{}.`$ (68)
The $`M_{13}`$, $`M_{14}`$, $`M_{23}`$, and $`M_{24}`$ contain the anomalous correlation functions in particular the $`\theta `$ correlation discussed in the introduction.
With the $`𝐌`$ and $`𝐈`$-matrices in hand, we now construct the energy matrix. From the structure of the normalization matrix and the symmetry properties of the $`M`$ matrix (Eq. (65), it follows that the energy matrix is of the form,
$`𝐄=\left(\begin{array}{cc}𝐀\hfill & \hfill 𝐁\\ 𝐁\hfill & \hfill 𝐀\end{array}\right),`$ (71)
where $`𝐀`$ and $`𝐁`$ are $`2\times 2`$ matrices. The off-diagonal blocks of the energy matrix determine the gap in the energy spectrum induced by pairing. To isolate the relevant parts of the normalization and $`𝐌`$ matrices that enter the energy gap, we consider the candidate pairing symmetries. No simplifications occur in the s-channel. However, we have verified that none of the anomalous correlation functions are non-zero either in this channel. The same state of affairs occurs in p-wave symmetry. What about d-wave symmetry? Here several significant simplifications occur. First, as a result of the nodes in the gap, no purely on-site anomalous correlation functions survive. As a consequence, $`c_ic_i=\mathrm{\Delta }_0=0`$. This leads to a vanishing of all the off-diagonal blocks of the normalization matrix. Consequently, only the off-diagonal blocks of the $`𝐌`$-matrix enter the off-diagonal elements of the energy matrix. However, further simplifications occur. Consider for example, $`\mathrm{\Delta }_{c\xi }`$ and $`\mathrm{\Delta }_{c\eta }`$ in Eqs. (62) and (63). These quantities are symmetrized and in addition involve a sum over the nearest neighbour sites in the Brillouin zone. Consequently, they vanish identically in the d-wave channel. The only anomalous correlation function that remains is $`\theta `$, Eq. (64). At this level of theory, this is the only correlation function that does not vanish by symmetry conditions. This is the only anomalous correlation function that enters the gap in the energy spectrum. Hence, as advertised, any pairing will necessarily be governed by a non-traditional correlation function. This correlation function involves a composite operator that lives on a cluster of nearest-neighbour sites. By defining $`\eta _{ij\sigma }`$ to be $`c_{i\sigma }n_{j\tau }`$, we rewrite Eq. (64) as Eq. (3) in which it is clear that pairing in $`\theta `$ occurs between two composite excitations.
### B Computational Procedure
What remains to be done is the calculation of the anomalous correlation functions. Our assumption of singlet pairing implies that
$`c_ic_in_{j\tau }=c_ic_in_{j\tau }=c_i^{}c_i^{}n_{j\tau }^{}.`$ (72)
Consequently, we can define $`\theta `$ as
$`\theta =2c_ic_in_{j\tau }`$ (73)
where the spin $`\tau `$ is arbitrary. As this correlation function involves more than two basis operators, it must be decoupled. We will follow a procedure analogous to that devised by Roth in her treatment of the strong-coupling limit of the Hubbard model. To implement the Roth method, we consider the Green function
$`H(i,j,t,t^{})=c_{}^{}(i,t);c_{}(i,t^{})^{}n_\sigma (j,t^{}).`$ (74)
Clearly, if this Green function is calculated, $`\theta `$ can be obtained directly from an expression analogous to Eq. (21). We proceed by constructing the series of Green functions,
$`A_\sigma (i,j,k,t,t^{})`$ $`=`$ $`\xi _{}(i,t);c_{}(j,t^{})^{}n_\sigma (k,t^{})`$ (75)
$`B_\sigma (i,j,k,t,t^{})`$ $`=`$ $`\eta _{}(i,t);c_{}(j,t^{})^{}n_\sigma (k,t^{})`$ (76)
$`F_\sigma (i,j,k,t,t^{})`$ $`=`$ $`\xi _{}^{}(i,t);c_{}(j,t^{})^{}n_\sigma (k,t^{})`$ (77)
$`G_\sigma (i,j,k,t,t^{})`$ $`=`$ $`\eta _{}^{}(i,t);c_{}(j,t^{})^{}n_\sigma (k,t^{})`$ (78)
in which the creation operator at time $`t`$ is replaced by each of the composite operators in the 4-component basis. The Green function of interest, $`H`$, is obtained by summing $`F`$ and $`G`$ and setting $`i=j`$. To calculate these quantities, we use the equations of motion for $`\psi `$, Eq. (11), and in particular the approximation that $`j=i_t\psi E\psi `$. Consequently, the equations of motion for the Green functions,
$`i_t\left(\begin{array}{c}A_\sigma \hfill \\ B_\sigma \hfill \\ F_\sigma \hfill \\ G_\sigma \hfill \end{array}\right)=E(i)\left(\begin{array}{c}A_\sigma \hfill \\ B_\sigma \hfill \\ F_\sigma \hfill \\ G_\sigma \hfill \end{array}\right)+i\delta (tt^{})\left(\begin{array}{c}f_{1\sigma }\hfill \\ f_{2\sigma }\hfill \\ f_{3\sigma }\hfill \\ f_{4\sigma }\hfill \end{array}\right)`$ (91)
are directly related to the energy matrix in real space, $`E(i)`$ and the term, $`f_{n\sigma }=\{\psi _n(i,t),c_{}(j,t)^{}n_\sigma (k,t)\}`$, that arises from the equal-time anticommutation of $`\psi `$ with the composite operators in the Green function. Note that in general, $`f_n`$ is a linear combination of correlation functions in the composite operator basis and correlations associated with $`A+B`$ or $`F+G`$. Consequently, we now have a closed set of equations from which $`\theta `$ can be obtained.
In the next step, we Fourier transform these equations of motion so that $`t\omega `$, $`r_kr_ik_1`$, and $`r_jr_ik_2`$. Noting that $`\omega E=IS^1`$, we obtain
$`\left(\begin{array}{c}A_\sigma (k_1,k_2,\omega )\hfill \\ B_\sigma (k_1,k_2,\omega )\hfill \\ F_\sigma (k_1,k_2,\omega )\hfill \\ G_\sigma (k_1,k_2,\omega )\hfill \end{array}\right)=S(k_1+k_2,\omega )I^1\left(\begin{array}{c}f_{1\sigma }(k_1,k_2)\hfill \\ f_{2\sigma }(k_1,k_2)\hfill \\ f_{3\sigma }(k_1,k_2)\hfill \\ f_{4\sigma }(k_1,k_2)\hfill \end{array}\right)`$ (100)
To extract $`\theta `$ from these equations, we sum $`F`$ and $`G`$, integrate over $`k_2`$ so that $`i=j`$ and then using Eq. (21), we find that
$`\theta ={\displaystyle \frac{2\zeta }{\varphi +1}}`$ (101)
where
$`\varphi `$ $`=`$ $`{\displaystyle \frac{n^24D}{n(2n)}}`$ (102)
$`\zeta `$ $`=`$ $`{\displaystyle \frac{2}{2n}}\left(C_{11}(i,j)+C_{12}(i,j)\right)\left(C_{13}(i,j)+C_{14}(i,j)\right)`$ (104)
$`+{\displaystyle \frac{2}{n}}\left(C_{22}(i,j)+C_{12}(i,j)\right)\left(C_{24}(i,j)+C_{14}(i,j)\right),`$
and $`D`$ is the on-site double occupancy
$`D=n_in_i={\displaystyle \frac{n}{2}}C_{22}(i,i).`$ (105)
In the expression for $`\varphi `$, $`(i,j)`$ denote nearest neighbour sites along the x-axis. We point out that Beenan and Edwards used a simple factorization scheme for $`\theta `$. As these authors indicate, there is no unique way of performing this procedure. In our scheme, there are two distinct ways of decoupling the Green functions that lead to $`\theta `$. That is, we could have started with another Green function
$`H^{}(i,j,t,t,t^{})=c_{}(i,t);c_{}(i,t^{})c_{}^{}(j,t^{})c_{}^{}(j,t^{}).`$ (106)
Following the procedure outlined above, we arrive at the alternative expression for $`\theta `$:
$`\theta ={\displaystyle \frac{2\zeta ^{}}{\varphi +1}}.`$ (107)
with
$`\zeta ^{}`$ $`=`$ $`{\displaystyle \frac{2}{2n}}\left(C_{11}(i,j)+C_{12}(i,j)\right)\left(C_{24}(i,j)+C_{14}(i,j)\right)`$ (109)
$`+{\displaystyle \frac{2}{n}}\left(C_{22}(i,j)+C_{12}(i,j)\right)\left(C_{13}(i,j)+C_{14}(i,j)\right).`$
In the following section, we compare the values of $`\theta `$ obtained by both methods.
Similarly it is possible to decouple the parameter $`p`$ defined in Eq. (61). Because Eq. (61) contains only symmetric combinations of the fermionic operators, it can be decoupled uniquely. From the procedure outlined above, we find that $`p`$ is given by
$`p={\displaystyle \frac{n^2}{4}}{\displaystyle \frac{\rho _1+\varphi \rho _2^2}{1\varphi ^2}}{\displaystyle \frac{\rho _1+\rho _2}{1\varphi }}{\displaystyle \frac{\rho _3}{1+\varphi }}`$ (110)
where
$`\rho _1`$ $`=`$ $`{\displaystyle \frac{2}{2n}}\left(C_{11}(i,j)+C_{12}(i,j)\right)^2+{\displaystyle \frac{2}{n}}\left(C_{22}+C_{12}\right)^2,`$ (111)
$`\rho _2`$ $`=`$ $`{\displaystyle \frac{2}{2n}}((C_{13}(i,j)+C_{14}(i,j))^2+{\displaystyle \frac{2}{n}}(C_{24}+C_{14})^2`$ (112)
and
$`\rho _3={\displaystyle \frac{4}{n(2n)}}\left(C_{11}(i,j)+C_{12}(i,j)\right)\left(C_{22}+C_{12}\right).`$ (114)
In the normal state, $`\rho _2=0`$ as it contains only anomalous correlations and $`p`$ reduces to the result derived by Roth. The expressions for $`\theta `$, $`p`$ and the self-consistent equation for the correlation functions constitute the working equations of this method.
To summarize, the equations for $`𝐌`$ and $`𝐈`$ contain four parameters, $`\mu `$, $`e`$, $`p`$ and $`\theta `$ that must be determined self-consistently. From Eq. (60) and the definition of the Hubbard operators, the self-consistent equations for $`e`$ and $`n`$ are
$`e`$ $`=`$ $`C_{11}(i,j)C_{22}(i,j)`$ (115)
$`n`$ $`=`$ $`2(1C_{11}(i,i)2C_{12}(i,i)C_{22}(i,i)).`$ (116)
In our computational procedure, we either used the two previous equations together with Eqs. (101) or (107). For the fourth self-consistent equation, we either used $`p`$ or the equation that imposes the Pauli principle, Eq. (28).
## IV Results
To set the stage for our results on the anomalous correlations, we discuss first the chemical potential. Shown in Fig. (1) are three different calculations of the chemical potential as a function of the filling: 1) solid line– present method, 2) dashed-dotted line–Beenan and Edwards, and 3) dashed line–Avella and co-workers. Here, $`n=1`$ corresponds to half-filling. The method of Avella and co-workers is identical to the method used here. However, as is evident, our results are significantly different. Particularly striking is the continuous increase of the chemical potential obtained by Avella and colleagues as the filling is increased. This is in stark contrast our solution which levels off in the vicinity of half-filling. The difference between our treatments lies in that there are two solutions to the self-consistent equations. Avella and co-workers chose the solution that is higher in energy in the filling range $`0.71.0`$. While the lower energy solution is the physically-relevant solution, this solution in the absence of pairing has a distinct negative compressibility in the vicinity of half-filling. It is for this reason that Avella and co-workers criticized, the method used by Roth as it also gives rise to a negative compressibility as seen from the dashed-dotted line of Beenan and Edwards. Avella and co-workers advocated that imposing the Pauli principle in the self-consistent solution to the equations of motion eliminates the negative compressibility.
Our work shows that even if the Pauli principle is maintained by means of Eq. (28), a second solution to the integral equations still exists along which the compressibility remains negative. To explore further the negative compressibility, we show in Fig. (2) the role of pairing on the compressibility. The dashed line corresponds to the chemical potential in the absence of pairing while along the solid line $`\theta 0`$. As is evident, pairing alleviates the negative compressibility almost entirely giving rise to a flattening of the chemical potential in the vicinity of half-filling as seen from the solid curve in Fig. (2). In fact, a key trend common to the curves shown in Fig. (1) is that the compressibility is most positive when pairing vanishes. This result is particularly important because a negative compressibility occurs in numerous dilute electron systems, such as the 2D electron gas for $`r_s>3`$. Our results also corroborate the earlier observation of Tandon, Wang, and Kotliar that the compressibility becomes negative in the $`U\mathrm{}`$ Hubbard model. For short-range Coulomb interactions, a negative compressibility signifies that the ground state of an electronic system is unstable relative to a uniform charge distribution. Hence, a negative compressibility is typically associated with phase separation or stripe formation. Our work explicitly shows that pairing alleviates this instability at least in the case of short-range Coulomb interactions. We have proposed that even in the case of long-range Coulomb interactions, pairing still obtains and alleviates the negative compressibility as well.
We now discuss explicitly the anomalous correlation functions. Shown in Fig. (3) are four different calculations of the correlation function involving pairing of composite excitations. The two solid lines in Fig. (3) were obtained from Eq. (107). On the upper curve, the Pauli principle was imposed whereas along the lower curve, the decoupling scheme was used to calculate the parameter $`p`$ defined in Eq. (61). Hence, the lower curve corresponds to the method of Beenan and Edwards. On the dashed curves, Eq. (101) was used to compute $`\theta `$. Once again, along the upper dashed curve, the Pauli principle was used whereas along the lower curve, the decoupling scheme (for $`p`$) was used. Our results indicate that the anomalous correlations are largest and most stable when the Pauli principle is imposed by using the integral equation, Eq. (28). When the decoupling scheme is used to obtain $`p`$, the two distinct decouplings for $`\theta `$ yield vastly different results. This will result in huge fluctuations in $`T_c`$ as the work of Beenan and Edwards illustrates. What our work establishes is that if the Pauli principle is imposed in the self-consistent procedure outlined here, consistent pairing solutions exist regardless of the decoupling scheme used to calculate $`\theta `$. Note also that the two lower curves are peaked around a doping level of 20$`\%`$. Beenan and Edwards associated great significance to this doping level as it corresponds to the filling at which the Fermi surface resembles the Fermi surface at half-filling for the non-interacting system. This appears to be an accident of their approximations as our more accurate method shows that the peak in the order parameter occurs at a doping level of 10$`\%`$.
The remaining figures constitute the primary results of this method. Shown in Fig. (4) is a calculation of $`\theta `$ in which the average of Eq. 101) and (107) as a function of filling was used. Our results show clearly that $`d_{x^2y^2}`$ pairing exists and such pairing is diminished as $`U`$ increases. As our method is valid only in the large $`U`$ limit, we cannot establish the lower value of $`U/t`$ for which $`\theta `$ is non-zero. Also of significance is the maximum in $`\theta `$ at roughly 10$`\%`$ doping for the range of on-site repulsions studied. In addition, the peak in $`\theta `$ moves to higher dopings as $`U`$ increases. We emphasize that the pair-formation found here does not include the effect of phase fluctuations. Hence, the acutal doping regime over which true superconductivity with long-range order occurs might be significantly smaller than that obtained here. For example, if long-range order obtains at all, the optimum doping will appear shifted to higher doping values relative to the maxima of the curves shown here. Preliminary results by Manske, Dahm and Bennemann using phenomenological spin-fluctuation approximations offer some support of this conclusion. Also of note is the fact that our treatment preserves the particle-hole symmetry of $`\theta `$. The temperature dependence of the order parameter is shown in Fig. (5). As is evident, the shape of $`\theta `$ is characteristic of any order parameter that vanishes at a transition temperature. For $`U=8t`$, we find that the onset temperature, $`T^{}=.021t`$. In the cuprates $`t=.5eV`$. Hence, we obtain an onset temperature of $`T^{}=100K`$ for pair formation. Phase coherence occurs at a lower temperature, $`T_c`$. As a consequence, the $`T^{}`$ calculated here should serve as a realistic estimate of the psuedogap temperature. At optimum doping, $`T^{}`$ should correspond to $`T_c`$. Experimentally, optimal doping corresponds to roughly 20%. At this doping level and for $`U=8t`$, we obtain that $`T_c=65K`$ leading to a ratio of $`T^{}/T_c=1.5`$. For YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.95</sub>, $`T_c=92K`$ while $`T^{}=110K`$. Typically in the cuprates, $`1.2<T^{}/T_c<2`$ in the underdoped regime. While our calculated values for $`T^{}`$ and $`T_c`$ are not rigorous estimamtes, it is encouraging that they are not far off from the experimental values.
## V Final Remarks
We have presented here an analysis ideally suited for the strong-coupling limit of the Hubbard model. Two key results are established. First, pairing in the Hubbard model occurs in the d$`_{x^2y^2}`$ channel. Second, composite excitations rather than electon-like excitations pair together. They live on a cluster of nearest-neighbour sites and are formed out of a hole and a singly-occupied site. Ultimately, such excitations could explain the absence of a well-defined electron-like peak in the ARPES experiments if the self-energy correction in Eq. (23) arising from the dynamical processes significantly broadens the energy levels near the Fermi energy. Sasaki, Matsumoto and Tachiki have evaluated such dynamical corrections in the context of the p-d model for the cuprates and found that the composite operator spectrum remains intact (thereby justifying the initial choice of the composite operator basis); however, broadening of all levels including those at the Fermi surface was observed. Matsumoto and Mancini have also performed similar calculations for the Hubbard model and observed a broadening of the levels at the Fermi surface.
Longer-range Coulomb interactions, necessitate the retension of higher-order composite excitations. The simplest of such excitations will involve three sites. Computing the equations of motion for the three-site composite excitations leads 4-site excitations. At each iteration of this procedure, the range of the composite excitations grows. We have verified that nearest-neighbour repulsive Coulomb interactions drastically diminish the anomalous pairing in $`\theta `$. Trivially, attractive nearest-neighbour Coulomb interactions enhance pairing. However, to completely settle the issue of pairing when repulsive next-nearest neighbour interactions are included, we must also retain the anomalous correlation functions that are generated in the presence of the longer-range interaction. Our work here suggests that such correlation functions should be computed to determine if composite particle pairing survives in the extended Hubbard model.
Nonetheless, within the on-site repulsive model, the correlation function ($`\theta `$) calculated here should be sufficient to describe the pairing process. Because $`\theta `$ involves a product of the form $`\eta _{ij}\eta _{ij}`$, the pairing mechanism is entirely local. From the form of $`\eta _{ij}`$, it is tempting to conclude that pairing requires a doubly-occupied site to neighbour a singly-occupied site. In such a configuration, double occupancy can be shared between the two sites. However, this is just one of the many types of local configurations that gives rise to a non-zero value of $`\theta `$. If long-range phase coherence obtains, one can think of the condensate as a coherent superposition of all such resonating structures. This suggests that the pair-pair correlation function for the cexons should be calculated to verify one way or another if phase coherence obtains.
Current Address: Theory Division, Los Alamos National Laboratory.
###### Acknowledgements.
We thank P. Wolynes, A. Yazdani, and T. Leggett for helpful conversations and the NSF grant No. DMR98-96134. |
warning/0001/hep-th0001137.html | ar5iv | text | # Massive chiral random matrix ensembles at 𝛽 = 1 & 4 : Finite-volume QCD partition functions
## A orthogonal ensemble
We use the identity
$`\mathrm{\Delta }(x_1,\mathrm{},x_N){\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{n}{}}}(x_jx_{N+i})`$ (39)
$`=`$ $`{\displaystyle \frac{\mathrm{\Delta }(x_1,\mathrm{},x_{N+n})}{\mathrm{\Delta }(x_{N+1},\mathrm{},x_{N+n})}}={\displaystyle \frac{\underset{1i,jN+n}{det}R_{i1}(x_j)}{\mathrm{\Delta }(x_{N+1},\mathrm{},x_{N+n})}},`$ (40)
where $`R_i(x)`$ is an arbitrary monic polynomial of the $`i`$-th order, and $`x_{N+i}m_i^20`$. We take $`\{R_i(x)\}`$ to be skew-orthogonal
$$R_{2i},R_{2j+1}_R=R_{2j+1},R_{2i}_R=h_i\delta _{ij},\text{others}=0,$$
(41)
with respect to the antisymmetric product
$`f,g_R`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑xx^{\frac{\nu 1}{2}}\mathrm{e}^xg(x){\displaystyle _0^x}𝑑yy^{\frac{\nu 1}{2}}\mathrm{e}^yf(y)`$ (43)
$`(fg).`$
When eq.(40) is integrated over $`x_1,\mathrm{},x_N`$ with the weight $`_{i=1}^N(\mathrm{e}^{x_i}x_i^{\frac{\nu 1}{2}})`$ in a cell $`0x_1x_2\mathrm{}x_N`$, it can be neatly expressed as a Pfaffian, due to the skew orthogonality (41) :
$`\mathrm{\Xi }_\nu ^{(1)}(\{m\})={\displaystyle \frac{\left(_{i=0}^{\frac{N+n}{2}1}h_i\right)\mathrm{Pf}(F)}{\mathrm{\Delta }(m_1^2,\mathrm{},m_n^2)}}(n:\text{even})`$ (45)
$`={\displaystyle \frac{\left(_{i=0}^{[\frac{N+n}{2}]1}h_i\right)\mathrm{Pf}\left(\begin{array}{cc}F& R\\ R^T& 0\end{array}\right)}{\mathrm{\Delta }(m_1^2,\mathrm{},m_n^2)}}(n:\text{odd}),`$ (48)
where
$`F^{ij}={\displaystyle \underset{k=0}{\overset{[\frac{N+n}{2}]1}{}}}{\displaystyle \frac{R_{2k}(m_i^2)R_{2k+1}(m_j^2)(ij)}{h_k}},`$
$`R^i=R_{N+n1}(m_i^2),1i,jn.`$
Explicit forms of the monic skew-orthogonal polynomials and their norms are known as :
$`R_{2k}(x)={\displaystyle \frac{(2k)!}{2^{2k+1}}}{\displaystyle \frac{d}{dx}}L_{2k+1}^{(\nu 1)}(2x),`$ (49)
$`R_{2k+1}(x)={\displaystyle \frac{(2k+1)!}{2^{2k+1}}}L_{2k+1}^{(\nu 1)}(2x)`$ (50)
$`{\displaystyle \frac{(2k)!}{2^{2k+2}}}(2k+\nu ){\displaystyle \frac{d}{dx}}L_{2k}^{(\nu 1)}(2x),`$ (51)
$`h_k=2^{4k\nu }(2k)!(2k+\nu )!.`$ (52)
In the microscopic limit (33) with $`\mu _i=2\sqrt{2N}m_i`$ fixed, the sum over the indices $`k`$ becomes an integral, and Laguerre polynomials approach modified Bessel functions:
$$L_k^{(\alpha )}(x)\left(\frac{k}{x}\right)^{\frac{\alpha }{2}}I_\alpha (2\sqrt{kx})\left(x=O(\frac{1}{k})<0\right).$$
(53)
Then the partition function is expressed as
$`Z_\nu ^{(1)}(\{\mu \})`$ $`=`$ $`\left({\displaystyle \underset{i=1}{\overset{n}{}}}\mu _i^\nu \right)\xi _\nu ^{(1)}(\{\mu \}),`$ (55)
$`\xi _\nu ^{(1)}(\{\mu \})`$ $`=`$ $`c_n{\displaystyle \frac{\mathrm{Pf}(f)}{\mathrm{\Delta }(\mu _1^2,\mathrm{},\mu _n^2)}}(n:\text{even})`$ (56)
$`=`$ $`c_n{\displaystyle \frac{\mathrm{Pf}\left(\begin{array}{cc}f& r\\ r^T& 0\end{array}\right)}{\mathrm{\Delta }(\mu _1^2,\mathrm{},\mu _n^2)}}(n:\text{odd}),`$ (59)
where
$`c_n`$ $`=`$ $`(1)^{\frac{n(n1)}{2}}2^{n^21}(n1)!{\displaystyle \underset{k=0}{\overset{n2}{}}}(2k+1)!(n:\text{even})`$
$`=`$ $`(1)^{\frac{n(n1)}{2}}2^{\frac{(n1)(7n+11)}{8}}(n1)!!{\displaystyle \underset{k=0}{\overset{n2}{}}}(2k+1)!`$
$`(n:\text{odd}),`$
$`f^{ij}`$ $`=`$ $`{\displaystyle _0^1}dtt^2{\displaystyle \frac{I_{\nu 1}(t\mu _i)}{\mu _i^{\nu 1}}}{\displaystyle \frac{I_\nu (t\mu _j)}{\mu _j^\nu }}(ij),`$
$`r^i`$ $`=`$ $`{\displaystyle \frac{I_\nu (\mu _i)}{\mu _i^\nu }},1i,jn.`$
The constant $`c_n`$ is conveniently determined as the above by requiring the small-$`\mu `$ behavior be in accord with that of the zero-dimensional $`\sigma `$ model,
$$Z_\nu ^{(1)}(\underset{n}{\underset{}{\mu ,\mathrm{},\mu }})\left(\frac{\mu }{2}\right)^{n\nu }\underset{k=0}{\overset{n1}{}}\frac{(2k)!}{(2k+\nu )!}(\mu 1).$$
(60)
## B symplectic ensemble
We concentrate on the case with an even $`n(2a)`$ number of flavors and pairwise degenerated mass parameters, corresponding to adjoint Dirac fermions in the QCD context.
We use the identity
$`\mathrm{\Delta }(x_1,\mathrm{},x_N)^4{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \underset{i=1}{\overset{a}{}}}(x_jx_{N+i})^2`$ (61)
$`={\displaystyle \frac{\left|\begin{array}{cccccc}\hfill 1& \hfill x_1& \hfill x_1^2& \hfill \mathrm{}& & x_1^{2N+a1}\\ \hfill 0& \hfill 1& \hfill 2x_1& \hfill \mathrm{}& \hfill (2N+a1)& x_1^{2N+a2}\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& & \mathrm{}\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& & \hfill \mathrm{}& \mathrm{}\\ \hfill 1& \hfill x_N& \hfill x_N^2& \hfill \mathrm{}& & x_N^{2N+a1}\\ \hfill 0& \hfill 1& \hfill 2x_N& \hfill \mathrm{}& \hfill (2N+a1)& x_N^{2N+a2}\\ \hfill 1& \hfill x_{N+1}& \hfill x_{N+1}^2& \hfill \mathrm{}& & x_{N+1}^{2N+a1}\\ \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& \hfill \mathrm{}& & \mathrm{}\\ \hfill 1& \hfill x_{N+a}& \hfill x_{N+a}^2& \hfill \mathrm{}& & x_{N+a}^{2N+a1}\end{array}\right|}{\mathrm{\Delta }(x_{N+1},\mathrm{},x_{N+a})}}`$ (71)
$`={\displaystyle \frac{\left|\begin{array}{cccc}Q_0(x_1)\hfill & Q_1(x_1)\hfill & \mathrm{}\hfill & Q_{2N+a1}(x_1)\hfill \\ Q_0^{}(x_1)\hfill & Q_1^{}(x_1)\hfill & \mathrm{}\hfill & Q_{2N+a1}^{}(x_1)\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & & \mathrm{}\hfill \\ Q_0(x_N)\hfill & Q_1(x_N)\hfill & \mathrm{}\hfill & Q_{2N+a1}(x_N)\hfill \\ Q_0^{}(x_N)\hfill & Q_1^{}(x_N)\hfill & \mathrm{}\hfill & Q_{2N+a1}^{}(x_N)\hfill \\ Q_0(x_{N+1})\hfill & Q_1(x_{N+1})\hfill & \mathrm{}\hfill & Q_{2N+a1}(x_{N+1})\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ Q_0(x_{N+a})\hfill & Q_1(x_{N+a})\hfill & \mathrm{}\hfill & Q_{2N+a1}(x_{N+a})\hfill \end{array}\right|}{\mathrm{\Delta }(x_{N+1},\mathrm{},x_{N+a})}},`$ (81)
where $`Q_i(x)`$ is an arbitrary monic polynomial of the $`i`$-th order, and $`x_{N+i}m_i^20`$. We take $`\{Q_i(x)\}`$ to be skew-orthogonal
$$Q_{2i},Q_{2j+1}_Q=Q_{2j+1},Q_{2i}_Q=h_i\delta _{ij},\text{others}=0,$$
(82)
with respect to the antisymmetric product
$`f,g_Q`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑xx^{2\nu +1}\mathrm{e}^{4x}(f(x)g^{}(x)f^{}(x)g(x)).`$ (83)
When eq.(81) is integrated over $`x_1,\mathrm{},x_N`$ with the weight $`_{i=1}^N(\mathrm{e}^{4x_i}x_i^{2\nu +1})`$, it can be neatly expressed as a Pfaffian, due to the skew orthogonality (82) :
$`\mathrm{\Xi }_\nu ^{(4)}(\{m\})={\displaystyle \frac{\left(_{i=0}^{N+\frac{a}{2}1}h_i\right)\mathrm{Pf}(F)}{\mathrm{\Delta }(m_1^2,\mathrm{},m_a^2)}}(a:\text{even})`$ (85)
$`={\displaystyle \frac{\left(_{i=1}^{N+[\frac{a}{2}]1}h_i\right)\mathrm{Pf}\left(\begin{array}{cc}F& Q\\ Q^T& 0\end{array}\right)}{\mathrm{\Delta }(m_1^2,\mathrm{},m_a^2)}}(a:\text{odd}),`$ (88)
where
$`F^{ij}={\displaystyle \underset{k=0}{\overset{N+[\frac{a}{2}]1}{}}}{\displaystyle \frac{Q_{2k}(m_i^2)Q_{2k+1}(m_j^2)(ij)}{h_k}},`$
$`Q^i=Q_{2N+a1}(m_i^2),1i,ja.`$
Explicit forms of the monic skew-orthogonal polynomials and their norms are known as :
$`Q_{2k}(x)={\displaystyle \frac{k!(k+\nu )!}{2^{2k}}}{\displaystyle \underset{l=0}{\overset{k}{}}}{\displaystyle \frac{(2l1)!!}{2^{2l}(l+\nu )!}}L_{2l}^{(2\nu )}(4x),`$ (89)
$`Q_{2k+1}(x)={\displaystyle \frac{(2k+1)!}{2^{4k+2}}}L_{2k+1}^{(2\nu )}(4x),`$ (90)
$`h_k=2^{8k4\nu 4}(2k+1)!(2k+2\nu +1)!.`$ (91)
In the microscopic limit with $`\mu _i=2\sqrt{2N}m_i`$ fixed, the sums over the indices $`k`$ and $`l`$ become integrals, and Laguerre polynomials approach modified Bessel functions. Then the partition function is expressed as
$`Z_\nu ^{(4)}(\{\mu \})`$ $`=`$ $`\left({\displaystyle \underset{i=1}{\overset{a}{}}}\mu _i^{2\nu }\right)\xi _\nu ^{(4)}(\{\mu \}),`$ (93)
$`\xi _\nu ^{(4)}(\{\mu \})`$ $`=`$ $`c_a{\displaystyle \frac{\mathrm{Pf}(f)}{\mathrm{\Delta }(\mu _1^2,\mathrm{},\mu _a^2)}}(a:\text{even})`$ (94)
$`=`$ $`c_a{\displaystyle \frac{\mathrm{Pf}\left(\begin{array}{cc}f& q\\ q^T& 0\end{array}\right)}{\mathrm{\Delta }(\mu _1^2,\mathrm{},\mu _a^2)}}(a:\text{odd}),`$ (97)
where
$`c_a`$ $`=`$ $`(1)^{\frac{a(a1)}{2}}{\displaystyle \underset{k=0}{\overset{a1}{}}}(2k+1)!`$
$`f^{ij}`$ $`=`$ $`{\displaystyle _0^1}dtt{\displaystyle \frac{I_{2\nu }(2t\mu _i)}{\mu _i^{2\nu }}}{\displaystyle _0^1}du{\displaystyle \frac{I_{2\nu }(2tu\mu _j)}{\mu _j^{2\nu }}}(ij),`$
$`q^i`$ $`=`$ $`{\displaystyle _0^1}𝑑t{\displaystyle \frac{I_{2\nu }(2t\mu _i)}{\mu _i^{2\nu }}},1i,ja.`$
The constant $`c_a`$ is conveniently determined as the above by requiring the small-$`\mu `$ behavior be in accord with that of the zero-dimensional $`\sigma `$ model,
$$Z_\nu ^{(4)}(\underset{2a}{\underset{}{\mu ,\mathrm{},\mu }})\mu ^{2a\nu }\underset{k=0}{\overset{a1}{}}\frac{(2k+1)!}{(2k+2\nu +1)!}(\mu 1).$$
(98)
It remains to confirm whether the above expressions for the $`\chi `$RME partition functions agree with those of the zero-dimensional $`\sigma `$ models. By taking all $`\mu `$’s to be identical, we obtain for $`\beta =1`$:
$`Z_\nu ^{(1)}(\underset{n}{\underset{}{\mu ,\mathrm{},\mu }})`$ $`=\stackrel{~}{c}_n{\displaystyle \frac{\mathrm{Pf}(\stackrel{~}{f})}{\mu ^{(n/21)n}}}(n:\text{even})`$ (103)
$`=\stackrel{~}{c}_n{\displaystyle \frac{\mathrm{Pf}\left(\begin{array}{cc}\stackrel{~}{f}& \stackrel{~}{r}\\ \stackrel{~}{r}^T& 0\end{array}\right)}{\mu ^{(n1)^2/2}}}(n:\text{odd}),`$
where
$`\stackrel{~}{c}_n`$ $`=`$ $`(1)^{\frac{n(n1)}{2}}2^{(n1)(\frac{n}{2}+1)}{\displaystyle \underset{k=0}{\overset{\frac{n}{2}1}{}}}{\displaystyle \frac{(2k+n1)!}{(2k)!}}(n:\text{even})`$
$`=`$ $`(1)^{\frac{n(n1)}{2}}2^{\frac{(n1)(3n+11)}{8}}{\displaystyle \frac{_{k=0}^{\frac{n3}{2}}(2k+n)!}{_{k=0}^{\frac{n5}{2}}(2k+1)!}}(n:\text{odd}),`$
$`\stackrel{~}{f}^{ij}`$ $`=`$ $`{\displaystyle _0^1}dtt^{i+j+2}I_{\nu +i1}(t\mu )I_{\nu +j}(t\mu )(ij),`$
$`\stackrel{~}{r}^i`$ $`=`$ $`I_{\nu +i}(\mu ),0i,jn1,`$
and for $`\beta =4`$:
$`Z_\nu ^{(4)}(\underset{2a}{\underset{}{\mu ,\mathrm{},\mu }})`$ $`=`$ $`\stackrel{~}{c}_a{\displaystyle \frac{\mathrm{Pf}(\stackrel{~}{f})}{\mu ^{(a1)a/2}}}(a:\text{even})`$ (105)
$`=`$ $`\stackrel{~}{c}_a{\displaystyle \frac{\mathrm{Pf}\left(\begin{array}{cc}\stackrel{~}{f}& \stackrel{~}{q}\\ \stackrel{~}{q}^T& 0\end{array}\right)}{\mu ^{(a1)a/2}}}(a:\text{odd}),`$ (108)
where
$`\stackrel{~}{c}_a`$ $`=`$ $`(1)^{\frac{a(a1)}{2}}{\displaystyle \underset{k=0}{\overset{a1}{}}}{\displaystyle \frac{(2k+1)!}{k!}},`$
$`\stackrel{~}{f}^{ij}`$ $`=`$ $`{\displaystyle _0^1}𝑑tt^{i+j+1}I_{2\nu +i}(2t\mu ){\displaystyle _0^1}𝑑uu^jI_{2\nu +j}(2tu\mu )`$
$`(ij),`$
$`\stackrel{~}{q}^i`$ $`=`$ $`{\displaystyle _0^1}𝑑tt^iI_{2\nu +i}(2t\mu ),0i,ja1.`$
We have numerically checked that, despite the appearances, the above expressions are identical to eqs.(22) and (24). Together with the $`\beta =2`$ case that has previously been confirmed , they explicitly show the equivalence between the $`\chi `$RMEs and the $`\sigma `$ models in Leutwyler-Smilga limit.
## C smallest eigenvalue distribution
The probability of finding no eigenvalue in the interval $`0x<s`$ is given by
$$E_\nu ^{(\beta )}(s;\{m\})=\frac{_s^{\mathrm{}}_{j=1}^N\left(dx_jw(x_j;\{m\})\right)|\mathrm{\Delta }(\{x\})|^\beta }{_0^{\mathrm{}}_{j=1}^N\left(dx_jw(x_j;\{m\})\right)|\mathrm{\Delta }(\{x\})|^\beta }.$$
(109)
The integral domain in the numerator can be traded to $`[0,\mathrm{})`$ with the weight function shifted by $`s`$, $`w(x+s;\{m\})`$. If the exponent in the weight function $`\frac{\beta }{2}(\nu +1)1`$ is an integer (excluding the case with $`\beta =1`$ and $`\nu `$ even), we can utilize the ‘flavor-topology duality’
$$\mathrm{\Xi }_\nu ^{(\beta )}(m_1,\mathrm{},m_n)=\mathrm{\Xi }_{\frac{2}{\beta }1}^{(\beta )}(m_1,\mathrm{},m_n,\underset{\frac{\beta }{2}(\nu +1)1}{\underset{}{0,\mathrm{},0}}),$$
(110)
to express $`E_\nu ^{(\beta )}(s;\{m\})`$ in terms of the partition functions:
$`E_\nu ^{(\beta )}(s;\{m\})=\mathrm{e}^{N\beta s}\times `$ (111)
$`{\displaystyle \frac{\mathrm{\Xi }_{\frac{2}{\beta }1}^{(\beta )}(\sqrt{s+m_1^2},\mathrm{},\sqrt{s+m_n^2},\stackrel{\frac{\beta }{2}(\nu +1)1}{\stackrel{}{\sqrt{s},\mathrm{},\sqrt{s}}})}{\mathrm{\Xi }_\nu ^{(\beta )}(m_1,\mathrm{},m_n)}}.`$ (112)
Now we change the picture back from Laguerre to chiral Gaussian, and take the microscopic limit with $`\zeta =\pi \rho (0)\sqrt{s}=2\sqrt{2Ns}`$ and $`\mu _i=2\sqrt{2N}m_i`$ fixed. The distribution of the smallest eigenvalue of chiral random matrices is then given by the first $`\zeta `$-derivative of $`E_\nu ^{(\beta )}`$:
$`P_\nu ^{(\beta )}(\zeta ;\{\mu \})={\displaystyle \frac{}{\zeta }}\{\mathrm{e}^{(\beta /8)\zeta ^2}\times `$ (113)
$`{\displaystyle \frac{\xi _{\frac{2}{\beta }1}^{(\beta )}(\sqrt{\zeta ^2+\mu _1^2},\mathrm{},\sqrt{\zeta ^2+\mu _n^2},.\stackrel{\frac{\beta }{2}(\nu +1)1}{\stackrel{}{\zeta ,\mathrm{},\zeta }})}{\xi _\nu ^{(\beta )}(\mu _1,\mathrm{},\mu _n)}}\}.`$ (114)
For $`\beta =1`$ and $`\nu `$ odd, eqs.(114) and (A) suffice to express the smallest eigenvalue distribution in a closed form. This prediction should in future be put in comparison with lattice QCD simulations with overlap dynamical quarks.
On the other hand, for $`\beta =4`$, the partition function in the numerator falls out of the range of this Letter, as the number of additional flavors is odd. A different formalism based on Fredholm determinant might be needed in order to overcome this limitation.
## D correlation function
In the case of even $`\beta `$, a $`p`$-level correlation function
$$\rho (\lambda _1,\mathrm{},\lambda _p;\{\mu \})=\underset{k=1}{\overset{p}{}}\mathrm{tr}\delta (\lambda _ki𝒟)$$
(115)
is expressed by construction as a ratio of partition functions with $`n`$ and $`n+\beta p`$ flavors . After taking the microscopic limit, In the $`\beta =4`$ case with $`a`$ pairs of doubly degenerated masses, it reads:
$$\rho _\nu ^{(4)}(\zeta _1,\mathrm{},\zeta _p;\{\mu \})=C_{a,\nu }^{(p)}\mathrm{\Delta }(\zeta _1^2,\mathrm{},\zeta _p^2)^4\underset{k=1}{\overset{p}{}}\left(\zeta _k^3\underset{i=1}{\overset{a}{}}(\zeta _k^2+\mu _i^2)^2\right)\frac{Z_\nu ^{(4)}(\mu _1,\mu _1,\mathrm{},\mu _a,\mu _a,\stackrel{4}{\stackrel{}{i\zeta _1,\mathrm{},i\zeta _1}},\mathrm{},\stackrel{4}{\stackrel{}{i\zeta _p,\mathrm{},i\zeta _p}})}{Z_\nu ^{(4)}(\mu _1,\mu _1,\mathrm{},\mu _a,\mu _a)}.$$
(116)
As our derivation of the partition function (B) is valid as well for negative values of $`m_i^2`$ or $`\mu _i^2`$, the above relationship suffices to express any $`p`$-level correlation function in a closed form, up to an overall constant independent of $`\mu `$’s.
SMN thanks P.H. Damgaard and E. Kanzieper for communications on various stages. This work was supported in part (SMN) by JSPS Research Fellowships for Young Scientists, and by Grant-in-Aid No. 411044 from the Ministry of Education, Science and Culture, Japan. |
warning/0001/astro-ph0001058.html | ar5iv | text | # Electron densities, temperatures and ionization rates in two interstellar clouds in front of 𝛽 Canis Majoris, as revealed by UV absorption lines observed with IMAPS
## 1 Introduction
The line of sight toward the B1 II-III star $`\beta `$ CMa in the direction ($`l_{II},b_{II}`$)=($`226\stackrel{}{.},14\stackrel{}{.}`$) has raised interest since its first study in the UV from Copernicus observations when Gry et al. (1985) showed that it contained very little neutral gas and that a great majority of the material on the line of sight was ionized. Indeed for a distance now set to 153 pc by Hipparcos measurements, Gry et al (1985) determined that $`N`$(H i) was somewhere within the range $`1\mathrm{2.2\; 10}^{18}`$cm<sup>-2</sup> based on the interstellar Ly$`\beta `$ absorption profile in the spectrum of $`\beta `$ CMa recorded by Copernicus, and showed that the total (neutral and ionized) hydrogen column density was ten times higher using the combined abundances of S ii and S iii compared to that of H i multiplied by the cosmic ratio of S to H. EUVE observations independently confirmed both findings. Cassinelli et al (1996) found N(H i)=$`2.0\pm \mathrm{0.2\; 10}^{18}`$cm<sup>-2</sup> by fitting the Lyman limit absorption in the EUV spectrum of the star. Evidence for ionization comes from the lower limits for the continuum absorption by He i in the spectrum of $`\beta `$ CMa which indicate that $`N(\text{He }\text{i})>\mathrm{6.\; 10}^{17}`$cm<sup>-2</sup> (Aufdenberg, et al. 1999) or $`>\mathrm{1.4\; 10}^{18}`$cm<sup>-2</sup> (Cassinelli, et al. 1996), depending on assumptions about the drop in stellar flux across the He i ionization edge. The fact that $`N(\text{He }\text{i})0.1N(\text{}\text{i})`$ indicates that a substantial fraction of the hydrogen toward $`\beta `$ CMa is ionized.
The weak presence of H i gas in this region of the sky is particularly evident in the neighbouring line of sight toward $`ϵ`$ CMa ($`l_{II}=239\stackrel{}{.}8`$, $`b_{II}=11\stackrel{}{.}3`$) where the Lyman limit absorption is even lower (Cassinelli et al, 1995). Gry et al (1995) found a H i column density upper limit of $`\mathrm{5\; 10}^{17}`$ and showed furthermore that almost all the neutral gas is very close to the Sun, distributed in two small components which are also detected with similar column densities in the line of sight to Sirius 2.7 pc away (Lallement et al.,1994). Because of the scarcity of H i in this region, $`ϵ`$ CMa and $`\beta `$ CMa are the two strongest sources of EUV radiation, and the two stars, especially $`ϵ`$ CMa, dominate the H-ionization field close to the Sun (Vallerga & Welsh 1995 ; Vallerga 1997). The conditions along these lines of sight, especially their ionization structure, strongly influence the nature and ionization of the local interstellar matter. Yet even with this in mind, we find that the ionization of the local ISM is not very well understood. For instance, with the radiation from $`ϵ`$ and $`\beta `$ CMa plus other easily identifiable sources of ionizing radiation, it is difficult to reconcile the fact that helium appears to be more ionized than hydrogen, as revealed by evidence from the EUVE spectra of several white dwarfs (Dupuis et al. 1995).
High resolution observations with GHRS on board HST allowed Dupin and Gry (1998) to show that the bulk of the gas in the $`\beta `$ CMa sight-line is distributed in two main components separated by 10 km/s, which are warm, only slightly depleted and both mostly ionized. However the analysis of the ionization processes was hampered by the fact that several species and especially the neutral species had been observed with a limited resolution of about 20 000 which did not allow for a clear separation of the two components. The relative distribution of matter in the two components therefore had uncertainties that allowed only lower limits to the ionization fractions to be derived.
Here we present new observations performed with the Interstellar Medium Absorption Profile Spectrograph (IMAPS) (Jenkins, et al. 1996), with the aim of gaining further insights on the ionization structure and various physical processes in the local interstellar medium. The higher resolution of 60 000 provided by IMAPS allowed us to resolve the absorptions of the two main components which in turn permitted the derivations of ionization fractions in each cloud. We also determined the clouds’ electron densities using the observation of C ii\*, which ultimately led to new insights on the ionization process.
## 2 Observations and data reduction
The observations were carried out by IMAPS when it was operated on the ORFEUS-SPAS II mission that flew in late 1996 (Hurwitz, et al. 1998). IMAPS is an objective-grating echelle spectrograph that was designed to record the spectra of bright, early-type stars over the wavelengths from $`950`$ Å to $`1150`$ Å with a high spectral resolution. For more details on the instrument see Jenkins et al. (1998, 1996).
The spectra were extracted from the echelle spectral images using special procedures developed by one of us (EBJ) and his collaborators on the IMAPS investigation team. A zero point in the wavelength calibration was made by measuring the O i* and O i\** telluric lines detected at 1040.943 Å and 1041.688 Å. Table 1 presents the interstellar lines detected in the data, with wavelengths and $`f`$-values taken from Morton (1991).
## 3 Analysis
The structure of velocity components in the sight-line toward $`\beta `$ CMa has already been determined by the analysis of GHRS HST data (Dupin & Gry 1998) : four components were detected, two of which (designated as Components C and D) were strongly dominant, and the remaining two (A and B) showed up only for the strongest lines. In the study conducted by Dupin & Gry (1998), only the lines between 1800 Å and 3000 Å (Fe ii, Mg ii, Mg i and Si ii lines) were observed at a resolution $`\lambda /\mathrm{\Delta }\lambda \mathrm{80\hspace{0.17em}000}`$. Because of the failure of the Echelle A grating at the time of the observations, all the species that had lines between 1150 Å and 1800 Å were observed at medium resolution ($`\mathrm{20\hspace{0.17em}000}`$). With a resolution of $`\mathrm{60\hspace{0.17em}000}`$, the IMAPS spectrograph now allows us to analyse lines of some important species like N i and O i at a resolution comparable to that of the GHRS Echelle data. As a consequence, we can resolve partially clouds C and D which are separated by only $`10`$km s<sup>-1</sup>. To facilitate a comparison between the IMAPS data and the HST data that were presented by Dupin & Gry (1998), we interpreted the column density vs. velocity profiles in terms of the components defined by them. However, the data are not precise enough to rule out other, slightly different combinations of components that could produce fits that are just as acceptable.
The relative velocities between the four clouds are precisely known from the Dupin & Gry (1998) study. The heliocentric radial velocities found in our study are larger than those of Dupin & Gry (1998) by about $`5.5`$km s<sup>-1</sup>. This difference is larger than the stated error for the GHRS velocity zero points, but we favor the velocities from the IMAPS data because we could use the telluric absorption lines as a reference. The velocity of C ii* is especially good since it is very close to the O i\** line at 1041.688Å in the spectral format of the IMAPS echelle and cross disperser gratings. The velocity dispersions $`b`$ for C ii, N i and O i in Components C and D are not well determined, since some blurring of the lines is caused by the instrument. However, they should be constrained by $`b`$(Mg ii) on the low end and the measured widths on the upper end. Thus, the principal free parameters are the column densities of various species in cloud C and cloud D. For the strong lines of N i and O i we see some evidence for the presence of Component B, and we can place upper limits for the strength of A.
The principal source of uncertainty for the column densities is the determination of the background intensity caused by scattered light from the echelle and cross disperser gratings in IMAPS. As is evident from Jenkins, et al. (1996 – see their Fig. 12) there is some overlap of energy from adjacent echelle orders, precluding the use of the interorder regions for measuring the background levels. We could determine the overall behavior of the scattered light by examining the trends for the bottoms of strong stellar absorption features. For the background level below the line of C ii* at 1037.018Å, we could use the nearby, saturated absorption from interstellar C ii at 1036.337Å.
Fig. 1 shows the features recorded by IMAPS after they were normalized to the stellar continuum. Overplotted on these data are the predicted absorption profiles for the column densities and central velocities given in Table 2. Selected results for Components C and D derived by Dupin & Gry (1998) are listed in Table 3, as an aid to following the interpretations that we present in §4. We made no attempt to measure the Si ii lines listed in Table 1, since they were saturated and not nearly as useful as the transition at 1808 Å observed by Dupin & Gry (1998). Conversely, our N i lines are not saturated, but the weakest member of the 1200 Å multiplet observed by Dupin & Gry had a central optical depth $`\tau _0=1.7`$ for Component C \[derived from our $`N`$(N i) and their $`b`$ value\]. The lack of saturation and higher resolution for the IMAPS recordings of N i absorption lead to our preferring the values derived here to those of Dupin & Gry, even though there is some uncertainty in the IMAPS background level.
The O i transition at 1302 Å observed by Dupin & Gry was hopelessly saturated, which led to their only being able to express a lower limit for $`N`$(O i). Fig. 1 shows that the IMAPS recording of the much weaker transition at 1039 Å shows a fairly strong saturation. If there are narrow, unresolved subcomponents within Component C that are far more saturated than what we see in the apparent residual intensity of the line, our listed column density may be below the true value. In principle, one could sense the presence of such components by observing that the N i absorptions do not grow as fast as their increases in transition strengths. However the lines that we observed here are too weak to show this effect well. Nevertheless, it is interesting to note that by observing the strong multiplet of N i at 1200 Å Dupin & Gry (1998) found $`N(\text{}\text{i})=8.1\pm \mathrm{0.8\; 10}^{13}`$cm<sup>-2</sup> for Component C (including the hidden local clouds)<sup>1</sup><sup>1</sup>1Note that in Dupin & Gry Table 2, the listed column densities for Component C had a contribution from the LIC and the other local cloud detected in both $`\alpha `$ CMa and $`ϵ`$ CMa sightlines subtracted off. This contribution is negligible for all species but O i and N i for which it represents about 25% of the total absorption in Component C., which is only 0.84 times the value that we obtained (see Table 2). (Our and their results for Component D agree however.) Thus, when lines reach a strength similar to that of the O i line at 1039 Å, the column densities might be understated by a factor of about 0.84.
## 4 Interpretation
### 4.1 Radial velocities
As noted earlier, there is a systematic offset between the velocities reported here and those given by Dupin & Gry (1988). Our new velocity for Component C at $`20`$km s<sup>-1</sup> makes it consistent with that of gas in the Local Interstellar Cloud (LIC), moving at a velocity of $`26`$km s<sup>-1</sup> towards the direction $`l=186^{}`$ and $`b=16^{}`$ (Lallement & Bertin 1992), whereas Dupin & Gry felt that the LIC was buried in the gap between Components C and D. However, there are two important problems for identifying the origin of Component C with the LIC. First, the column densities of Mg ii and Fe ii in Component C (see Table 3) are considerably larger than $`N(\text{Mg }\text{ii})=\mathrm{1.65\; 10}^{12}`$cm<sup>-2</sup> and $`N(\text{Fe }\text{ii})=\mathrm{8.5\; 10}^{11}`$cm<sup>-2</sup> found in the LIC component at $`18.8`$km s<sup>-1</sup> in the spectrum of $`\alpha `$ CMa (Lallement, et al. 1994) which is only $`5.5^{}`$ away in the sky, and than $`N(\text{Mg }\text{ii})=\mathrm{3.0\; 10}^{12}`$cm<sup>-2</sup> and $`N(\text{Fe }\text{ii})=\mathrm{1.35\; 10}^{12}`$cm<sup>-2</sup> found in the LIC component at $`17`$km s<sup>-1</sup> in the spectrum of $`ϵ`$ CMa (Gry, et al. 1995), $`17^{}`$ away in the sky. Also, $`N(\text{}\text{i})=\mathrm{5\; 10}^{17}`$cm<sup>-2</sup> in front of $`\alpha `$ CMa (Holberg, et al. 1998) is only about a third as large as the H i column density that we identify with Component C \[$`N(\text{}\text{i})=\mathrm{1.6\; 10}^{18}`$cm<sup>-2</sup>, taken from $`N`$(N i) and $`N`$(O i) after correcting for depletion – see §4.2 below.\] Even though $`\alpha `$ CMa is only 2.6 pc away, it seems to lie beyond the boundary of the LIC because Lallement, et al. (1994) detected another (somewhat weaker) component at $`v=13`$km s<sup>-1</sup>. Second, we show in §4.5 that the ionizing radiation field that is needed to sustain the ionization is considerably higher than that found for the local vicinity. Thus, the agreement of velocity of Component C with that of the LIC is probably just coincidental.
### 4.2 Column densities, depletions, and ionization fractions
We define a depletion factor $`\delta (X)`$ of an element $`X`$ in terms of its column density $`N(X)`$, relative to that of hydrogen $`N(\mathrm{H})`$, by the expression
$$\delta (X)=\left(\frac{N(X)}{N(\mathrm{H})}\right)\left(\frac{\mathrm{H}}{X}\right)_{\mathrm{cosmic}}$$
(1)
where $`(\mathrm{H}/X)_{\mathrm{cosmic}}`$ is the cosmic abundance ratio.<sup>2</sup><sup>2</sup>2Throughout this paper, we use the cosmic abundances of Anders & Grevesse (1989) For such a depletion to be meaningful in any given circumstance, one must account for differences in the observed ionization stage(s) of element $`X`$ relative to that adopted for H (either H i or H<sub>total</sub>). We address this consideration in the cases that follow, starting with the simplest ones involving the elements nitrogen, oxygen and sulphur.
The relative ionizations of N i and O i are closely coupled to that of H i through the large charge exchange reaction rates that arise from their nearly equal ionization potentials (Field & Steigman 1971; Butler & Dalgarno 1979). In the diffuse interstellar medium, these two elements are only mildly depleted below their cosmic abundances relative to hydrogen (Hibbert, et al. 1985; Cardelli, et al. 1991a,b; Meyer, et al. 1994, 1997, 1998), so they can serve as reasonably good indicators for H i in a partially ionized region. Sulphur is another element that has little or no depletion onto dust grains, but its second ionization potential is very high (23 eV). If the photoionization rate arising from photons with $`E>23`$ eV is not very large, we expect that S ii should be a good indicator of the total hydrogen in the line of sight, both in the neutral and ionized forms. In essence, for N, O and S we can make use of our general understanding of their depletions to arrive at conclusions on the relative ionization of hydrogen.
With our ability to identify how much of the absorptions by N i and O i can be assigned to Components C and D, we can differentiate between the relative fractions of H in neutral form in each case, after making comparisons with Dupin & Gry’s (1998) values for S ii in these same components. In doing this, we invoke two assumptions about the depletions: (1) the depletions of N and O onto dust grains are the same for Components C and D, and (2) S is undepleted. Assumption (1) seems reasonable in the light of evidence presented by Meyer, et al. (1997, 1998) and (2) seems justified from the evidence summarized by Savage & Sembach (1996) and Fitzpatrick & Spitzer (1997). If either of these assumptions is not quite correct, small errors in our conclusions may arise. However, they are probably not much worse than the uncertainty in identifying the relative contributions of Components C and D in the low resolution spectra of the S ii features.
For nitrogen, we obtain a general depletion factor $`\delta (\mathrm{N})=0.55`$ through Eq. 1 by taking $`N(\text{}\text{i})_{\mathrm{C}+\mathrm{D}}`$, dividing it by $`N(\text{}\text{i})=\mathrm{2.\; 10}^{18}`$cm<sup>-2</sup>, and then multiplying the result by the cosmic ratio of H to N. The same analysis for O leads to $`\delta (\mathrm{O})=0.58`$. These values, incidentally, are not much different from those shown by Meyer, et al. (1997, 1998). We now estimate the hydrogen neutral fractions in each component from the expression (for the case of N i)
$$n(\text{}\text{i})/n(\mathrm{H}_{\mathrm{total}})=\left(\frac{\mathrm{S}}{\mathrm{N}}\right)_{\mathrm{cosmic}}\frac{N(\text{}\text{i})}{\delta (\mathrm{N})N(\text{}\text{ii})}$$
(2)
and likewise for O i. For both N i and O i, we obtain $`n(\text{}\text{i})/n(\mathrm{H}_{\mathrm{total}})=0.25`$ for Component C and 0.035 for Component D.
The high column density of Si iii identified with Component D by Dupin & Gry (1998) presents a special problem (see Table 3). The rate constant for charge exchange between doubly ionized Si and neutral H is about $`\mathrm{3.0\; 10}^9`$cm<sup>3</sup> s<sup>-1</sup> at the temperatures of interest to us (Gargaud, et al. 1982). From the electron density derived later in §4.3 and $`n(\text{}\text{i})/n(\mathrm{H}_{\mathrm{total}})`$ for Component D, we find that $`n(\text{}\text{i})0.005`$cm<sup>-3</sup> (for an applicable temperature that we derive in §4.4). This density times the charge exchange rate constant is considerably larger than the expected ionization rate $`\mathrm{\Gamma }(\text{Si }\text{ii})=\mathrm{5.\; 10}^{14}\mathrm{s}^1`$ that we infer from the hydrogen ionization and the shapes of the EUV spectra of $`ϵ`$ and $`\beta `$ CMa (Vallerga & Welsh 1995; Aufdenberg, et al. 1999) which should dominate the photoionizing radiation field. It is hard to reconcile this inequality with Dupin & Gry’s finding that $`N(\text{Si }\text{iii})/N(\text{Si }\text{ii})2`$ for Component D. We propose that this contradiction can be overcome by having the Si iii exclusively within a different part of the region, perhaps one that is much closer to some ionizing source or, alternatively, within the cloud’s conduction/evaporation front at an interface with gaseous material at a much higher temperature. There is empirical evidence that the association of Si iii with lower ionization stages is not unusual. For instance, Cowie, et al. (1979) found that the velocity endpoints of O vi, Si iii, and N ii features were mutually correlated, and they used this evidence to suggest that the interfaces between cool and hot gas were conspicuous in the interstellar absorption lines.
In sections that follow, we will ignore the existence of the Si iii-bearing region and propose that the contradiction that we have noted justifies our regarding it as unrelated to the gas that holds most of the lower ionization states.
### 4.3 Electron densities from $`N`$(C ii\*)/$`N`$(C ii)
The relative populations of the fine-structure states of C ii are governed by the balance between collisions and the radiative decay of the upper level. If electrons are the dominant projectiles for excitation and de-excitation, the rate coefficient for de-excitations is
$$\gamma _{2,1}=\frac{\mathrm{8.63\; 10}^6\mathrm{\Omega }_{1,2}}{g_2T^{0.5}}\mathrm{cm}^3\mathrm{s}^1$$
(3)
(Spitzer 1978, p. 73), with the reverse rate given by detailed balancing, $`\gamma _{1,2}=(g_2/g_1)\mathrm{exp}(E_{1,2}/kT)\gamma _{2,1}`$. The statistical weights of the levels are $`g_1=2`$ and $`g_2=4`$, and the temperature equivalent for the difference in energy levels $`E_{1,2}/k=94.9`$K. Thus we find that the condition for equilibrium,
$$n(e)\gamma _{1,2}n(\text{}\text{ii})=[n(e)\gamma _{2,1}+A_{2,1}]n(\text{}\text{ii})$$
(4)
will lead to an equation for the electron density
$$n(e)=\frac{g_2A_{2,1}T^{0.5}\left[\frac{n(\text{}\text{ii})}{n(\text{}\text{ii})}\right]}{\mathrm{8.63\; 10}^6\mathrm{\Omega }_{1,2}\left\{\left(\frac{g_2}{g_1}\right)\mathrm{exp}\left(\frac{E_{1,2}}{kT}\right)\left[\frac{n(\text{}\text{ii})}{n(\text{}\text{ii})}\right]\right\}}$$
(5)
where the radiative decay probability for the upper level is $`A_{2,1}=\mathrm{2.29\; 10}^6\mathrm{s}^1`$ (Nussbaumer & Storey 1981), and the collision strength $`\mathrm{\Omega }_{1,2}=2.81`$ (Hayes & Nussbaumer 1984). Optical pumping of the C ii fine structure levels is unlikely to happen under the circumstances where there is a very high optical depth in the line (Sarazin, et al. 1979). Also, the density of pumping radiation must be very large – even larger than the elevated levels considered later in §4.5. Finally, unacceptably high values for $`n(\text{}\text{i})`$ would be needed for collisions by neutrals to have any importance (Keenan, et al. 1986).
As is clear from Fig. 1, the C ii feature at 1036 Å is far too saturated to allow a determination of $`N(\text{}\text{ii})`$ for either Components C or D. Thus, our determination of the column densities must be indirect. A good surrogate for carbon is sulphur. Our repeat of calculations of the type done by Sofia & Jenkins (1998) indicate that in partially ionized gases these two elements have about the same fraction of atoms elevated to higher (unseen) stages of ionization for a wide range of conditions. The depletion of carbon atoms in dense clouds is typically $`\delta (\mathrm{C})=0.39`$ (Cardelli, et al. 1993, 1996; Sofia, et al. 1998), and, learning from the example of $`\tau `$ CMa (Sofia, et al. 1997), this level of depletion seems to hold also for low density lines of sight. On the assumption that carbon toward $`\beta `$ CMa is depleted by this factor and there is no depletion of sulphur, we can arrive at $`N(\text{}\text{ii})`$ from the product $`(\mathrm{C}/\mathrm{S})_{\mathrm{cosmic}}\delta (\mathrm{C})N(\text{}\text{ii})`$. Doing so gives us the values $`N(\text{}\text{ii})=\mathrm{8.3\; 10}^{14}`$cm<sup>-2</sup> for Component C and $`\mathrm{1.7\; 10}^{15}`$cm<sup>-2</sup> for Component D, leading to $`N(\text{}\text{ii})/N(\text{}\text{ii})=0.029_{0.012}^{+0.020}`$ and $`0.0071_{0.0023}^{+0.0034}`$ for Components C and D, respectively.
Before we can derive values for $`n(e)`$ from Eq. 5, we must determine $`T`$. To do this, we rely on another method for measuring $`n(e)`$, but one that has a different temperature dependence. The balance between the ionization of Mg i and the recombination of Mg ii is a fundamentally different process from that which governs the fine-structure equilibrium of C ii, but both are driven by the value for $`n(e)`$. In the next section (§4.4), we shall make use of the difference to constrain other free parameter, the temperature $`T`$.
### 4.4 Electron densities from Mg i/Mg ii
The equation for the equilibrium of the lowest 2 ionization levels of Mg is given by
$`\left[\mathrm{\Gamma }(\text{Mg }\text{i})+C(\text{Mg }\text{ii})n(H^+)\right]n(\text{Mg }\text{i})=`$
$`\alpha (\text{Mg }\text{i})n(e)n(\text{Mg }\text{ii})`$ (6)
For the charge exchange rate $`C(\text{Mg }\text{ii})`$ that applies to the reaction Mg i \+ H$`{}_{}{}^{+}`$ Mg ii \+ H, we used the analytical approximation $`C(\text{Mg }\text{ii})=\mathrm{1.74\; 10}^9\mathrm{exp}(\mathrm{2.21\; 10}^4/T)`$ derived by Allan, et al. (1988). Collisional ionization of Mg i is negligible at the temperatures of interest in this study. We computed $`\mathrm{\Gamma }(\text{Mg }\text{i})`$ at the position of the Sun using mean of the stellar radiation fields estimated by Jura (1974) $`\lambda F_\lambda =\mathrm{3.0\; 10}^3\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ and Mathis, et al. (1983) $`\lambda F_\lambda =\mathrm{2.5\; 10}^3\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, multiplied by the photoionization cross section given by Verner, et al. (1996), to arrive at $`\mathrm{\Gamma }(\text{Mg }\text{i})=\mathrm{6.1\; 10}^{11}\mathrm{s}^1`$. (However, this number increases to $`\mathrm{9.5\; 10}^{11}\mathrm{s}^1`$ when we consider locations closer to $`ϵ`$ and $`\beta `$ CMa at a later stage of the analysis.) For $`\alpha (\text{Mg }\text{i})`$ we used the radiative and dielectronic recombination rates given by Shull & van Steenberg (1982), supplemented by the additional contributions from low-lying resonance states computed by Nussbaumer & Storey (1986).
It is immediately evident that the Echelle B observations of Mg i and Mg ii by Dupin & Gry (1998) both show an excellent signal-to-noise ratio, a smooth stellar continuum flux, and absorption features that are reasonably well resolved. Nevertheless, there is a difference between the two. On the one hand, the 2853 Å feature of Mg i represents an absorption that is unsaturated and thus straightforward to analyze. On the other hand, the two features of Mg ii (2796 and 2803 Å) were heavily saturated, which leads to a large uncertainty in $`N(\text{Mg }\text{ii})`$. Even though Dupin & Gry could use the well-determined $`b`$ value for Mg i to help in the interpretation of the Mg ii features, the error in $`N`$(Mg ii) is still significant (see Table 3). We will use here a less direct method of deriving $`N`$(Mg ii) from similar but more accurately determined species, along with an application of some empirical relationships for interstellar depletions, that will ultimately lead to a completely independent estimate for $`N`$(Mg ii).
As with C and S discussed in §4.3, Mg and Si are a good pair of elements that have very similar photoionization properties (although the rate coefficients for charge exchange of their doubly ionized forms with neutral hydrogen are very different – see below). Fitzpatrick (1997) shows a general relationship between the logarithms of the interstellar depletions of Mg and Si. If we assume that S is generally undepleted, we can substitute it for H and then derive the quantity $`\mathrm{log}\delta (\mathrm{Mg})`$ from our observed $`\mathrm{log}\delta (\mathrm{Si})`$ and his best fit to the trend $`\mathrm{log}\delta (\mathrm{Mg})=0.82\mathrm{log}\delta (\mathrm{Si})0.17`$ dex. When we do this, we find that $`\mathrm{log}\delta (\mathrm{Mg})=0.63`$ dex for Component C and $`0.80`$ dex for Component D. Products of these derived $`\delta `$’s, the S ii column densities, and the cosmic values for Mg relative to S give $`N(\text{Mg }\text{ii})=\mathrm{5.4\; 10}^{13}`$cm<sup>-2</sup> and $`\mathrm{7.6\; 10}^{13}`$cm<sup>-2</sup> for Components C and D, respectively. A moderate adjustment in $`N(\text{Mg }\text{ii})`$ for Component D must be made to account for the fact that about 24% of the Mg is probably in the doubly ionized form, while the fractional amount of Si ii is much lower because it has a significantly larger rate coefficient for recombining via charge exchange with neutral hydrogen. This conclusion was derived quantitatively after making preliminary calculations of the ionization balance, as described later in §4.5. No such correction is needed for Component C. The estimate we made here for $`N(\text{Mg }\text{ii})`$ for Component C is somewhat larger than the one given in Table 3 – however is marginally consistent with it when the error of 0.16 dex attached to our present method is considered– while our (adjusted) value of $`N(\text{Mg }\text{ii})=\mathrm{5.8\; 10}^{13}`$cm<sup>-2</sup> is remarkably close to the value derived by Dupin & Gry. Finally, we arrive at $`N(\text{Mg }\text{i})/N(\text{Mg }\text{ii})`$ equal to $`0.0015_{0.0005}^{+0.0007}`$ for Component C and $`0.0043_{0.0013}^{+0.0020}`$ for Component D. Note that while Dupin & Gry’s values for $`N`$(S ii) have large estimated relative errors (36%, 22%), such errors have a very weak influence in the conclusions because the slope of the empirical relationship between $`\mathrm{log}\delta (\mathrm{Mg})`$ and $`\mathrm{log}\delta (\mathrm{Si})`$ is not much different than 1.0. The same applies to errors in the assumption that $`\mathrm{log}\delta (\mathrm{S})=0.0`$ dex.
As a check that the depletions toward $`\beta `$ CMa are not out of the ordinary, we carried out a similar exercise, comparing the measured $`\mathrm{log}\delta (\mathrm{Fe})=(1.20,1.60)`$ dex and their respective $`\mathrm{log}\delta (\mathrm{Si})=(0.56,0.77)`$ dex for consistency with the trend shown by Fitzpatrick (1996). Both cases fell within the body of points that defined this trend.
The curved bands in Fig. 2 show, as a function of temperature $`T`$, the expected outcomes for $`\mathrm{log}N(\text{Mg }\text{i})\mathrm{log}N(\text{Mg }\text{ii})`$ found by solving the ionization equilibrium equation (Eq. 4.4) with the values of $`n(e)`$ obtained from Eq. 5 and the measured values of $`N`$(C ii\*)/$`N`$(C ii) given in §4.3. From the intersections of these curves with the logarithms of the observed $`N(\text{Mg }\text{i})/N(\text{Mg }\text{ii})`$ presented above (straight, horizontal bands in the figure), we find that $`400<T<6500`$K for Component C and $`8000<T<\mathrm{14\hspace{0.17em}000}`$K for Component D. To compute the upper limits for $`n(e)`$, we evaluate Eq. 5 with the upper limit for $`N(\text{}\text{ii})/N(\text{}\text{ii})`$ at the higest temperature that is consistent with this upper limit (not the highest temperature allowed in general). For the lower limits, we take the lowest permissible $`N(\text{}\text{ii})/N(\text{}\text{ii})`$ and apply it to Eq. 5 with the lowest temperature. These conditions for the limits are shown by solid dots in Fig. 2. For Component C we find that $`0.08<n(e)<0.6`$cm<sup>-3</sup>, and for Component D we arrive at $`0.09<n(e)<0.2`$cm<sup>-3</sup>. (While our lower limit for $`n(e)`$ for Component C is formally allowed, the real value of $`n(e)`$ is probably much closer to the upper limit because the temperature of the gas is probably much higher than $`400`$K.) These represent the worst possible extremes in $`n(e)`$ permitted by the data.
### 4.5 Ionizing radiation
Shortward of the Lyman limit, the general ionizing radiation field near the Sun is dominated by the output of $`ϵ`$ and $`\beta `$ CMa supplemented by many nearby white dwarf stars, cataclysmic variables, and the line emission from some late-type stars (Vallerga 1997). In addition to stellar sources, some diffuse energetic radiation from the hot gas that surrounds the LIC may contribute additional ionization (Cheng & Bruhweiler 1990). The intensity of this diffuse radiation is rather uncertain: Jelinsky, et al. (1995) have observations that place a meaningful upper limit for this radiation if it produces spectral features in agreement with some specific predictions of hot-gas emission models.
If we calculate the hydrogen ionization produced by Vallerga’s (1997) composite EUV stellar radiation field, supplemented by Cheng & Bruhweiler’s (1990) hot gas radiation attenuated by absorption from a hydrogen column $`N(\text{}\text{i})=\mathrm{2\; 10}^{17}`$cm<sup>-2</sup>, we find that the neutral fraction of hydrogen is 0.89 and 0.47, for Components C and D, respectively, if we require that $`n(e)=0.31`$cm<sup>-3</sup> and $`0.13`$cm<sup>-3</sup> as we found from the analysis of $`N(\text{}\text{ii})/N(\text{}\text{ii})`$ and $`N(\text{Mg }\text{i})/N(\text{Mg }\text{ii})`$ in §§ 4.3 and 4.4. This is clearly inconsistent with our finding $`N(\text{}\text{i})/N(\mathrm{H}_{\mathrm{total}})=0.25`$ and 0.035 for the two components in §4.2. The model calculations followed the scheme outlined by Sofia & Jenkins (1998) that included the ionization of helium and allowed for the extra electrons coming from He. (In contrast to the LIC, where helium ionization is important, the models indicated that the high values of $`n(e)`$ effectively suppressed the ionization of helium.)
To overcome the fact that the predicted ionization of H is less than observed, we must find additional sources of ionization. One possibility is that there is vestigial ionization left over from a previous event that heated the gas and collisionally ionized it (Lyu & Bruhweiler 1996), such as the passage of a shock front in less a few times $`10^5`$yr ago. If this were the case, it would be difficult to explain why the gas is not moving at a substantial velocity. Another tactic is to propose that most of the gas that we see is rather close to $`ϵ`$ and $`\beta `$ CMa, where the hydrogen (and helium) ionization rates must be much higher than in the local vicinity.
We experimented with ionization models that allowed for the enhancement of $`\mathrm{\Gamma }(\text{}\text{i})`$ by placing the clouds closer to $`\beta `$ CMa. For one possible location, at the nearest point to $`\alpha `$ CMa with its white dwarf companion, the increase in the radiation field is not much: the strength of the radiation from Sirius B (Holberg, et al. 1998) reaches a maximum of only 0.4 times that of $`\beta `$ CMa as seen near the Sun. However this applies only to fluxes observed at wavelengths above the Lyman limit. If there were substantially less H i and He i between Sirius B and the cloud in front of $`\beta `$ CMa, the ionizing fluxes could be larger.
We expect that at a distance of 130 pc from the Sun, a cloud is at the smallest possible distance from $`ϵ`$ CMa (33 pc) and only 26 pc from $`\beta `$ CMa (assuming that the most probable distances between us and the stars derived from Hipparcos are correct). Under this condition, the flux of $`ϵ`$ CMa is enhanced by a factor of 17, and the flux from $`\beta `$ CMa increases by a factor of 36. With the spectral distributions given by Vallerga & Welsh (1995) for $`ϵ`$ CMa and Aufdenberg, et al. (1999) for $`\beta `$ CMa (after correction for absorption by the H and He between these stars and the Sun), we find that the enhanced fluxes at this position produce hydrogen neutral fractions of 0.25 and 0.037 for Components C and D. These values are close to what we observed and reported in §4.2.
To some degree, our results are dependent on the assumed $`N`$(H i) and $`N`$(He i) between the ionizing sources and most of the gas in the components. Absorption in front of $`\beta `$ CMa is constrained by the column densities that we observed, whereas those for $`ϵ`$ CMa are arbitrary. For Component C our adopted columns were $`N(\text{}\text{i})=\mathrm{1.1\; 10}^{18}`$cm<sup>-2</sup> for $`\beta `$ CMa and $`\mathrm{1.1\; 10}^{17}`$cm<sup>-2</sup> for $`ϵ`$ CMa (this low value was required to raise the ionization high enough). For Component D, we made the columns in front of both $`ϵ`$ CMa and $`\beta `$ CMa equal to $`N(\text{}\text{i})=\mathrm{2\; 10}^{17}`$cm<sup>-2</sup>. For both components, we made the He i to H i ratio in the absorbing gas consistent with what we found from the ionization equilibrium calculations.
Our ionization equilibrium calculations indicate that both clouds are located in a region close to $`ϵ`$ and $`\beta `$ CMa where the radiation fields are strongly enhanced. Note also that a physically separate region, perhaps an outer layer of the cloud representing Component D that has evolved to a more fully ionized state, is revealed by the presence of Si iii.
The fractional ionization of He that we compute is not very substantial. Singly ionized He accounts for only 6% (Component C) and 19% (Component D) of the total He, even when the material is placed at the position where the radiation fields from $`ϵ`$ and $`\beta `$ CMa are enhanced by the large factors given above. The total column density $`N(\text{He }\text{i})_{\mathrm{C}+\mathrm{D}}=\mathrm{1.6\; 10}^{18}`$cm<sup>-2</sup>, which is above the lower limits stated in §4.2. For Component C, the computed He i/H i of 0.38 agrees with the general results of Wolff, et al. (1999), while the much higher ionization of H in Component D produces a considerably larger He i/H i=2.2.
###### Acknowledgements.
The observations reported in this paper are from a guest investigator program that was approved by NASA as a part of the US share of observing time on the ORFEUS-SPAS II flight in 1996, a joint undertaking of the US and German space agencies, NASA and DARA. The successful execution of our observations was the product of efforts over many years by engineering teams at Princeton University Observatory, Ball Aerospace Systems Group (the industrial subcontractor for the IMAPS instrument) and Daimler-Benz Aerospace (the German firm that built the ASTRO-SPAS spacecraft and conducted mission operations). Contributions to the success of IMAPS also came from the generous efforts by many members of the Optics Branch of the NASA Goddard Space Flight Center (grating coatings and testing) and from O. H. W. Siegmund and S. R. Jelinsky at the Berkeley Space Sciences Laboratory (deposition of the photocathode material). We are grateful to Martin Lemoine for providing his absorption line fitting software that we used to derive the component column densities. This research was supported by NASA grant NAG5$``$616 to Princeton University. |
warning/0001/nlin0001034.html | ar5iv | text | # Chaotic Field Theory: a Sketch
## 1 Unstable recurrent patterns in classical field theories
Field theories such as 4-dimensional QCD or gravity have many dimensions, symmetries, tensorial indices. They are far too complicated for exploratory forays into this forbidding terrain. We start instead by taking a simple spatio-temporally chaotic nonlinear system of physical interest, and investigate the nature of its solutions.
One of the simplest and extensively studied spatially extended dynamical systems is the Kuramoto-Sivashinsky system
$$u_t=(u^2)_xu_{xx}\nu u_{xxxx}$$
(1)
which arises as an amplitude equation for interfacial instabilities in a variety of contexts. The “flame front” $`u(x,t)`$ has compact support, with $`x[0,2\pi ]`$ a periodic space coordinate. The $`u^2`$ term makes this a nonlinear system, $`t`$ is the time, and $`\nu `$ is a fourth-order “viscosity” damping parameter that irons out any sharp features. Numerical simulations demonstrate that as the viscosity decreases (or the size of the system increases), the “flame front” becomes increasingly unstable and turbulent. The task of the theory is to describe this spatio-temporal turbulence and yield quantitative predictions for its measurable consequences.
Armed with a computer and a great deal of skill, one can obtain a numerical solution to a nonlinear PDE. The real question is; once a solution is found, what is to be done with it? The periodic orbit theory is an answer to this question.
Dynamics drives a given spatially extended system through a repertoire of unstable patterns; as we watch a “turbulent” system evolve, every so often we catch a glimpse of a familiar pattern:
$``$ other swirls $``$
For any finite spatial resolution, the system follows approximately for a finite time a pattern belonging to a finite alphabet of admissible patterns, and the long term dynamics can be thought of as a walk through the space of such patterns, just as chaotic dynamics with a low dimensional attractor can be thought of as a succession of nearly periodic (but unstable) motions. The periodic orbit provides the machinery that converts this intuitive picture into precise calculation scheme that extracts asymptotic time predictions from the short time dynamics. For extended systems the theory gives a description of the asymptotics of partial differential equations in terms of recurrent spatio-temporal patterns.
Putkaradze has proposed that the Kuramoto-Sivashinsky system (1) be used as a laboratory for exploring such ideas. We now summarize the results obtained so far in this direction by Christiansen et al. and Zoldi and Greenside .
The solution $`u(x,t)=u(x+2\pi ,t)`$ is periodic on the $`x[0,2\pi ]`$ interval, so one (but by no means only) way to solve such equations is to expand $`u(x,t)`$ in a discrete spatial Fourier series
$$u(x,t)=i\underset{k=\mathrm{}}{\overset{+\mathrm{}}{}}a_k(t)e^{ikx}.$$
(2)
Restrict the consideration to the subspace of odd solutions $`u(x,t)=u(x,t)`$ for which $`a_k`$ are real. Substitution of (2) into (1) yields the infinite ladder of evolution equations for the Fourier coefficients $`a_k`$:
$$\dot{a}_k=(k^2\nu k^4)a_kk\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}a_ma_{km}.$$
(3)
$`u(x,t)=0`$ is a fixed point of (1), with the $`k^2\nu <1`$ long wavelength modes of this fixed point linearly unstable, and the short wavelength modes stable. For $`\nu >1`$, $`u(x,t)=0`$ is the globally attractive stable fixed point; starting with $`\nu =1`$ the solutions go through a rich sequence of bifurcations, and myriad unstable periodic solutions whose number grows exponentially with time.
The essential limitation on the numerical studies undertaken so far have been computational constraints: in truncation of high modes in the expansion (3), sufficiently many have to be retained to ensure the dynamics is accurately represented. Christiansen et al. have examined the dynamics for values of the damping parameter close to the onset of chaos, while Zoldi and Greenside have explored somewhat more turbulent values of $`\nu `$. With improvement of numerical codes considerably more turbulent regimes should become accessible.
One pleasant surprise is that even though one is dealing with (infinite dimensional) PDEs, for these strong dissipation values of parameters the spatio-temporal chaos is sufficiently weak that the flow can be visualised as an approximately 1-dimensional Poincaré return map $`sf(s)`$ from the unstable manifold of the shortest periodic point onto its neighborhood, see figure 1(a). This representation makes it possible to systematically determine all nearby periodic solutions up to a given maximal period.
So far some 1,000 prime cycles have been determined numerically for various values of viscosity. The rapid contraction in the nonleading eigendirections is illustrated in figure 1(b) by the plot of the first 16 eigenvalues of the $`\overline{1}`$-cycle. As the length of the orbit increases, the magnitude of contracting eigenvalues falls off very quickly. In figure 2 we plot $`u_0(x,t)`$ corresponding to the $`\overline{0}`$-cycle. The difference between this solution and the other shortest period solution is of the order of 50% of a typical variation in the amplitude of $`u(x,t)`$, so the chaotic dynamics is already exploring a sizable swath in the space of possible patterns even so close to the onset of spatio-temporal chaos. Other solutions, plotted in the configuration space, exhibit the same overall gross structure. Together they form the repertoire of the recurrent spatio-temporal patterns that is being explored by the turbulent dynamics.
## 2 Periodic orbit theory
Now we turn to the central issue; qualitatively, these solutions demonstrate that the recurrent patterns program can be implemented, but how is this information to be used quantitatively? This is what the periodic orbit theory is about; it offers the machinery that assembles the topological and the quantitative information about individual solutions into accurate predictions about measurable global averages, such as the Lyapunov exponents and correlation functions.
Very briefly (for a detailed exposition the reader is referred to ref. ), the task of any theory that aspires to be a theory of chaotic, turbulent systems is is to predict the value of an “observable” $`a`$ from the spatial and time averages evaluated along dynamical trajectories $`x(t)`$
$$a=\underset{t\mathrm{}}{lim}\frac{1}{t}A^t,A^t(x)=_0^t\tau a(x(\tau )).$$
The key idea of the periodic orbit theory is to extract this average from the leading eigenvalue of the evolution operator
$$^t(x,y)=\delta (yx(t))\mathrm{}^{\beta A^t(x)}$$
via the trace formula
$$\mathrm{tr}^t=\underset{p}{}\text{}_p=\underset{p}{}\underset{r=1}{\overset{\mathrm{}}{}}\frac{T_p\delta \left(trT_p\right)}{\left|\text{det}\left(\mathrm{𝟏}𝐉_p^r\right)\right|}\mathrm{}^{r\beta A_p}$$
(4)
which relates the spectrum of the evolution operator to a sum over prime periodic solutions $`p`$ of the dynamical system and their repeats $`r`$.
What does this formula mean? Prime cycles partition the dynamical space into neighborhoods, each cycle enclosed by a tube whose volume is the product of its length $`T_p`$ and its thickness $`\left|\text{det}(\mathrm{𝟏}𝐉_p)\right|^1`$. The trace picks up a periodic orbit contribution only when the time $`t`$ equals a prime period or its repeat, a constraint enforced here by $`\delta \left(trT_p\right)`$. $`𝐉_p`$ is the linear stability of cycle $`p`$, so for long cycles $`\left|\text{det}\left(\mathrm{𝟏}𝐉_p^r\right)\right|`$ (product of expanding eigenvalues), and the contribution of long and very unstable cycles are exponentially small compared to the short cycles which dominate trace formulas. The number of contracting directions and the overall dimension of the dynamical space is immaterial; that is why the theory can also be applied to PDEs. All this information is purely geometric, intrinsic to the flow, coordinate reparametrization invariant, and the same for any average one might wish to compute. The information related to a specific observable is carried by the weight $`\mathrm{}^{\beta A_p}`$, the periodic orbit estimate of the contribution of $`\mathrm{}^{\beta A^t(x)}`$ from the $`p`$-cycle neighborhood.
The intuitive meaning of a trace formula is that it expresses the average $`\mathrm{}^{\beta A^t}`$ as a discretized integral
$`\genfrac{}{}{0pt}{}{\text{smooth}}{\text{dynamics}}\text{}\genfrac{}{}{0pt}{}{\text{linearized}}{\text{neighborhoods}}`$
over the dynamical space partitioned topologically into a repertoire of spatio-temporal patterns, each weighted by the likelihood of pattern’s occurrence in the long time evolution of the system.
Periodic solutions are important because they form the skeleton of the invariant set of the long time dynamics, with cycles ordered hierarchically; short cycles give good approximations to the invariant set, longer cycles refinements. Errors due to neglecting long cycles can be bounded, and for nice hyperbolic systems they fall off exponentially or even super-exponentially with the cutoff cycle length . Short cycles can be accurately determined and global averages (such as Lyapunov exponents and escape rates) can be computed from short cycles by means of cycle expansions.
The Kuramoto-Sivashinsky periodic orbit calculations of Lyapunov exponents and escape rates demonstrate that the periodic orbit theory predicts observable averages for deterministic but classically chaotic spatio-temporal systems. The main problem today is not how to compute such averages — periodic orbit theory as well as direct numerical simulations can handle that — but rather that there is no consensus on what the sensible experimental observables worth are predicting.
It should be obvious, and it still needs to be said: the spatio-temporally periodic solutions are not to be thought of as eigenmodes, a good linear basis for expressing solutions of the equations of motion. Something like a dilute instant approximation makes no sense at all for strongly nonlinear systems that we are considering here. As the equations are nonlinear, the periodic solutions are in no sense additive, and their linear superpositions are not solutions.
$`\genfrac{}{}{0pt}{}{A\text{ }\text{}\text{ }+B\text{ }\text{}\text{ }+\mathrm{}\text{ }\text{u(x,t)}}{\text{ a solution a solution not a solution}}`$
Instead, it is the trace formulas and spectral determinants of the periodic orbit theory that prescribe how the repertoire of admissible spatio-temporal patterns is to be systematically explored, and how these solutions are to be put together in order to predict measurable observables.
Suppose that the above program is successfully carried out for classical solutions of some field theory. What are we to make of this information if we are interested in the quantum behavior of the system? In the semiclassical quantization the classical solutions are the starting approximation.
## 3 Stochastic evolution
For the same pragmatic reasons that we found it profitable to shy away from facing the 4-dimensional QCD head on in the above exploratory foray into a strongly nonlinear field theory, we shall start out by trying to understand the structure of perturbative corrections for systems radically simpler than a full-fledged quantum field theory. First, instead of perturbative corrections to the quantum problem, we shall start by exploring the perturbative corrections to weakly stochastic flows. Second, instead of continuous time flows, we shall start by a study of a discrete time process.
For discrete time dynamics a Langevin trajectory in presence of additive noise is generated by iteration
$$x_{n+1}=f(x_n)+\sigma \xi _n,$$
(5)
where $`f(x)`$ is a map, $`\xi _n`$ a random variable, and $`\sigma `$ parametrizes the noise strength. In what follows we assume that $`\xi _n`$ are uncorrelated, and that the mapping $`f(x)`$ is one-dimensional and expanding, but we expect that the form of the results will remain the same for higher dimensions, including the field theory example of the preceding section.
Tracking an individual noisy trajectory does not make much sense; what makes sense is the Fokker-Planck formulation, where one considers instead evolution of an ensemble of trajectories. An initial density of trajectories $`\varphi _0(x)`$ evolves with time as
$$\varphi _{n+1}(y)=\left(\varphi _n\right)(y)=𝑑x(y,x)\varphi _n(x)$$
(6)
where $``$ is the evolution operator
$$(y,x)=\delta (yf(x)\sigma \xi )P(\xi )𝑑\xi =\sigma ^1P\left[\sigma ^1(yf(x))\right],$$
(7)
and $`\xi _n`$ a random variable with the normalized distribution $`P(\xi )`$, centered on $`\xi =0`$.
If the noise is weak, the goal of the theory is to compute the perturbative corrections to the eigenvalues $`\nu `$ of $``$ order by order in the noise strength $`\sigma `$,
$$\nu (\sigma )=\underset{m=0}{\overset{\mathrm{}}{}}\nu ^{(m)}\frac{\sigma ^m}{m!}.$$
One way to get at the spectrum of $``$ is to consider the discrete Laplace transform of $`^n`$, or the resolvent
$$\underset{n=1}{\overset{\mathrm{}}{}}z^n\mathrm{tr}^n=\mathrm{tr}\frac{z}{1z}=\underset{\alpha =0}{\overset{\mathrm{}}{}}\frac{z\nu _\alpha }{1z\nu _\alpha }$$
(8)
which has a pole at every $`z=\nu _\alpha ^1`$.
The effects of weak noise are of interest in their own right, as any deterministic evolution that occurs in nature is affected by noise. However, what is most important in the present context is the fact that the form of perturbative corrections for the stochastic problem is the same as for the quantum problem, and still the actual calculations are sufficiently simple that one can explore many more orders in perturbation theory than would be possible for a full-fledged field theory, and develop new perturbative methods.
The first method we try is the standard Feynman-diagrammatic expansion. For semiclassical quantum mechanics of a classically chaotic system such calculation was first carried out by Gaspard . The stochastic version described here, implemented by Dettmann , reveals features not so readily apparent in the quantum calculation.
The Feynman diagram method becomes unwieldy at higher orders. The second method, introduced by Vattay , is based on Rugh’s explicit matrix representation of the evolution operator. If one is interested in evaluating numerically many orders of perturbation theory and many eigenvalues, this method is unsurpassed.
The third approach, the smooth conjugations introduced by Mainieri , is perhaps an altogether new idea in field theory. In this approach the neighborhood of each saddlepoint is rectified by an appropriate nonlinear field transformation, with the focus shifted from the dynamics in the original field variables to the properties of the conjugacy transformation. The expressions obtained are equivalent to sums of Feynman diagrams, but are more compact.
## 4 Feynman diagrammatic expansions
We start our computation of the weak noise corrections to the spectrum of $``$ by calculating the trace of the $`n`$-th iterate of the stochastic evolution operator $``$. A convenient choice of noise is Gaussian, $`P(\xi )=e^{\xi ^2/2}/\sqrt{2\pi },`$ with the trace given by an $`n`$-dimensional integral on $`n`$ points along a discrete periodic chain
$`\mathrm{tr}^n`$ $`=`$ $`{\displaystyle 𝑑x_0\mathrm{}𝑑x_{n1}(x_0,x_{n1})\mathrm{}(x_1,x_0)}`$ (9)
$`=`$ $`{\displaystyle [dx]\mathrm{exp}\left\{\frac{1}{2\sigma ^2}\underset{a}{}\left[x_{a+1}f(x_a)\right]^2\right\}}`$
$`x_n=x_0,[dx]={\displaystyle \underset{a=0}{\overset{n1}{}}}{\displaystyle \frac{dx_a}{\sqrt{2\pi \sigma ^2}}}.`$
The choice of Gaussian noise is not essential, as the methods that we develop here apply equally well to other noise distributions, and more generally to the space dependent noise distributions $`P(x,\xi )`$. As the neighborhood of any trajectory is nonlinearly distorted by the flow, the integrated noise is anyway never Gaussian, but colored.
If the classical dynamics is hyperbolic, periodic solutions of given finite period $`n`$ are isolated. Furthermore, if the noise broadening $`\sigma `$ is sufficiently weak they remain distinct, and the dominant contributions come from neighborhoods of periodic points, the tubes sketched in the trace formula (4). In the saddlepoint approximation the trace (9) is given by the sum over neighborhoods of periodic points
$$\mathrm{tr}^n\mathrm{tr}^n|_{\text{sc}}=\underset{x_c\text{Fix}f^n}{}e^{W_c}=\underset{p}{}n_p\underset{r=1}{\overset{\mathrm{}}{}}\delta _{n,n_pr}e^{W_{p^r}}.$$
(10)
As traces are cyclic, $`e^{W_c}`$ is the same for all periodic points in a given cycle, independent of the choice of the starting point $`x_c`$, and the periodic point sum can be rewritten in terms of prime cycles $`p`$ and their repeats. In the deterministic, $`\sigma 0`$ limit this is the discrete time version of the classical trace formula (4). Effects such as noise induced tunnelling are not included in the weak noise approximation.
We now turn to the evaluation of $`W_{p^r}`$, the weight of the $`r`$-th repeat of prime cycle $`p`$. The contribution of the cycle point $`x_a`$ neighborhood is best expressed in an intrinsic coordinate system, by centering the coordinate system on the cycle points,
$$x_ax_a+\varphi _a.$$
(11)
From now on $`x_a`$ will refer to the position of the $`a`$-th periodic point, $`\varphi _a`$ to the deviation of the noisy trajectory from the deterministic one, $`f_a(\varphi _a)`$ to the map (5) centered on the $`a`$-th cycle point, and $`f_a^{(m)}`$ to its $`m`$-th derivative evaluated at the $`a`$-th cycle point:
$$f_a(\varphi _a)=f(x_a+\varphi _a)x_{a+1},f_a^{}=f^{}(x_a),f_a^{\prime \prime }=f^{\prime \prime }(x_a),\mathrm{}.$$
(12)
Rewriting the trace in vector notation, with $`x`$ and $`f(x)`$ $`n`$-dimensional column vectors with components $`x_a`$ and $`f(x_a)`$ respectively, expanding $`f`$ in Taylor series around each of the periodic points in the orbit of $`x_c`$, separating out the quadratic part and integrating we obtain
$`e^{W_c}`$ $`=`$ $`{\displaystyle _c}[d\varphi ]e^{\left(\mathrm{\Delta }^1\varphi V^{}(\varphi )\right)^2/2\sigma ^2}={\displaystyle _c}[d\varphi ]e^{\frac{1}{2\sigma ^2}\varphi ^T\frac{1}{\mathrm{\Delta }^T\mathrm{\Delta }}\varphi +(\mathrm{})}`$ (13)
$`=`$ $`|\text{det}\mathrm{\Delta }|{\displaystyle _c}[d\phi ]e^{{\scriptscriptstyle {\scriptscriptstyle \frac{1}{k}}\mathrm{tr}\left(\mathrm{\Delta }V^{\prime \prime }(\varphi )\right)^k}}e^{\phi ^2/2\sigma ^2}`$
The \[$`n`$$`\times `$$`n`$\] matrix $`\mathrm{\Delta }`$ arises from the quadratic part of the exponent, while all higher powers of $`\varphi _a`$ are collected in $`V(\varphi )`$:
$$\mathrm{\Delta }_{ab}^1\varphi _b=f_a^{^{}}\varphi _a+\varphi _{a+1},V(\varphi )=\underset{a}{}\underset{m=2}{\overset{\mathrm{}}{}}f_a^{(m)}\frac{\varphi _a^{m+1}}{(m+1)!}.$$
(14)
The saddlepoint expansion is most conveniently evaluated in terms of Feynman diagrams, by drawing $`\mathrm{\Delta }`$ as a directed line $`\mathrm{\Delta }_{ab}=\text{}`$, and the derivatives of $`V`$ as the “interaction” vertices <sup>1</sup><sup>1</sup>1PC: fix birdtrack height
$$f_a^{^{\prime \prime }}=\text{},f_a^{^{\prime \prime \prime }}=\text{},\mathrm{}.$$
In the jargon of field theory, $`\mathrm{\Delta }`$ is the “free propagator”. Its determinant
$$\left|\text{det}\mathrm{\Delta }\right|=\frac{1}{|\mathrm{\Lambda }_c1|},\mathrm{\Lambda }_c=\underset{a=0}{\overset{n1}{}}f_a^{^{}}$$
(15)
is the 1-dimensional version of the classical stability weight $`\left|\text{det}(\mathrm{𝟏}𝐉)\right|^1`$ in (4), with $`\mathrm{\Lambda }_c`$ the stability of the $`n`$-cycle going through the periodic point $`x_c`$.
Standard methods now yield the perturbation expansion in terms of the connected “vacuum bubbles”
$`W_c`$ $`=`$ $`\mathrm{ln}|\mathrm{\Lambda }_c1|+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}W_{c,2k}\sigma ^{2k}`$ (16)
$`W_{c,2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}+{\displaystyle \frac{1}{2}}\text{}`$
$`W_{c,4}`$ $`=`$ $`\mathrm{}.`$
In the usual field-theoretic calculations the $`W_{c,0}`$ term corresponds to an overall volume term that cancels out in the expectation values. In contrast, as explained in sect. 2, here the $`e^{W_{c,0}}=|\mathrm{\Lambda }_c1|^1`$ term is the classical volume of cycle $`c`$. Not only does this weight not cancel out in the expectation value formulas, it plays the key role both in classical and semiclassical trace formulas.
In the diagrams sketched above a propagator line connects $`x_a`$ at time $`a`$ with $`x_b`$ at later time $`b`$ by a deterministic trajectory. At time $`b`$ noise induces a kick whose strength depends on the local curvature of the flow. A penalty of a factor $`\sigma `$ is paid, $`m1`$ deterministic trajectories originate in the neighborhood of $`x_b`$ from vertex $`V^{(m)}(x_b)`$, and the process repeats itself, each vertex carrying a penalty of $`\sigma `$, and higher derivatives of the $`f_b`$. Summing over all noise kick sequences encoded by a given diagram and using the periodicity of the trace integral (9) Dettmann obtains expressions such as
$$\frac{r}{2}\frac{\mathrm{\Lambda }_p^{2r}1}{\mathrm{\Lambda }_p^21}\frac{\mathrm{\Lambda }_p^r}{(\mathrm{\Lambda }_p^r1)^3}\underset{ab}{}\left(\frac{f_a^{{}_{}{}^{\prime \prime }2}}{f_a^{}_{}{}^{}2}\frac{f_a^{^{\prime \prime \prime }}}{f_a^{^{}}}\right)\underset{d=b+1}{\overset{a1}{}}f_d^{}_{}{}^{}2.$$
(17)
This particular sum is the Feynman diagram $`\sigma ^2`$ correction to $`r`$-th repeat of prime cycle $`p`$. More algebra leads to similar contributions from the remaining diagrams. But the overall result is surprising; the dependence on the repeat number $`r`$ factorizes, with each diagram yielding the same prefactor depending only on $`\mathrm{\Lambda }_p^r`$. This remarkable fact will be explained in sect. 6. The result of the Feynman-diagrammatic calculations is the stochastic trace formula
$$\mathrm{tr}\frac{z}{1z}|_{\text{sc}}=\underset{p}{}\underset{k=0}{\overset{\mathrm{}}{}}\frac{n_pt_{p,k}}{1t_{p,k}},t_{p,k}=\frac{z^{n_p}}{|\mathrm{\Lambda }_p|\mathrm{\Lambda }_p^k}e^{\frac{\sigma ^2}{2}w_{p,k}^{(2)}+O(\sigma ^4)},$$
(18)
where $`t_{p,k}`$ is the $`k`$-th local eigenvalue evaluated on the $`p`$ cycle. The deterministic, $`\sigma =0`$ part of this formula is the stochastic equivalent of the Gutzwiller semiclassical trace formula . The $`\sigma ^2`$ correction $`w_{p,k}^{(2)}`$ is the stochastic analogue of Gaspard’s $`\mathrm{}`$ correction . At the moment the explicit formula is sufficiently unenlightening that we postpone writing it down to sect. 6.
While the diagrams are standard, the chaotic field theory calculations are considerably more demanding than is usually the case in field theory. Here there is no translational invariance along the chain, so the vertex strength depends on the position, and the free propagator is not diagonalized by a Fourier transform. Furthermore, here one is neither “quantizing” around a trivial vacuum, nor a countable infinity of analytically explicit soliton saddles, but around an infinity of nontrivial unstable hyperbolic saddles.
Two aspects of the above perturbative results are a priori far from obvious: (a) that the structure of the periodic orbit theory should survive introduction of noise, and (b) a more subtle and surprising result, repeats of prime cycles can be re-summed and theory reduced to the dynamical zeta functions and spectral determinants of the same form as for deterministic systems.
Pushing the Feynman-diagrammatic approach to higher orders is laborious, and has not been attempted for this class of problems. As we shall now see, it is not smart to keep pushing it, either, as one can compute many more orders of perturbation theory by means of a matrix representation for $``$.
## 5 Evolution operator in a matrix representation
An expanding map $`f(x)`$ takes an initial smooth distribution $`\varphi (x)`$ defined on a subinterval, stretches it out and overlays it over a larger interval. Repetition of this process smoothes the initial distribution $`\varphi (x)`$, so it is natural to concentrate on smooth distributions $`\varphi _n(x)`$, and represent them by their Taylor series. By expanding both $`\varphi _n(x)`$ and $`\varphi _{n+1}(y)`$ in (6) in Taylor series Rugh derived a matrix representation of the evolution operator
$$𝑑x(y,x)\frac{x^m}{m!}=\underset{m^{}}{}\frac{y^m^{}}{m^{}!}𝐋_{m^{}m},m,m^{}=0,1,2,\mathrm{}$$
which maps the $`x^m`$ component of the density of trajectories $`\varphi _n(x)`$ in (6) to the $`y^m^{}`$ component of the density $`\varphi _{n+1}(y)`$ one time step later. The matrix elements follow by differentiating both sides with $`^m^{}/y^m^{}`$ and evaluating the integral
$$𝐋_{m^{}m}=\frac{^m^{}}{y^m^{}}𝑑x(y,x)\frac{x^m}{m!}|_{y=0}.$$
(19)
In (7) we have written the evolution operator $``$ in terms of the Dirac delta function in order to emphasize that in the weak noise limit the stochastic trajectories are concentrated along the classical trajectory $`y=f(x)`$. Hence it is natural to expand the kernel in a Taylor series in $`\sigma `$
$$(y,x)=\delta (yf(x))+\underset{n=2}{\overset{\mathrm{}}{}}\frac{(\sigma )^n}{n!}\delta ^{(n)}(yf(x))\xi ^nP(\xi )𝑑\xi ,$$
(20)
where $`\delta ^{(n)}(y)=\frac{^n}{y^n}\delta (y).`$ This yields a representation of the evolution operator centered along the classical trajectory, dominated by the deterministic Perron-Frobenius operator $`\delta (yf(x))`$, with corrections given by derivatives of delta functions weighted by moments of the noise distribution $`P_n=P(\xi )\xi ^n𝑑\xi `$. We again center the coordinate system on the cycle points as in (11), and also introduce a notation for the operator (7) centered on the $`x_ax_{a+1}`$ segment of the classical trajectory
$$_a(\varphi _{a+1},\varphi _a)=(x_{a+1}+\varphi _{a+1},x_a+\varphi _a).$$
The weak noise expansion (20) for the $`a`$-th segment operator is given by
$$_a(\varphi _{a+1},\varphi _a)=\delta (\varphi _{a+1}f_a(\varphi _a))+\underset{n=2}{\overset{\mathrm{}}{}}\frac{(\sigma )^n}{n!}P_n\delta ^{(n)}(\varphi _{a+1}f_a(\varphi _a)).$$
(21)
As the evolution operator has a simple $`\delta `$-function form, the local matrix representation of $`_a`$ centered on the $`x_ax_{a+1}`$ segment of the deterministic trajectory can be evaluated recursively in terms of derivatives of the map $`f`$:
$`\left(𝐋_a\right)_{m^{}m}`$ $`=`$ $`{\displaystyle \underset{n}{\overset{\mathrm{}}{}}}P_n{\displaystyle \frac{(\sigma )^n}{n!}}(𝐁_a)_{m^{}+n,m},n=\text{max}(mm^{},0)`$
$`(𝐁_a)_{m^{}m}`$ $`=`$ $`{\displaystyle 𝑑\varphi \delta ^{(m^{})}(\varphi _{a+1}f_a(\varphi ))\frac{\varphi ^m}{m!}}`$
$`=`$ $`{\displaystyle \frac{1}{|f_a^{}|}}\left({\displaystyle \frac{d}{d\varphi }}{\displaystyle \frac{1}{f_a^{}(\varphi )}}\right)^m^{}{\displaystyle \frac{\varphi ^m}{m!}}|_{\varphi =0}.`$
The matrix elements vanish for $`m^{}<m`$, so $`𝐁`$ is a lower triangular matrix. The diagonal and the successive off-diagonal matrix elements are easily evaluated iteratively by computer algebra
$$(𝐁_a)_{mm}=\frac{1}{|f_a^{}|(f_a^{})^m},(𝐁_a)_{m+1,m}=\frac{(m+2)!f_a^{\prime \prime }}{2m!|f_a^{}|(f_a^{})^{m+2}},\mathrm{}.$$
For chaotic systems the map is expanding, $`|f_a^{}|>1`$. Hence the diagonal terms drop off exponentially, as $`1/|f_a^{}|^{m+1}`$, the terms below the diagonal fall off even faster, and truncating $`𝐋_a`$ to a finite matrix introduces only exponentially small errors.
The trace formula (8) takes now a matrix form
$$\mathrm{tr}\frac{z}{1z}|_{\text{sc}}=\underset{p}{}n_p\mathrm{tr}\frac{z^{n_p}𝐋_p}{1z^{n_p}𝐋_p},$$
(23)
where $`𝐋_p=𝐋_{n_p}𝐋_2\mathrm{}𝐋_1`$ is the contribution of the $`p`$ cycle. The subscript sc is a reminder that this is a saddlepoint or semiclassical approximation, valid as an asymptotic series in the limit of weak noise. Vattay interprets the local matrix representation of the evolution operator as follows. The matrix identity log det = tr log together with the trace formula (23) yields
$`\text{det}(1z)|_{\text{sc}}={\displaystyle \underset{p}{}}\text{det}(1z^{n_p}𝐋_p),`$ (24)
so in the saddlepoint approximation the spectrum of the global evolution operator $``$ is pieced together from the local spectra computed cycle-by-cycle on neighborhoods of individual prime cycles with periodic boundary conditions. The meaning of the $`k`$-th term in the trace formula (18) is now clear; it is the $`k`$-th eigenvalue of the local evolution operator restricted to the $`p`$-th cycle neighborhood.
Using this matrix representation Palla and Søndergaard were able to compute corrections to order $`\sigma ^{12}`$, a feat simply impossible along the Feynman-diagrammatic line of attack. In retrospect, the matrix representation method for solving the stochastic evolution is eminently sensible — after all, that is the way one solves a close relative to stochastic PDEs, the Schrödinger equation. What is new is that the problem is being solved locally, periodic orbit by periodic orbit, by translation to coordinates intrinsic to the periodic orbit. It is this natural local basis that makes the matrix representation so simple.
Mainieri takes this observation one step further; as the dynamics is nonlinear, why not search for a nonlinear coordinate transformation that makes the intrinsic coordinates as simple as possible?
## 6 Smooth conjugacies
This step injects into field theory a method standard in the construction of normal forms for bifurcations . The idea is to perform a smooth nonlinear coordinate transformation $`x=h(y)`$, $`f(x)=h(g(h^1(x)))`$ that flattens out the vicinity of a fixed point and makes the map linear in an open neighborhood, $`f(x)g(y)=𝐉y`$.
$`\genfrac{}{}{0pt}{}{\text{ }\text{}\text{ }}{\text{an arbitrary coordinatization}}`$ $``$ $`\genfrac{}{}{0pt}{}{\text{ }\text{}\text{ }}{\text{intrinsic, flat coordinates}}`$
The key idea of flattening the neighborhood of a saddlepoint can be traced back to Poincaré’s celestial mechanics, and is perhaps not something that a field theorist would instinctively hark to as a method of computing perturbative corrections. This local rectification of a map can be implemented only for isolated non-degenerate fixed points (otherwise higher terms are required by the normal form expansion around the point), and only in finite neighborhoods, as the conjugating functions in general have finite radia of convergence.
We proceed in two steps. First, substitution of the weak noise perturbative expansion of the evolution operator (21) into the trace centered on cycle $`c`$ generates products of derivatives of $`\delta `$-functions:
$$\mathrm{tr}^n|_c=\mathrm{}+[d\varphi ]\left\{\mathrm{}\delta ^{(m^{})}(\varphi ^{\prime \prime }f_a(\varphi ^{}))\delta ^{(m)}(\varphi ^{}f_{a1}(\varphi ))\mathrm{}\right\}+\mathrm{}.$$
The integrals are evaluated as in (5), yielding recursive derivative formulas such as
$$𝑑x\delta ^{(m)}(y)=\frac{1}{|y^{}(x)|}\left(\frac{d}{dx}\frac{1}{y^{}(x)}\right)^m|_{y=0},y=f(x)x.$$
(25)
or $`n`$-point integrals, with derivatives distributed over $`n`$ different $`\delta `$-functions.
Next we linearize the neighborhood of the $`a`$-th cycle point. For a 1-dimensional map $`f(x)`$ with a fixed point $`f(0)=0`$ of stability $`\mathrm{\Lambda }=f^{}(0)`$, $`|\mathrm{\Lambda }|1`$ we search for a smooth conjugation $`h(x)`$ such that:
$$f(x)=h(\mathrm{\Lambda }h^1(x)),h(0)=0,h^{}(0)=1.$$
(26)
In higher dimensions $`\mathrm{\Lambda }`$ is replaced by the Jacobian matrix $`𝐉`$. For a periodic orbit each point around the cycle has a differently distorted neighborhood, with differing second and higher derivatives, so the conjugation function $`h_a`$ has to be computed point by point,
$$f_a(\varphi )=h_{a+1}(f_a^{}h_a^1(\varphi )).$$
An explicit expression for $`h_a`$ in terms of $`f`$ is obtained by iterating around the whole cycle, and using the chain rule (15) for the cycle stability $`\mathrm{\Lambda }_p`$
$$f_a^{n_p}(\varphi )=h_a(\mathrm{\Lambda }_ph_a^1(\varphi )),$$
(27)
so each $`h_a`$ is given by some combination of $`f_a`$ derivatives along the cycle. Expand $`f(x)`$ and $`h(x)`$
$$f(x)=\mathrm{\Lambda }x+x^2f_2+x^3f_3+\mathrm{},h(y)=y+y^2h_2+y^3h_3+\mathrm{},$$
and equate recursively coefficients in the functional equation $`h(\mathrm{\Lambda }y)=f(h(y))`$ expansion
$$h(\mathrm{\Lambda }u)\mathrm{\Lambda }h(u)=\underset{n=2}{\overset{\mathrm{}}{}}f_m\left(h(u)\right)^m.$$
(28)
This yields the expansion for the conjugation function $`h`$ in terms of the mapping $`f`$
$$h_2=\frac{f_2}{\mathrm{\Lambda }(\mathrm{\Lambda }1)},h_3=\frac{2f_2^2+\mathrm{\Lambda }(\mathrm{\Lambda }1)f_3}{\mathrm{\Lambda }^2(\mathrm{\Lambda }1)(\mathrm{\Lambda }^21)},\mathrm{}.$$
(29)
The periodic orbit conjugating functions $`h_a`$ are obtained in the same way from (27), with proviso that the cycle stability is not marginal, $`|\mathrm{\Lambda }_p|1`$.
What is gained by replacing the perturbation expansion in terms of $`f^{(m)}`$ by still messier perturbation expansion for the conjugacy function $`h`$? Once the neighborhood of a fixed point is linearized, the conjugation formula for the repeats of the map
$$f^r(x)=h(\mathrm{\Lambda }^rh^1(x))$$
can be used to compute derivatives of a function composed with itself $`r`$ times. The expansion for arbitrary number of repeats depends on the conjugacy function $`h(x)`$ computed for a single repeat, and all the dependence on the repeat number is carried by polynomials in $`\mathrm{\Lambda }^r`$, a result that emerged as a surprise in the Feynman diagrammatic approach of sect. 4. The integrals such as (25) evaluated on the $`r`$-th repeat of prime cycle $`p`$
$$y(x)=f^{n_pr}(x)x$$
(30)
have a simple dependence on the conjugating function $`h`$
$`{\displaystyle \frac{1}{3!}}{\displaystyle \frac{^2}{y^2}}{\displaystyle \frac{1}{y^{}\left(0\right)}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }^r\left(1+\mathrm{\Lambda }^r\right)}{\left(\mathrm{\Lambda }^r1\right)^3}}\left(2h_2^2h_3\right)`$ (31)
$`{\displaystyle \frac{1}{4!}}{\displaystyle \frac{^3}{y^3}}{\displaystyle \frac{1}{y^{}\left(0\right)}}`$ $`=`$ $`5\mathrm{\Lambda }^r{\displaystyle \frac{\left(\mathrm{\Lambda }^r+1\right)^2}{\left(\mathrm{\Lambda }^r1\right)^4}}h_2^3+\mathrm{\Lambda }^r{\displaystyle \frac{5\mathrm{\Lambda }^{2r}+8\mathrm{\Lambda }^r+5}{\left(\mathrm{\Lambda }^r1\right)^4}}h_2h_3\mathrm{\Lambda }^r{\displaystyle \frac{\mathrm{\Lambda }^{2r}+\mathrm{\Lambda }^r+1}{\left(\mathrm{\Lambda }^r1\right)^4}}h_4`$
$`\mathrm{}`$ $`=`$ $`\mathrm{}`$
The evaluation of $`n`$-point integrals is more subtle . The final result of all these calculations is that expressions of form (31) depend on the conjugation function determined from the iterated map, with the saddlepoint approximation to the spectral determinant given by
$$\text{det}(1z_\sigma )|_{\text{sc}}=\underset{p}{}\underset{k=0}{\overset{\mathrm{}}{}}(1t_{p,k})$$
in terms of local $`p`$-cycle eigenvalues
$`t_{p,k}`$ $`=`$ $`{\displaystyle \frac{z^{n_p}}{|\mathrm{\Lambda }_p|\mathrm{\Lambda }_p^k}}\mathrm{}^{\frac{\sigma ^2}{2}P_2w_{p,k}^{(2)}+\frac{\sigma ^3}{3!}P_3w_{p,k}^{(3)}+\frac{\sigma ^4}{4!}P_4w_{p,k}^{(4)}+O(\sigma ^6)}`$
$`w_{p,k}^{(2)}`$ $`=`$ $`(k+1)^2{\displaystyle \underset{a}{}}(2h_{a,2}^2h_{a,3}),w_{p,k}^{(3)}=\mathrm{},\mathrm{}.`$
accurate up to order $`\sigma ^4`$. $`w^{(3)}`$, $`w^{(4)}`$ are also computed by Dettmann, but we desist from citing them here; the reader is referred to ref. . What is remarkable about these results is their simplicity when expressed in terms of the conjugation function $`h`$, as opposed to the Feynman diagram sums, in which each diagram contributes a sum like the one in (17), or worse. Furthermore, both the conjugation and the matrix approaches are easily automatized, as they require only recursive evaluation of derivatives, as opposed to the handcrafted Feynman diagrammar.
Simple minded as they might seem, discrete stochastic processes are a great laboratory for testing ideas that would otherwise be hard to test. Dettmann, Palla and Søndergaard have used a 1-dimensional repeller of bounded nonlinearity and complete binary symbolic dynamics to check numerically the above results, and computed the leading eigenvalue of $``$ by no less than five different methods. As anticipated by Rugh , the evolution operator eigenvalues converge super-exponentially with the cycle length; addition of cycles of period $`(n`$+1$`)`$ to the set of all cycles up to length $`n`$ doubles the number of significant digits in the perturbative prediction. However, as the series is asymptotic, for realistic values of the noise strength summations beyond all orders are needed .
## 7 Summary
The periodic orbit theory approach to turbulence is to visualize turbulence as a sequence of near recurrences in a repertoire of unstable spatio-temporal patterns. The investigations of the Kuramoto-Sivashinsky system discussed above are first steps in the direction of implementing this program. So far, existence of a hierarchy of spatio-temporally periodic solutions of spatially extended nonlinear system has been demonstrated, and the periodic orbit theory has been tested in evaluation of global averages for such system. The parameter ranges tested so far probe the weakest nontrivial “turbulence”, and it is an open question to what extent the approach remains implementable as the system goes more turbulent.
The most important lesson of this investigation is that the unstable spatio-temporally periodic solutions do explore systematically the repertoire of admissible spatio-temporal patterns, with the trace and spectral determinants formulas and their cycle expansions being the proper tools for extraction of quantitative predictions from the periodic orbits data.
We formulat next a semiclassical perturbation theory for stochastic trace formulas with support on infinitely many chaotic saddles. The central object of the periodic orbit theory, the trace of the evolution operator, is a discrete path integral, similar to those found in field theory and statistical mechanics. The weak noise perturbation theory, likewise, resembles perturbative field theory, and can be cast into the standard field-theoretic language of Feynman diagrams. However, we found out that both the matrix and the nonlinear conjugacy perturbative methods are superior to the standard approach. In contrast to previous perturbative expansions around vacua and instanton solutions, the location and local properties of each saddlepoint must be found numerically.
The key idea in the new formulation of perturbation theory is this: Instead of separating the action into quadratic and “interaction” parts, one first performs a nonlinear field transformation which turns the saddle point into an exact quadratic form. The price one pays for this is the Jacobian of the nonlinear field transformation — but it turns out that the perturbation expansion of this Jacobian in terms of the conjugating function is order-by-order more compact than the Feynman-diagrammatic expansion.
## Acknowledgements
I am indebted to my collaborators C.P. Dettmann, G. Vattay, F. Christiansen, V. Putkaradze, G. Palla, N. Søndergaard, R. Mainieri and H.H. Rugh for co-suffering through all the details omitted in this overview. I am grateful to E.A. Spiegel, L. Tuckerman and M.J. Feigenbaum for patient instruction. I am not grateful to those directors and gentlemen of committees who do not find theoretical physics a vibrant subject. |
warning/0001/physics0001040.html | ar5iv | text | # Black-Scholes option pricing within Itô and Stratonovich conventions
## 1 Introduction
An European option is a financial instrument giving to its owner the right but not the obligation to buy (European call) or to sell (European put) a share at a fixed future date, the maturing time $`T`$, and at a certain price called exercise or striking price $`x_C`$. In fact, this is the most simple of a large variety of contracts that can be more sophisticated. One of those possible extensions is the American option which gives the right to exercise the option at any time until the maturing time. In a certain sense, options are a security for the investor thus avoiding the unpredictable consequences of operating with risky speculative stocks.
The trading of options and their theoretical study have been known for long, although they were relative obscure and unimportant financial instruments until the early seventies. It was then when options experimented an spectacular development. The Chicago Board Options Exchange, created in 1973, is the first attempt to unify options in one market and trade them on only a few stock shares. The market rapidly became a tremendous success and led to a series of innovations in option trading .
The main purpose in option studies is to find a fair and presumably riskless price for these instruments. The first solution to the problem was given by Bachelier in 1900 , and several option prices were proposed without being completely satisfactory . However, in the early seventies it was finally developed a complete option valuation based on equilibrium theoretical hypothesis for speculative prices. The works of Fisher Black, Myron Scholes and Robert Merton were the culmination of this great effort, and left the doors open for extending the option pricing theory in many ways. In addition, the method has been proved to be very useful for investors and has helped to option markets to have the importance that they have nowadays in finance .
The option pricing method obtains the so-called Black-Scholes equation which is a partial differential equation of the same kind as the diffusion equation. In fact, it was this similarity that led Black and Scholes to obtain their option price formula as the solution of the diffusion equation with the initial and boundary conditions given by the option contract terms. Incidentally, these physics studies applied to economy have never been disrupted and there still is a growing effort of the physics community to understand the dynamics of finance from approaches similar to those that tackle complex systems in physics .
The economic ideas behind the Black-Scholes option pricing theory translated to the stochastic methods concepts are as follows. First, the option price depends on the stock price and this is a random variable evolving with time. Second, the efficient market hypothesis , i.e., the market incorporates instantaneously any information concerning future market evolution, implies that the random term in the stochastic equation must be delta-correlated. That is: speculative prices are driven by white noise . It is known that any white noise can be written as a combination of the derivative of the Wiener process and white shot noise . In this framework, the Black-Scholes option pricing method was first based on the geometric Brownian motion , and it was lately extended to include white shot noise .
As is well known, any stochastic differential equation (SDE) driven by a state dependent white noise, such as the geometric Brownian motion, is meaningless unless an interpretation of the multiplicative noise term is given. Two interpretations have been presented: Itô and Stratonovich . To our knowledge, all derivations of the Black-Scholes equation starting from a SDE are based on the Itô interpretation. A possible reason is that mathematicians prefer this interpretation over the Stratonovich’s one, being the latter mostly preferred among physicists. Nonetheless, as we try to point out here, Itô framework is perhaps more convenient for finance being this basically due to the peculiarities of trading (see Sect. 4). In any case, as Van Kampen showed some time ago no physical reason can be attached to the interpretation of the SDE modelling price dynamics. However, the same physical process results in two different SDEs depending on the interpretation chosen. In spite of having different differential equations as starting point, we will show that the resulting Black-Scholes equation is the same regardless the interpretation of the multiplicative noise term, and this constitutes the main result of the paper. In addition, the mathematical exercise that represents this translation into the Stratonovich convention provides a useful review, specially to physicists, of the option pricing theory and the “path-breaking” Black-Scholes method.
The paper is divided in 5 sections. After the Introduction, a summary of the differences between Itô and Stratonovich calculus is developed in Section 2. The following section is devoted to explain the market model assumed in Black-Scholes option pricing method. Section 4 concentrates in the derivation of the Black-Scholes equation using both Itô and Stratonovich calculus. Conclusions are drawn in Section 5, and some technical details are left to appendices.
## 2 Itô vs. Stratonovich
It is not our intention to write a formal discussion on the differences between Itô and Stratonovich interpretations of stochastic differential equations since there are many excellent books and reviews on the subject . However, we will summarize those elements in these interpretations that change the treatment of the Black-Scholes option pricing method. In all our discussion, we use a notation that it is widely used among physicists.
The interpretation question arises when dealing with a multiplicative stochastic differential equation, also called multiplicative Langevin equation,
$$\dot{X}=f(X)+g(X)\xi (t),$$
(1)
where $`f`$ and $`g`$ are given functions, and $`\xi (t)`$ is Gaussian white noise, that is, a Gaussian and stationary random process with zero mean and delta correlated. Alternatively, Eq. (1) can be written in terms of the Wiener process $`W(t)`$ as
$$dX=f(X)dt+g(X)dW(t),$$
(2)
where $`dW(t)=\xi (t)dt`$. When $`g`$ depends on $`X`$, Eqs. (1) and (2) have no meaning, unless an interpretation of the multiplicative term $`g(X)\xi (t)`$ is provided. These different interpretations of the multiplicative term must be given because, due to the extreme randomness of white noise, it is not clear what value of $`X`$ should be used even during an infinitesimal timestep $`dt`$. According to Itô, that value of $`X`$ is the one before the beginning of the timestep, i.e., $`X=X(t)`$, whereas Stratonovich uses the value of $`X`$ at the middle of the timestep: $`X=X(t+dt/2)=X(t)+dX(t)/2`$.
Before proceeding further with the consequences of the above discussion, we will first give a precise meaning of the differential of random processes driven by Gaussian white noise and its implications. Obviously, the differential of any random process $`X(t)`$ is defined by
$$dX(t)X(t+dt)X(t).$$
(3)
On the other hand, the differential $`dX(t)`$ of any random process is equal (in the mean square sense) to its mean value if its variance is, at least, of order $`dt^2`$ : $`[dX(t)dX(t)]^2=O(dt^2)`$. We observe that from now on all the results of this paper must be interpreted in the mean square sense. The mean square limit relation can be used to show that $`|dW(t)|^2=dt`$ . We thus have from Eq. (2) that
$$|dX|^2=|g(X)|^2dt+O(dt^2),$$
(4)
and we symbolically write
$$dX(t)=O\left(dt^{1/2}\right).$$
(5)
Let us now turn our attention to the differential of the product of two random processes since this differential adopts a different expression depending on the interpretation (Itô or Stratonovich) chosen. In accordance to Eq. (3), we define
$$d(XY)[(X+dX)(Y+dY)]XY.$$
(6)
This expression can be rewritten in many different ways. One possibility is
$$d(XY)=\left(X+\frac{dX}{2}\right)dY+\left(Y+\frac{dY}{2}\right)dX,$$
(7)
but it is also allowed to write the product as
$$d(XY)=XdY+YdX+dXdY.$$
(8)
Therefore, we say that the differential of a product reads in the Stratonovich interpretation when
$$d(XY)X_SdY+Y_SdX,$$
(9)
where
$$X_S(t)X(t+dt/2)=X(t)+dX(t)/2,$$
(10)
and similarly for $`Y_S(t)`$. Whereas we say that the differential of a product follows the Itô interpretation when
$$d(XY)X_IdY+Y_IdX+dXdY,$$
(11)
where
$$X_I(t)X(t),$$
(12)
and $`Y_I(t)Y(t)`$. Note that Eq. (9) formally agrees with the rules of calculus while Eq. (11) does not. Note also that Eqs. (9) and (11) can easily be generalized to the product of two functions, $`U(X)`$ and $`V(X)`$, of the random process $`X=X(t)`$. Thus
$$d(UV)=U(X_S)dV(X)+V(X_S)dU(X),$$
(13)
where $`X_S`$ is given by Eq. (10), and $`dV(X)=V(X+dX)V(X)`$ with an analogous expression for $`dU(X)`$. Within Itô convention we have
$$d(UV)=U(X)dV(X)+V(X)dU(X)+dU(X)dV(X).$$
(14)
Let us now go back to Eq. (1) and see that one important consequence of the above discussion is that the expected value of the multiplicative term, $`g(X)\xi (t)`$, depends on the interpretation given. In the Itô interpretation, it is clear that $`g(X)\xi (t)=0`$ because the value of $`X`$ (and, hence the value of $`g(X)`$) anticipates the jump in the noise. In other words, $`g(X)`$ is independent of $`\xi (t)`$. On the other hand, it can be proved that within the Stratonovich framework the average of the multiplicative term reads $`g(X)g^{}(X)/2`$ where the prime denotes the derivative . The zero value of the average $`g(X)\xi (t)`$ makes Itô convention very appealing because then the deterministic equation for the mean value of $`X`$ only depends on the drift term $`f(X)`$. In this sense, note that any multiplicative stochastic differential equation has different expressions for the functions $`f(X)`$ and $`g(X)`$ depending on the interpretation chosen. In the Stratonovich framework, a SDE of type Eq. (2) can be written as
$$dX=f^{(S)}(X_S)dt+g^{(S)}(X_S)dW(t),$$
(15)
where $`X_S=X+dX/2`$. In the Itô sense we have
$$dX=f^{(I)}(X_I)dt+g^{(I)}(X_I)dW(t),$$
(16)
where $`X_I=X`$. Note that $`f^{(S)}`$ and $`f^{(I)}`$ are not only evaluated at different values of $`X`$ but are also different functions depending on the interpretation given, and the same applies to $`g^{(S)}`$ and $`g^{(I)}`$. One can easily show from Eq. (10) and Eqs. (15)-(16) that, after keeping terms up to order $`dt`$, the relation between $`f_S`$ and $`f_I`$ is
$$f^{(I)}(X)=f^{(S)}(X)\frac{1}{2}g^{(S)}(X)\frac{g^{(S)}(X)}{X},$$
(17)
while the multiplicative functions $`g^{(S)}`$ and $`g^{(I)}`$ are equal
$$g^{(I)}(X)=g^{(S)}(X).$$
(18)
Conversely, it is possible to pass from a Stratonovich SDE to an equivalent Itô SDE . Note that the difference between both interpretation only affects the drift term given by the function $`f`$ while the function $`g`$ remains unaffected. In addition, we see that for an additive SDE, i.e., when $`g`$ is independent of $`X`$, the interpretation question is irrelevant.
Finally, a crucial difference between Itô and Stratonovich interpretations appears when a change of variables is performed on the original equation. Then it can be proved that, using Stratonovich convention, the standard rules of calculus hold, but new rules appear when the equation is understood in the Itô sense. From the point of view of this property, the Stratonovich criterion seems to be more convenient. For the sake of completeness, we remind here what are the rules of change of variables in each interpretation. Let $`h(X,t)`$ be an arbitrary function of $`X`$ and $`t`$. In the Itô sense, the differential of $`h(X,t)`$ reads
$$dh=\frac{h(X,t)}{X}dX+\left[\frac{h(X,t)}{t}+\frac{1}{2}g^2(X,t)\frac{^2h(X,t)}{X^2}\right]dt,$$
(19)
whereas in the Stratonovich sense, we have the usual expression
$$dh=\frac{h(X_S,t)}{X_S}dX+\frac{h(X_S,t)}{t}dt,$$
(20)
where
$$\frac{h(X_S,t)}{X_S}=\frac{h(X,t)}{X}|_{X=X_S},$$
and $`X_S`$ is given by Eq. (10).
Equation (19) is known as the Itô’s lema and it is extensively used in mathematical finance books .
The information on the properties of the Itô and Stratonovich interpretation of SDE contained in this brief summary is sufficient to follow the derivations of the next sections.
## 3 Market model
Option pricing becomes a problem because market prices or indexes change randomly. Therefore, any possible calculation of an option price is based on a model for the stochastic evolution of the market prices. The first analysis of price changes was given one hundred years ago by Bachelier who, studying the option pricing problem, proposed a model assuming that price changes behave as an ordinary random walk . Thus, in the continuum limit (continous time finance ) speculative prices obey a Langevin equation. In order to include the limited liability of the stock prices, i.e., prices cannot be negative, Osborne proposed the geometric or log-Brownian motion for describing the price changes . Mathematically, the market model assumed by Osborne can be written as a stochastic equation of type Eq. (2):
$$dR(t)=\mu dt+\sigma dW(t),$$
(21)
where $`R(t)`$ is the so-called return rate after a period $`t`$. Therefore, $`dR(t)`$ is the infinitessimal relative change in the stock share price $`X(t)`$ (see Eq. (22) below), $`\mu `$ is the average rate per unit time, and $`\sigma ^2`$ is the volatility per unit time of the rate after a period $`t`$, i.e., $`dR=\mu dt`$ and $`(dRdR)^2=\sigma ^2dt`$. There is no need to specify an interpretation (Itô’s or Stratonovich’s) for Eq. (21) because $`\sigma `$ is constant and we are thus dealing with an additive equation. The rate is compounded continuously and, therefore, an initial price $`X_0`$ becomes after a period $`t`$:
$$X(t)=X_0\mathrm{exp}[R(t)].$$
(22)
This equation can be used as a change of variables to derive the SDE for $`X(t)`$ given that $`R(t)`$ evolves according to Eq. (21). However, as it becomes multiplicative, we have to attach the equation to an interpretation. Indeed, using Stratonovich calculus (see Eq. (20)), it follows that $`X(t)`$ evolves according to the equation
$$dX=\mu X_Sdt+\sigma X_SdW(t),$$
(23)
where $`X_S=X+dX/2`$. In the Itô sense (see Eq. (19)), the equation for $`X(t)`$ becomes
$$dX=\left(\mu +\sigma ^2/2\right)Xdt+\sigma XdW(t).$$
(24)
Therefore, the Langevin equation for $`X(t)`$ is different depending on the sense it is interpreted. The main objective of this paper is to show that no matter which equation is used to derive the Black-Scholes equation the final result turns out to be the same.
Before proceeding further, we point out that the average index price after a time $`t`$ is $`X(t)=X_0\mathrm{exp}(\mu +\sigma ^2/2)t`$, regardless the convention being used. In fact, the independence of the averages on the interpretation used holds for moments of any order .
## 4 The Black-Scholes equation
There are several different approaches for deriving the Black-Scholes equation starting from the stochastic differential equation point of view. These different derivations only differ in the way the portfolio (i.e., a collection of different assets for diversifying away financial risk) is defined . In order to get the most general description of the concepts underlying in the Black-Scholes theory, our portfolio is similar to the one proposed by Merton , and it is based on one type of share whose price is the random process $`X(t)`$. The portfolio is compounded by a certain amount of shares, $`\mathrm{\Delta }`$, a number of calls, $`\mathrm{\Psi }`$, and, finally, a quantity of riskless securities (or bonds) $`\mathrm{\Phi }`$. We also assume that short-selling, or borrowing, is allowed. Specifically, we own a certain number of calls worth $`\mathrm{\Psi }C`$ dollars and we owe $`\mathrm{\Delta }X+\mathrm{\Phi }B`$ dollars. In this case, the value $`P`$ of the porfolio reads
$$P=\mathrm{\Psi }C\mathrm{\Delta }X\mathrm{\Phi }B,$$
(25)
where $`X`$ is the share stock price, $`C`$ is the call price to be determined, and $`B`$ is the bond price whose evolution is not random and is described according to the value of $`r`$, the risk-free interest rate ratio. That is
$$dB=rBdt.$$
(26)
The so-called “portfolio investor’s strategy” decides the quantity to be invested in every asset according to its stock price at time $`t`$. This is the reason why the asset amounts $`\mathrm{\Delta },\mathrm{\Psi },`$ and $`\mathrm{\Phi }`$ are functions of stock price and time, although they are “nonanticipating” functions of the stock price. This somewhat obscure concept is explained in the Appendix A. All derivations of Black-Scholes equation assume a “frictionless market”, that is, there are no transaction costs for each operation of buying and selling .
According to Merton we assume that, by short-sales, or borrowing, the portfolio (25) is constrained to require net zero investment, that is, $`P=0`$ for any time $`t`$ . Then, from Eq. (25) we have
$$C=\delta _nX+\varphi _nB,$$
(27)
where, $`\delta _n\mathrm{\Delta }/\mathrm{\Psi }`$ and $`\varphi _n\mathrm{\Phi }/\mathrm{\Psi }`$ are respectively the number of shares per call and the number of bonds per call. As we have mentioned above, $`\delta _n`$ and $`\varphi _n`$ are nonanticipating functions of the stock price (see Appendix A). Note that Eq. (27) has an interesting economic meaning, since tells us that having a call option is equivalent to possess a certain number, $`\delta _n`$ and $`\varphi _n`$, of shares and bonds thus avoiding any arbitrage opportunity . Equation (27), which is called “the replicating portfolio” , is the starting point of our derivation that we separate into two subsections according to Itô or Stratonovich interpretations.
### 4.1 The Black-Scholes equation derivation (Itô)
We need first to obtain, within the Itô interpretation, the differential of the call price $`C`$. This is done in the Appendix B and we show there that
$$dC=\delta dX+\varphi dB+Xd\delta _n+Bd\varphi _n+O(dt^{3/2}),$$
(28)
where the relationship between $`\delta `$, $`\varphi `$ and $`\delta _n`$, $`\varphi _n`$ is given in Appendix A (cf. Eq. (46)). We assume we follow a “self-financing strategy” , that is, variations of wealth are only due to capital gains and not to the withdrawal or infusion of new funds. In other words, we increase \[decrease\] the number of shares by selling \[buying\] bonds in the same proportion. We then have (see Appendix A for more details)
$$Xd\delta _n=Bd\varphi _n,$$
(29)
and Eq. (28) reads
$$dC=\delta dX+\varphi dB.$$
(30)
Moreover, from Eqs. (26)-(27) one can easily show that
$$\varphi dB=r(C\delta X)dt+O(dt^{3/2}),$$
(cf. Eq. (5) and Eq. (46) of Appendix A). Therefore,
$$dC=\delta dX+r(C\delta X)dt+O(dt^{3/2}).$$
(31)
On the other hand, since the call price $`C`$ is a function of share price $`X`$ and time, $`C=C(X,t)`$, and $`X`$ obeys the (Itô) SDE (24), then $`dC`$ can be evaluated from the Itô lemma (19) with the result
$$dC=\left(\frac{C}{t}+\frac{1}{2}\sigma ^2X^2\frac{^2C}{X^2}\right)dt+\frac{C}{X}dX.$$
(32)
Substituting Eq. (31) into Eq. (32) yields
$$\left(\delta \frac{C}{X}\right)dX=\left[\frac{C}{t}r(C\delta X)+\frac{1}{2}\sigma ^2X^2\frac{^2C}{X^2}\right]dt.$$
(33)
Note that this is an stochastic equation because of its dependence on the Wiener process enclosed in $`dX`$. We can thus turn Eq. (33) into a deterministic equation that will give the call price functional dependence on share price and time by equating to zero the term multiplying $`dX`$. This, in turn, will determine the “investor strategy”, that is the number of shares per call, the so called “delta hedging”:
$$\delta =\frac{C(x,t)}{x}.$$
(34)
The substitution of Eq. (34) into Eq. (33) results in the Black-Scholes equation:
$$\frac{C}{t}=rCrx\frac{C}{x}\frac{1}{2}(\sigma x)^2\frac{^2C}{x^2}.$$
(35)
A final observation, in Eqs. (34)-(35) we have set $`X=x`$, since, as explained above, Eq. (35) gives the functional dependence of the call price $`C`$ on $`X`$ and $`t`$ regardless whether the share price $`X`$ is random or not.
### 4.2 The Black-Scholes equation derivation (Stratonovich)
Let us now derive the Black-Scholes equation, assuming that the underlying asset obeys the Stratonovich SDE (23). In the Appendix B we present part of this derivation using the concept of nonanticipating function within the Stratonovich interpretation. Nevertheless, here we perform an alternative derivation that uses the Itô interpretation as starting point. We thus begin with Eq. (31) that we write in the form
$$dC=\delta (X,t)dX(t)+r\left[C(X,t)\delta (X,t)X\right]dt+O(dt^{3/2}).$$
(36)
Now, we have to express the function $`\delta `$ within Stratonovich interpretation. Note that $`X=X_SdX/2`$. Hence $`\delta (X,t)=\delta (X_SdX/2,t)`$, whence
$$\delta (X,t)=\delta (X_S,t)\frac{1}{2}\frac{\delta (X_S,t)}{X_S}dX+O(dX^2).$$
(37)
Analogously $`C(X,t)=C(X_S,t)+O(dX)`$. Therefore, from Eqs. (36)-(37) and taking into account Eq. (4) we have
$`dC=\delta (X_S,t)dX+[rC(X_S,t)`$ $``$ $`rX_S\delta (X_S,t)`$ (38)
$``$ $`{\displaystyle \frac{1}{2}}\sigma ^2X_S^2{\displaystyle \frac{\delta (X_S,t)}{X_S}}]dt+O(dt^{3/2}).`$
On the other hand, $`dC`$ will also be given by Eq. (20)
$$dC=\frac{C(X_S,t)}{t}dt+\frac{C(X_S,t)}{X_S}dX,$$
(39)
From these two equations we get
$`\left[\delta (X_S,t){\displaystyle \frac{C(X_S,t)}{X_S}}\right]dX`$ $`=`$ $`[{\displaystyle \frac{C(X_S,t)}{t}}rC(X_S,t)`$ (40)
$`+`$ $`rX_S\delta (X_S,t)+{\displaystyle \frac{1}{2}}\sigma ^2X_S^2{\displaystyle \frac{\delta (X_S,t)}{X_S}}]dt.`$
Again, this equation becomes non-stochastic if we set
$$\delta (X_S,t)=\frac{C(X_S,t)}{X_S}.$$
(41)
In this case, the combination of Eqs. (40)-(41) agrees with Eq. (35). Although the call price is evaluated at a different value of the share price, this is irrelevant for the reason explained right after Eq. (35). Therefore, the Stratonovich calculus results in the same call price formula and equation than the Itô calculus.
We have used the stochastic differential equation technique in order to derive the option price equation. However, this is only one of the possible routes. Another way, which was also proposed in the original paper of Black and Scholes , uses the Capital Asset Pricing Model (CAPM) where, adducing equilibrium reasons in the asset prices, it is assumed the equality of the so-called “Sharpe ratio” of the stock and the option respectively. The Sharpe ratio of an asset can be defined as its normalized excess of return, therefore CAPM assumption applied to option pricing reads
$$\frac{\alpha r}{\sigma }=\frac{\alpha _Cr}{\sigma _C},$$
where $`\alpha =dX/X`$, $`\sigma ^2=\text{Var}(dX/X)`$, $`\alpha _C=dC/C`$, and $`\sigma _C^2=\text{Var}(dC/C)`$. From this equality it is quite straightforward to derive the Black-Scholes equation . As remarked at the end of Sect. 3, moments are independent of the interpretation chosen, we thus clearly see the equivalence between Itô and Stratonovich calculus for the Black-Scholes equation derivation.
### 4.3 The Black-Scholes formula for the European call
For the sake of completeness, let us now finish the paper by shortly deriving from Eq. (35) the well-known Black-Scholes formula. Note that the Black-Scholes equation is a backward parabolic differential equation, we therefore need one “final” condition and, in principle, two boundary conditions in order to solve it . In fact, Black-Scholes equation is defined on the semi-infinite interval $`0x<\mathrm{}`$. In this case, since $`C(x,t)`$ is assumed to be sufficiently well behaved for all $`x`$, we only need to specify one boundary condition at $`x=0`$ (see and ), although we specify below the boundary condition at $`x=\mathrm{}`$ as well.
We also note that all financial derivatives (options of any kind, forwards, futures, swaps, etc…) have the same boundary conditions but different initial or final condition . Let us first specify the boundary conditions. We see from the multiplicative character of Eq. (2) that if at some time the price $`X(t)`$ drops to zero then it stays there forever. In such a case, it is quite obvious that the call option is worthless:
$$C(0,t)=0.$$
(42)
On the other hand, as the share price increases without bound, $`X\mathrm{}`$, the difference between share price and option price vanishes, since option is more and more likely to be exercised and the value of the option will agree with the share price, that is,
$$\underset{x\mathrm{}}{lim}\frac{C(x,t)}{x}=1.$$
(43)
In order to obtain the “final” condition for Eq. (35), we need to specify the following two parameters: the expiration or maturing time $`T`$, and the striking or exercise price $`x_C`$ that fixes the price at which the call owner has the right to buy the share at time $`T`$. If we want to avoid arbitrage opportunities, it is clear that the value of the option $`C`$ of a share that at time $`T`$ is worth $`x`$ dollars must be equal to the payoff for having the option . This payoff is either $`0`$ or the difference between share price at time $`T`$ and option striking price, that is, $`\mathrm{max}(xx_C,0)`$. Hence, the “final” condition for the European call is
$$C(x,t=T)=\mathrm{max}(xx_C,0).$$
(44)
In the Appendix C we show that the solution to the problem given by Eq. (35) and Eqs. (42)-(44) is
$$C(x,t)=xN(d_1)x_Ce^{r(Tt)}N(d_2),$$
(45)
$`(0tT)`$, where
$$N(z)=\frac{1}{\sqrt{2\pi }}_{\mathrm{}}^ze^{u^2/2}𝑑u,$$
is the probability integral,
$$d_1=\frac{\mathrm{ln}(x/x_c)+(r+\sigma ^2/2)(Tt)}{\sigma \sqrt{Tt}},$$
and
$$d_2=d_1\sigma \sqrt{Tt}.$$
## 5 Conclusions
We have updated the option pricing theory from the point of view of a physicist. We have centered our analysis of option pricing to the Black-Scholes equation and formula for the European call, extensions to other kind of options can be straightforward in many cases and are found in several good finance books . We have reviewed Black-Scholes theory using Itô calculus, which is standard to mathematical finance, with a special emphasis in explaining and clarifying the many subtleties of the calculation. Nevertheless, we have not limit ourselves only to review option pricing, but to derive, for the first time to our knowledge, the Black-Scholes equation using the Stratonovich calculus which is standard to physics, thus bridging the gap between mathematical finance and physics.
As we have proved, the Black-Scholes equation obtained using Stratonovich calculus is the same as the one obtained by means of the Itô calculus. In fact, this is the result we expected in advance because Itô and Stratonovich conventions are just different rules of calculus. Moreover, from a practical point of view, both interpretations differ only in the drift term of the Langevin equation and the drift term does not appear in the Black-Scholes equation and formula. But, again, we think that this derivation is still interesting and useful for all the reasons explained above.
This work has been supported in part by Dirección General de Investigación Científica y Técnica under contract No. PB96-0188 and Project No. HB119-0104, and by Generalitat de Catalunya under contract No. 1998 SGR-00015.
## Appendix A Nonanticipating functions and self-financing strategy
The functionals $`\varphi _n`$ and $`\delta _n`$ representing normalized asset quantities are nonaticipating functions with respect to the stock price $`X`$. This means that these functionals are in some way independent of $`X(t)`$ implying a sort of causality in the sense that unknown future stock price cannot affect the present portfolio strategy. The physical meaning of this translated to financial markets is: first buy or sell according to the present stock price $`X(t)`$ and right after the portfolio worth changes with variation of the prices $`dX`$, $`dB`$, and $`dC`$. In other words, the investor strategy does not anticipate the stock price change . Therefore, in the Itô sense, the functionals $`\delta _n`$ and $`\varphi _n`$ representing the number of assets in the portfolio solely depend on the share price right before time $`t`$, i.e., they do not depend on $`X(t)`$ but on $`X(tdt)=XdX`$. That is,
$$\delta _n(X,t)\delta (XdX,t),$$
(46)
and similarly for $`\varphi _n`$ (recall that all equalities must be understood in the mean square sense explained in Sect. 2).
The expansion of Eq. (46) yields (see Eq. (5))
$$\delta _n(X,t)=\delta (X,t)\frac{\delta (X,t)}{X}dX+O(dt),$$
but from the Itô lema (19) we see that
$$\frac{\delta (X,t)}{X}dX=d\delta (X,t)+O(dt),$$
and finally
$$\delta _n(X,t)=\delta (X,t)d\delta (X,t)+O(dt).$$
(47)
Analogously,
$$\delta (X,t)=\delta _n(X,t)+d\delta _n(X,t)+O(dt),$$
(48)
and a similar expresion for $`\varphi (X,t)`$.
As to the self-financing strategy, Eq. (29), we observe that $`\delta (X,t+dt)`$ is the number of shares we have at time $`t+dt`$, while $`\delta (XdX,t)`$ is that number at time $`t`$. Therefore,
$$X(t)d\delta (XdX,t)=[\delta (X,t+dt)\delta (XdX,t)]X(t)$$
is the money we need or obtain for buying or from selling shares at time $`t`$. Analogously, $`B(t)d\varphi (XdX,t)`$ is the money, needed or obtained at time $`t`$, coming from bonds. If we follow a self-financing strategy, both quantities are equal but with different sign, i.e.,
$$X(t)d\delta (XdX,t)=B(t)d\varphi (XdX,t)$$
(49)
which agrees with Eq. (29).
## Appendix B The differential of the option price
Let us derive the differential of the call price, $`dC`$, using either Itô and Stratonovich interpretations. The starting point for both derivations is the replicating portfolio, Eq. (27),
$$C(X,t)=X(t)\delta _n(X,t)+B(t)\varphi _n(X,t).$$
(50)
Taking into account the Itô product rule Eq. (11), we have
$`dC=[\delta _n(X,t)+d\delta _n(X,t)]dX`$ $`+`$ $`[\varphi _n(X,t)+d\varphi _n(X,t)]dB`$
$`+`$ $`X(t)d\delta _n(X,t)+B(t)d\varphi _n(X,t),`$
which, after using Eq. (48), reads
$`dC=\delta (X,t)dX+\varphi (X,t)dB`$ $`+`$ $`X(t)d\delta _n(X,t)`$
$`+`$ $`B(t)d\varphi _n(X,t)+O(dt^{3/2}),`$
and this agrees with Eq. (28).
Within the Stratonovich calculus, the differential of Eq. (50) reads
$$dC=X_S(t)d\delta _n+B(t)d\varphi _n+\delta _n(X_S,t)dX+\varphi _n(X_S,t)dB.$$
(51)
From Eq. (46) we have
$$\delta _n(X_S,t)=\delta (X_S,t)\frac{\delta (X_S,t)}{X_S}dX+O(dX^2),$$
(52)
and analogously for $`\varphi _n`$. Substituting Eq. (52) into Eq. (51), and taking into account Eqs. (4)-(5), (10) and (26) we obtain
$`dC=[X(t)`$ $`+`$ $`dX/2]d\delta _n+B(t)d\varphi _n+\delta (X_S,t)dX`$
$`+`$ $`\left[rB(t)\varphi (X_S,t)\sigma ^2X_S^2{\displaystyle \frac{\delta (X_S,t)}{X_S}}\right]dt+O(dt^{3/2}).`$
But from Eq. (46) and the self-financing strategy (49), we see that $`X(t)d\delta _n+B(t)d\varphi _n=0`$. Hence
$`dC={\displaystyle \frac{1}{2}}dXd\delta _n`$ $`+`$ $`\delta (X_S,t)dX`$ (53)
$`+`$ $`\left[rB(t)\varphi (X_S,t)\sigma ^2X_S^2{\displaystyle \frac{\delta (X_S,t)}{X_S}}\right]dt+O(dt^{3/2}).`$
The substitution of the Stratonovich rule Eq. (20),
$$d\delta _n=\frac{\delta _n(X_S,t)}{X_S}dX+\frac{\delta _n(X_S,t)}{t}dt,$$
yields
$`dC=\delta (X_S,t)dX`$ $`+`$ $`[rB(t)\varphi (X_S,t)`$ (54)
$``$ $`{\displaystyle \frac{1}{2}}\sigma ^2X_S^2{\displaystyle \frac{\delta (X_S,t)}{X_S}}]dt+O(dt^{3/2}),`$
where we have taken into account Eq. (4) and the fact that $`\delta _n/X_S=\delta /X_S+O(dt^{1/2})`$. Eq. (54) agrees with Eq. (38) and the rest of the derivation is identical to that of the main text.
## Appendix C Solution to the Black-Scholes equation
In this appendix we outline the solution to the Black-Scholes equation (35) under conditions (42)-(44).
We first transform Eq. (35) into a forward parabolic equation with constant coefficients by means of the change of variables
$$z=\mathrm{ln}(x/x_C),t^{}=Tt.$$
(55)
We have
$$\frac{C}{t^{}}=rC(z,t^{})+\left(r\frac{1}{2}\sigma ^2\right)\frac{C}{z}+\frac{1}{2}\sigma ^2\frac{^2C}{z^2},$$
(56)
$`(\mathrm{}<z<\mathrm{},\mathrm{\hspace{0.17em}0}<t^{}<T)`$. Moreover, the definition of a new dependent variable:
$$u(z,t^{})=\mathrm{exp}\left[\frac{1}{2}\left(1\frac{2r}{\sigma ^2}\right)z+\frac{1}{8}\sigma ^2\left(1+\frac{2r}{\sigma ^2}\right)(Tt^{})\right]C(z,t^{}),$$
(57)
turns Eq. (56) into the ordinary diffusion equation in an infinite medium
$$\frac{u}{t^{}}=\frac{1}{2}\sigma ^2\frac{^2u}{z^2},$$
(58)
with a constant diffusion coefficient given by $`\sigma ^2/2`$, and initial condition:
$`u(z,0)=x_C\mathrm{exp}[{\displaystyle \frac{1}{2}}(1`$ $``$ $`{\displaystyle \frac{2r}{\sigma ^2}})z`$ (59)
$`+`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(1+{\displaystyle \frac{2r}{\sigma ^2}})T\left]\mathrm{max}\right(e^z1,0).`$
The solution of problem (58)-(59) is standard and reads
$$u(z,t^{})=\frac{1}{\sqrt{2\pi \sigma ^2t^{}}}_{\mathrm{}}^{\mathrm{}}u(y,0)e^{(zy)^2/2\sigma ^2t^{}}𝑑y.$$
(60)
If we substitute the initial condition (59) into the right hand side of this equation and undo the changes of variables we finally obtain the Black-Scholes formula Eq. (45). |
warning/0001/astro-ph0001484.html | ar5iv | text | # A High Sensitivity Measurement of the MeV 𝛾-Ray Spectrum of Cygnus X-1
## 1 Introduction
Cygnus X-1 is generally considered to be one of the most well-established candidates for a stellar-mass black hole. Having been discovered as an X-ray source more than 30 years ago, it has been studied extensively at X-ray and $`\gamma `$-ray energies. At energies approaching 1 MeV, it is one of the brightest sources in the sky. The spectrum at MeV energies is so steep, however, that it has been poorly measured at energies near 1 MeV and above. An accurate characterization of the spectrum at these high energies will facilitate a more complete understanding of the underlying physics of this source. This, in turn, will surely have an impact on our understanding of the $`\gamma `$-ray emission from all black hole sources, including the (stellar mass) soft X-ray transients and the (supermassive) Active Galactic Nuclei (AGN).
It has long been recognized that the soft X-ray emission ($`10`$ keV) generally varies between two discrete levels (e.g., Priedhorsky, Terrell & Holt, 1983; Ling et al., 1983; Liang & Nolan, 1983). Cygnus X-1 seems to spend most ($`90\%`$) of its time in the so-called low X-ray state, characterized by a relatively low flux of soft X-rays and a relatively high flux of hard X-rays ($`100`$ keV). This state is sometimes referred to as the hard state, based on the nature of its soft X-ray spectrum. On occasion, it moves into the so-called high X-ray state, characterized by a relatively high soft X-ray flux and a relatively low hard X-ray flux. This state is sometimes referred to as the soft state, based on the nature of its soft X-ray spectrum. There are, however, some exceptions to this general behavior. For example, HEAO-3 observed, in 1979, a relatively low hard X-ray flux coexisting with a low level of soft X-ray flux (Ling et al., 1983, 1987). Ubertini et al. (1991a) observed a similar behavior in 1987.
Early X-ray spectra of Cygnus X-1 (in its low X-ray state) supported the notion that the high energy emission resulted from the accretion of matter onto a stellar-mass black hole. The emission was generally interpreted to be the result of the Comptonization of a soft thermal photon flux by an energetic electron population. The Comptonization model of Sunyaev & Titarchuk (1980, hereafter ST80) became the standard model for the interpretation of high energy spectra from Cygnus X-1 and other similar sources. This model assumed some (unspecified) source of soft photons interacting in an optically thick Comptonization region ($`\tau 1`$) with nonrelativistic electrons ($`kT_em_ec^2`$). It was successfully used to interpret many of the early hard X-ray observations at energies below $`200`$ keV. The data for Cygnus X-1 indicated that the Comptonization was taking place in a region with an electron temperature ($`kT_e`$) in the 30–60 keV range and an optical depth ($`\tau `$) of 1–5 (e.g., Sunyaev & Trümper, 1979; Steinle et al., 1982; Ubertini et al., 1991b).
As the observations improved, especially at energies extending beyond 200 keV, it became increasingly difficult to model the broad-band continuum spectrum with a single-temperature (or single-component) inverse Compton model (e.g., Nolan et al., 1981; Nolan & Matteson, 1983; Grebenev et al., 1993). Additional spectral components, such as a second Comptonization region, were invoked to improve the spectral fits. Grebenev et al. (1993) suggested that the discrepancy at higher energies could also be explained as a result of the limitations of the ST80 model. In particular, they showed that Monte Carlo simulations of the Comptonization process (which were not restricted by the assumptions of the analytical model) could provide accurate fits to the broad-band GRANAT data extending up to several hundred keV.
Meanwhile, there were several efforts designed to expand upon the fundamentals provided by the ST80 model. Zdziarski (1984, 1985, 1986) considered both bremsstrahlung and synchrotron radiation as soft photon sources. Further analytical improvements to the ST80 Comptonization model were developed by Titarchuk (1994, hereafter, T94). (See also Hua & Titarchuk (1995) and Titarchuk & Lyubarskij (1995).) Incorporating various relativistic corrections, this so-called generalized Comptonization model was applicable over a much wider range of parameter space. The ability of the T94 model to accurately predict the spectrum over a broader range of spectral parameters has been subsequently verified via Monte Carlo simulations by both Titarchuk & Hua (1995) and by Skibo et al. (1995). The simulations of Skibo et al. (1995) included both bremsstrahlung and annihilation radiation as soft photon sources, in addition to some (unspecified) external soft photon source. They concluded that the ST80 model agreed well with Monte Carlo simulations for $`kT_e200`$ keV and $`\tau 2`$, whereas the T94 model agreed well with Monte Carlo simulations for $`kT_e300`$ keV and $`\tau 0.2`$. These same simulations also demonstrated that, under certain conditions, the analytical model of Zdziarski (1985) could also be used.
Source geometry is also an important factor that must be considered in the interpretation of broad-band spectra. For example, Haardt et al. (1993) argued, based on Monte Carlo studies, that improved fits to broad-band spectra could be achieved by incorporating a reflection component along with the inverse Compton component. The reflection component would result from the Compton scattering of hard X-rays from an accretion disk corona off a cool optically-thick accretion disk. The reprocessing leads to an enhancement of emission in the 10–100 keV energy range. The requirement for such a reprocessed component is also supported by the observation of a fluorescence line from neutral iron (e.g., Ebisawa et al., 1996). Assuming a geometry consisting of an optically thin corona above an optically-thick accretion disk, Haardt et al. (1993) showed that the model was consistent with data from EXOSAT, SIGMA (Salotti et al., 1992) and OSSE (Grabelsky et al., 1993), suggesting a much higher coronal electron temperature ($`kT_e150`$ keV) and an even smaller optical depth ($`\tau 0.3`$) than suggested by Comptonization models alone. The reflection component was also incorporated into an accretion disk corona (ADC) model by Dove et al. (1997). In this case, the geometry consisted of a hot inner spherical corona with an exterior accretion disk. Reasonable comparisons with the broad-band spectrum from 1 keV up to several hundred keV were demonstrated with an electron temperature of $`kT_e90`$ keV and an optical depth of $`\tau 1.5`$. A similar model was used by Gierlinski et al. (1997) to fit a combined Ginga-OSSE spectrum, but they found that, with fixed normalization between the Ginga and OSSE spectra, a second Comptonization component was needed to improve the fit at the highest energies. (It has been pointed out by Poutanen (2000), however, that a single temperature Comptonization model fits the Ginga-OSSE data quite well if the relative normalization between the two spectra is left as a free parameter; see Figure 6 of Poutanen (1998).) Moskalenko, Collmar & Schönfelder (1998) also developed a multi-component model to explain the spectrum from soft X-rays into the $`\gamma `$-ray region. In this case, a central spherical corona surrounds the black hole itself, outside of which is the optically-thick accretion disk, with a much hotter outer corona surrounding the entire system.
The use of two-component Comptonization models to improve spectral fits is motivated by the reasonable assumption that the Comptonization region will not be isothermal. A simple two-temperature model may not be very physical, however, since it implicitly assumes that there are two distinct non-interacting regions in which Comptonization is taking place. A more realistic approach would be to assume a continuum of Comptonization parameters. Skibo & Dermer (1995) argued that a thermally-stratified black hole atmosphere could act to harden the spectrum. They simulated a model involving a high-energy spherical core surrounded by an optically thick accretion disk. If the inner core of the high-energy region were hot enough, then a hard tail extending into the MeV region might be produced. A similar model was used to interpret the first combined spectral results from the BATSE and COMPTEL experiments on CGRO (Ling et al., 1997) . Thermal stratification is also an essential concept in the transition disk model of Misra & Melia (1996). This model involves large thermal gradients in the inner region of an optically thick accretion disk. These gradients represent a transition between the cold optically thick disk and the hot plasma which exists near the inner part of the disk. This model has been used to provide a good fit to spectra in the 2–500 keV energy band (Misra et al., 1997, 1998). Dove et al. (1997) also incorporated thermal stratification into the structure of the inner corona of their model.
The net effect of thermal gradients is to produce a non-Maxwellian electron energy distribution. A high energy tail in the electron distribution leads, via the the inverse Compton process, to a high energy tail in the photon distribution. Several models have sought to explain hard-tail emissions by invoking some kind of nonthermal process to generate a high energy (possibly relativistic) tail to the thermal Maxwellian electron energy distribution. Many of these hybrid thermal/non-thermal models (e.g., Bednarek et al., 1990; Crider et al., 1997; Gierlinski et al., 1997; Poutanen, 1998; Poutanen & Coppi, 1998; Coppi, 1999) tend to be rather ad hoc, in that they assume some accelerated (or non-Maxwellian) particle population and then proceed to explore the subsequent consequences. Typically, although not always, the nonthermal population takes the form of a power-law in the electron energy distribution. Such a distribution is similar to that seen, for example, in solar flares (e.g., Coppi, 1999). Both Crider et al. (1997) and Poutanen & Coppi (1998) have shown that a Maxwellian plus power-law form for the electron energy distribution can be adapted to fit a composite COMPTEL-OSSE spectrum of Cygnus X-1.
Others have considered physical mechanisms by which non-thermal electron distributions might be developed. For example, both stochastic particle acceleration (Dermer, Miller, & Li, 1996; Li, Kusunose & Liang, 1996) and MHD turbulence (Li & Miller, 1997) have been proposed as mechanisms for directly accelerating the electrons. The ion popoulation might also contribute to the non-thermal electron distribution in the case where a two-temperature plasma develops (e.g., Dahlbacka et al., 1974; Shapiro, Lightman and Eardley, 1976; Chakrabarti & Titarchuk, 1995). In this situation, the ion population may reach a temperature of $`kT_i10^{12}`$ K, resulting in $`\pi ^o`$ production from proton-proton interactions (e.g., Eilik, 1980; Eilik & Kafatos, 1983). The $`\pi ^o`$ component then leads, via photon-photon interactions between the $`\pi ^o`$-decay photons and the X-ray photons, to production of energetic (nonthermal) $`e^+e^{}`$ pairs. Jordain & Roques (1994) used this concept to fit the hard X-ray tails of not only Cygnus X-1, but also GRO J0422+32 and GX 339-4, as measured by both SIGMA and OSSE. While retaining the standard ST80 spectrum to explain the emission at energies below 200 keV, they used $`\pi ^o`$ production to generate the nonthermal pairs needed to fit the spectrum at energies above $`200`$ keV.
The history of Cygnus X-1 is also riddled with unconfirmed reports of very intense, very broad line emission in the region around 1 MeV, far exceeding that which would be expected from a simple extrapolation of the low energy continuum spectrum (e.g., Ling et al., 1987; McConnell et al., 1989; Owens & McConnell, 1992). The broad MeV feature observed by HEAO-3 (Ling et al., 1987) occured under an unusual source condition when both the hard and soft X-ray fluxes were low. Liang & Dermer (1988) interpreted the feature as evidence for the presence of a very hot ($`kT_e400`$ keV) pair-dominated cloud in the inner region of the accretion disk. In this model, pairs may escape the system and produce a weak narrow annihilation feature in the cold surrounding medium (Dermer & Liang, 1989). Such a feature was also suggested in the HEAO-3 spectrum (Ling & Wheaton, 1989). Melia & Misra (1993) extended this model by considering a more realistic thermally stratified cloud. These “MeV bumps,” have not been seen by any of the experiments on the Compton Gamma-Ray Observatory (CGRO) (McConnell et al., 1994; Phlips et al., 1996; Ling et al., 1997). Any such emission must therefore be time-variable (c.f., Harris et al., 1993).
Although there have been no recent observations of an “MeV bump” in the spectrum of Cygnus X-1, the continuum flux levels that are now generally observed around 1 MeV still indicate a substantial hardening of the high energy spectrum. The extent to which the spectrum hardens at energies approaching 1 MeV has now become an important issue for theoretical modeling of the spectrum. Data from CGRO, in particular the COMPTEL experiment on CGRO, offer the best opportunity for more precisely defining the highest energy parts of the spectrum. Such measurements are critically important in our efforts to determine the nature of the complete hard X-ray / $`\gamma `$-ray spectrum.
In this paper, we present a composite low X-ray state spectrum of Cygnus X-1 compiled from contemporaneous CGRO observations that spans the energy range from $``$30 keV up to $`5`$ MeV. This spectrum should add new insights into the modeling of the broad band $`\gamma `$-ray spectrum of Cygnus X-1 and especially its emissions just above 1 MeV. In $`\mathrm{\S }`$ 2 we briefly describe the COMPTEL experiment, the data from which have been the focus of this investigation. The observations and data selection are described in $`\mathrm{\S }`$ 3. In $`\mathrm{\S }`$ 4 we discuss the analysis of the COMPTEL data. The analysis is expanded in $`\mathrm{\S }`$ 5 to include data from both OSSE and BATSE. A discussion of the results is then presented in $`\mathrm{\S }`$ 6.
## 2 The COMPTEL Experiment
The COMPTEL experiment images 0.75–30 MeV $`\gamma `$-radiation within a field-of-view of $`1`$ sterdian. It consists of two independent layers of detector modules separated by 150 cm. The upper (D1) layer is composed of 7 independent NE213A liquid scintillator detectors, each 28 cm in diameter and 8.5 cm thick. The lower (D2) layer is composed of 14 independent NaI(Tl) detectors, each 28 cm in diameter and 7.5 cm thick. An event is defined as a coincident interaction in a single D1 detector and a single D2 detector. For each event, the total energy is estimated as the sum of the energy losses in both D1 and D2. The interaction locations in both detectors are determined using the relative responses of the various PMTs which are affixed to each D1/D2 detector. In this manner, the interaction locations can be determined with an uncertainty of about 2.0 cm. The measurement of the time-of-flight (TOF) between D1 and D2 and the pulse-shape (PSD) in D1 allow for the rejection of a large fraction of background events, including both upward-moving and neutron-induced events. A more detailed description of COMPTEL can be found in Schönfelder et al. (1993).
The basic principle of COMPTEL imaging is governed by the physics of Compton scattering. If the energy loss of the Compton scattered electron is completely contained within the D1 module and if the scattered photon is completed absorbed by the D2 module, then we have an accurate measure of both the scattered electron energy and the scattered photon energy. Knowledge about the interaction locations in both the upper (D1) and lower (D2) detector layers provides information about the path of the scattered photon. Without more complete knowledge regarding the direction of the scattered electron, the possible arrival direction of the incident photon is confined to a annular region on the sky. This event circle is defined to have an angular radius ($`\overline{\varphi }`$) determined by the Compton scattering formula:
$$\mathrm{cos}\overline{\varphi }=1\frac{m_ec^2}{E_2}+\frac{m_ec^2}{E_1+E_2},$$
(1)
where $`m_ec^2`$ is the electron rest mass energy and $`E_1`$ and $`E_2`$ are the energy deposits measured in D1 and D2, respectively. The superposition of many event circles can lead to the crude localization of a source by determining the direction in which the majority of the event circles intersect (c.f., Figure 7 of Winkler et al. 1992). In practice, the analysis of COMPTEL data is far more complex. Many instrumental effects complicate this simplified approach. For example, incomplete energy absorption in either D1 or D2 can render equation (1) invalid. Effects such as these are taken into account in the final analysis via an appropriate instrumental point spread function (PSF).
## 3 The Observations
Since the launch of CGRO in April of 1991, several observations of the Cygnus region have been carried out with COMPTEL. In order to assemble a broad-band picture of the $`\gamma `$-ray emission from Cygnus X-1, we set out to combine COMPTEL data with data from the BATSE and OSSE experiments on CGRO. The BATSE experiment is an uncollimated array of NaI scintillation detectors covering $`4\pi `$ steradian and operating in the $``$30 keV – 1.8 MeV energy range. The spectral analysis of a given point source is achieved using techniques that rely on the occultation of the source flux by the Earth. There are currently two approaches to spectral analysis with BATSE. The BATSE experiment team routinely uses a technique that uses data only from a limited time period about each Earth occultation (both ingress and egress; Harmon et al. (1992)). Another approach, developed by Ling et al. (1996), makes use of a more extensive set of data in an effort to model the combined effect of some 65 celestial sources of $`\gamma `$-radiation and various models of the instrumental background. This method of BATSE source analysis is embodied in a software system known as the Enhanced BATSE Occultation Package (EBOP; Ling et al., 1996, 2000). The OSSE experiment is a collimated array of NaI detectors that are used for both on-source and off-source measurements. OSSE operates in the 50 keV – 10 MeV energy range.
In selecting those CGRO observations to use in our broad-band analysis, we required contemporaneous observations by all three instruments (COMPTEL, BATSE and OSSE). The COMPTEL FoV is centered on the CGRO pointing direction (the z-axis) and can generally study sources that are within $`40\mathrm{°}`$ of the pointing direction. OSSE has a more limited FoV of $`11\mathrm{°}`$ by $`4\mathrm{°}`$ and is restricted to pointing directions along the spacecraft X-Z plane. BATSE, on the other hand, can observe Cygnus X-1 continually by means of Earth occultation techniques. The data selection was therefore driven by the availability of contemporaneous COMPTEL and OSSE data. The selection was also confined to the first three cycles of CGRO observations (1991 – 1994), based on the quality of the COMPTEL data that was available at the time this study was undertaken. The initial selection of CGRO observation periods, based on these criteria, is listed in Table 1. In addition to the dates of each observation, the table also gives the COMPTEL viewing angle (the offset angle from the COMPTEL pointing direction) and a crude estimate of COMPTEL’s total effective on-axis exposure to Cyg X-1 (measured in days). The exposure estimate is based on an approximated angular response for COMPTEL and includes estimates for the effects of Earth occultation times and telemetry contact times.
The initial analaysis of the broad-band spectrum produced from the full set of observations in Table 1 led to a significant discrepancy between the BATSE and OSSE spectra (McConnell et al., 1998). This has been largely resolved by a further selection based on the hard X-ray flux. Figure 1 shows the hard X-ray flux, as measured by the BATSE occultation technique, during each of the observations in Table 1. These data, taken from the BATSE web site at Marshall Space Flight Center <sup>1</sup><sup>1</sup>1http://www.batse.msfc.nasa.gov/batse/, include statistical errors only. Also indicated in Figure 1 are the $`\gamma _0`$, $`\gamma _1`$, and $`\gamma _2`$ flux levels as defined by Ling et al. (1987, 1997). In general, the hard X-ray flux varied between the $`\gamma _1`$ and $`\gamma _2`$ levels. However, the hard X-ray flux was considerably below-average (at the $`\gamma _0`$ level) during Viewing Period (VP) 318.1 and, to a lesser extent, during VP 331.5. Due to the way that the data were collected, these low hard X-ray flux observations were weighted differently by the various CGRO instruments. These different weightings were directly responsible for the discrepancies noted during our initial analysis. For the final analysis reported here, we excluded the data from VP 318.1 and VP 331.5. The remaining data corresponded to a relatively constant hard X-ray flux of 0.1 photons cm<sup>-2</sup> s<sup>-1</sup> in the 45–100 keV energy band. It is assumed, based on the consistent level of hard X-ray flux, that the corresponding soft X-ray flux was also consistent during these observations and that these data all correspond to the canonical low X-ray state of Cygnus X-1. During the first few months of the CGRO mission (from the launch in April of 1991 until October of 1991) all-sky monitoring data from Ginga (1–20 keV) is available that confirms our assumption that Cygnus X-1 was in its low state (Kitamoto et al., 2000). For the period from October of 1991 until December of 1995 (when RXTE was launched), the data archives at the High Energy Astrophysics Science Archive Research Center (HEASARC) <sup>2</sup><sup>2</sup>2http://heasarc.gsfc.nasa.gov/ show only sporadic pointed observations by ASCA or ROSAT, none of which correspond precisely to any of the COMPTEL observation times.
## 4 COMPTEL Data Analysis
The analysis of COMPTEL data can be logically divided into a spatial (or imaging) analysis and a spectral analysis. In practice, however, the spatial and spectral analysis of the data is inextricably linked via the PSF and its dependence on the incident photon spectrum. The goal of the spatial analysis is to define, for a given range of photon energy loss values, the corresponding photon intensity distribution on the sky. The final source spectrum is then derived, in a self-consistent manner, from the results of a spatial analysis in several distinct energy loss intervals.
### 4.1 Event Selections and the COMPTEL Dataspace
The spatial analysis is performed by first selecting data within some range of measured (total) energy loss values. Within the chosen energy range, events are carefully selected in order to reduce the background contributons and to maximize the signal-to-noise. These event selections include the following: a) restrictions on D1 pulse shape (PSD) to select only those events consistent with incident photons; b) a TOF selection to reject all but “forward” scattered photon events; c) a scatter angle ($`\overline{\varphi }`$) selection to restrict the analysis to a range of values which is dominated by source events rather than by background events; and d) a selection to reject any event whose event circle passes within $`10\mathrm{°}`$ of the Earth’s disk (thus minimizing the background of earth albedo $`\gamma `$-rays).
Once the event data have been selected, the subsequent analysis is carried out within a three-dimensional dataspace that is defined by the fundamental quantities that represent each COMPTEL event. The first two of these parameters, the angles that are arbitarily referred to as $`\chi `$ and $`\psi `$, define the direction of the photon that is scattered from D1 to D2. The third parameter defining this dataspace is the Compton scatter angle ($`\overline{\varphi }`$), as estimated from equation (1) using the measured energy losses in both D1 and D2.
For a point source, the distribution of events in the 3-dimensional ($`\chi `$, $`\psi `$, $`\overline{\varphi }`$) dataspace is generally contained within the interior of a cone whose apex corresponds to the direction of the source. This distribution corresponds to the PSF of COMPTEL. The details of the PSF depend not only on the energy range of the analysis, but also on the shape of the assumed input spectrum, especially as it extends to higher energies. In the present analysis, we have derived PSFs from Monte Carlo simulations (Stacy et al., 1996). The current uncertainties in the PSF are estimated to contribute a systematic error of not more than $``$15–20% to the flux uncertainties. This error is dominated by uncertainties in the physical modeling of COMPTEL and not by Monte Carlo statistics.
### 4.2 COMPTEL Spatial Analysis (Imaging)
The derivation of an image from binned COMPTEL event data (binned into the three-dimensional dataspace) starts by folding an assumed source distribution through the instrumental response (PSF). The PSF incorporates all of the various physical processes alluded to in $`\mathrm{\S }`$ 2 as well as the various event selections which have been imposed on the data. The known exposure and geometric factors (which are specific to a given observation and include, for example, the on/off status of the various modules) are then included. Finally, a ‘background’ model, which may consist of several components, is added in. This process results in an estimate of the event distribution in the three-dimensional COMPTEL dataspace, which can be compared directly with the real data. Subsequent iterations of this process lead to a more precise estimate of the source distribution on the sky.
In practice, the imaging analysis of the COMPTEL data is performed in one of two ways. One technique employs the maximum entropy method to derive a source distribution on the sky (Strong et al., 1992; Schönfelder et al., 1993). As presently implemented for COMPTEL data analysis, this algorithm does not provide quantitative error estimates. A more quantitative analysis can be made using a maximum likelihood technique (de Boer et al., 1992; Schönfelder et al., 1993). This approach compares the relative probability of a model which contains only background to the probability of a model which contains both the background plus a single point source at the given location (the likelihood ratio). This method provides quantitative information regarding both the source location and the flux together with their associated errors.
The background model used to generate an image is a crucial component of the analysis. These background models consists of components which describe both the (internal) instrumental background and the (external) sky background. The term ’background’, in this case, refers to all photon sources other than the source of interest (Cygnus X-1). To date, the most successful approach for estimating the instrumental background involves a smoothing technique which suppresses point-source signals while preserving the general background structure (Bloemen et al., 1994). The sky background component typically includes both known (or suspected) point sources and diffuse sources (such as the galactic diffuse emission). The maximum likelihood analysis provides a best-fit normalization factor (and associated uncertainty) for each background component.
The location of Cygnus X-1 in galactic coordinates ($`l=71.3\mathrm{°},b=3.1\mathrm{°}`$) places it in a rather complex region. The background modeling therefore included a number of different models for the spatial distribution of the celestial photon emission. These models corresponded either to established models for celestial $`\gamma `$-ray emission (e.g., models based on the distribution of atomic and molecular gas within the galaxy) or on ad-hoc photon distributions as derived directly from COMPTEL images. Images were generated using a variety of such models in various combinations. In all cases, a point source model at the location of Cygnus X-1 was included. The resulting distribution of flux values for the Cygnus X-1 point source model provides a direct measure of the errors in the point source analysis, incorporating both the statistical errors from counting statistics and the systematic errors introduced by the background modeling.
For the purposes of background modeling, there are two known sources of $`\gamma `$-ray emission that should be noted. The first of these is the diffuse $`\gamma `$-ray emission from the galaxy. We have made use of a model that is consistent with global studies of COMPTEL data (e.g., Strong et al., 1996; Bloemen et al., 1999, 2000). This model includes estimates of the $`\gamma `$-ray emission which results from the interaction of cosmic rays with both $`HI`$ and $`H_2`$, as well as the contribution from inverse Compton emission off cosmic ray electrons (Strong & Youseffi, 1995; Strong et al., 1996; Strong, Moskalenko, & Reimer, 2000). An $`E^2`$ power law spectrum is assumed for the diffuse $`\gamma `$-ray spectrum. A second known source of $`\gamma `$-ray emission is the 39.5 msec pulsar PSR 1951+32, which lies only $`2.8\mathrm{°}`$ from Cygnus X-1 (at galactic coordinates $`l=68.8\mathrm{°},b=2.8\mathrm{°}`$). Although not readily apparent in spatial studies, a timing analysis of COMPTEL data integrated over the full 750 keV to 30 MeV energy band independently provides evidence for PSR 1951+32 (Kuiper et al., 1998). Its close proximity to Cygnus X-1 makes this source an important component of the background models.
A sample of COMPTEL imaging data is shown in Figure 2, where we present a maximum likelihood map derived from data integrated over the energy loss range of 0.75 – 2.0 MeV. The contours represent constant values of the quantity $`2\mathrm{ln}\lambda `$, where $`\lambda `$ is the likelihood ratio. In a search for single point sources, $`2\mathrm{ln}\lambda `$ has a chi-square distribution with 3 degrees of freedom. (For instance, a $`3\sigma `$ detection corresponds to $`2\mathrm{ln}\lambda `$ = 13.9.) Cygnus X-1 is clearly visible. The likelihood reaches a value of $`2\mathrm{ln}\lambda =93.2`$ at the position of Cygnus X-1, which corresponds to a detection significance of $`9.7\sigma `$. These same data were used to derive the $`1\sigma `$, $`2\sigma `$ and $`3\sigma `$ location confidence contours shown in Figure 3, which demonstrate the ability of COMPTEL to locate the source of emission. In defining constraints on the source location, $`2\mathrm{ln}\lambda `$ has a chi-square distribution with 2 degrees of freedom. So the $`1\sigma `$, $`2\sigma `$ and $`3\sigma `$ location confidence contours correspond to a change in likelihood of 2.3, 6.2, and 11.8, respectively. The contours reflect only the statistical uncertainties; systematic effects are not included.
### 4.3 COMPTEL Spectral Analysis
A photon spectrum of Cyg X-1 was assembled using flux values derived from the spatial (imaging) analysis of COMPTEL data in five distinct energy bands. Since the flux measurements are, at some level, dependent on the instrumental PSF (and the spectral form assumed for that PSF), the resulting spectrum is also dependent on the PSF (and the spectral form assumed for that PSF). The analysis therefore included a careful check on the consistency of the resulting photon spectrum with the spectral form assumed in generating each PSF.
Initial fluxes were derived using PSFs based on an $`E^2`$ power-law spectrum. The resulting photon spectrum was fit with a power-law spectrum of the form,
$$\frac{dN}{dE}=AE^\alpha $$
(2)
The fit gave a photon index of $`\alpha =3.2(\pm 0.4)`$. This suggested the need to use PSFs based on a steeper power-law source spectrum. Subsequent results were derived using PSFs based on an $`E^3`$ power-law spectrum. A power-law fit to this second spectrum gave an index of $`\alpha =3.3(\pm 0.4)`$. From this result, we conclude that, at least within this range of PSFs, the resulting flux values are relatively insensitive to changes in the PSFs. In other words, the derived photon spectrum is not very compliant with respect to the assumed input spectrum. The robust nature of the extracted photon spectrum means that spectral fits in photon dataspace (rather than energy-loss dataspace) can be performed with a high degree of reliability.
The final flux data points from the analysis of the COMPTEL data are given in Table 2. These results were derived using PSFs based on an $`E^3`$ power-law spectrum. The quoted uncertainties include both statistical and systematic errors. The statistical errors are those due to the counting statistics of the measurement. The systematic errors are based on results from using various background models (as described in $`\mathrm{\S }`$ 4.2) and, in some cases, are comparable to the statistical errors, especially at lower energies. (We have not included here the systematic error associated with the PSF calculations, as discussed in $`\mathrm{\S }`$ 4.1, because these are considered to be negligible relative to other sources of error.)
A plot of the COMPTEL spectra, along with the best-fit power-law fit, is shown in Figure 4. The best-fit power-law gives a spectral index ($`\alpha `$) of $`3.3\pm 0.4`$ and an amplitude of $`5.1(\pm 0.8)\times 10^4`$ cm<sup>-2</sup> s<sup>-1</sup> MeV<sup>-1</sup> at 1 MeV. These data provide convincing evidence for emission up to 2 MeV, with a marginal detection in the 2–5 MeV energy band. There is no evidence for emission at energies above 5 MeV. These conclusions are supported by the images generated from these same data. The resulting flux values are consistent with previously published results (e.g., McConnell et al., 1994).
## 5 Broad-Band CGRO Spectral Analysis
Several different models can be used to fit the COMPTEL data, but the lack of low energy data points results in very poor constraints on the model parameters. We therefore must incorporate additional spectral data (at lower energies) in order to more precisely define the broad-band spectrum and hence gain some important insight into the physics of the emission region. The additional data used in this case are the contemporaneous BATSE and OSSE data. As previously discussed, these data have been selected based both on the requirement of contemporaneous observations and on the requirement of a consistent level of hard X-ray flux.
For these broad-band spectra, power-law models clearly do not provide an adequate description of the data. We chose four alternative spectral forms to help describe the data. The first form is an exponentiated power-law,
$$\frac{dN}{dE}=AE^\alpha e^{E/E_c}$$
(3)
This particular function, defined by a power-law index ($`\alpha `$) and a characteristic energy cutoff ($`E_c`$), is not based on any underlying physical model, but its functional form approximates the spectrum at these energies (e.g., Phlips et al., 1996). The data was also fit using both single-component and two-component Comptonization models. In this case, we used the analytical Comptonization model of Titarchuk (1994), which is characterized by a normalization factor, an electron temperature ($`kT_e`$) and an optical depth ($`\tau `$). A fourth functional form is the hybrid thermal/non-thermal model of Poutanen & Svensson (1996) (see also Poutanen, 1998; Poutanen & Coppi, 1998), which describes the photon spectrum resulting from a thermal (Maxwellian) electron distribution plus a non-thermal (power-law) electron distribution at higher energies. The thermal component is described by a characteristic temperature ($`kT_e`$). The non-thermal component is characterized as a power law with spectral slope $`p_e`$ extending from an electron Lorentz factor $`\gamma _{min}`$ (where the Maxwellian transforms to the power-law tail) up to a Lorentz factor of $`\gamma _{max}=1000`$. The electron population is assumed to reside in a accretion disk corona with an optical depth of $`\tau `$. This particular model is useful in that it permits a quantitative description of the underlying electron distribution.
We have already demonstrated that the reproduced COMPTEL spectrum is relatively insensitive to the shape of the assumed spectrum, thus giving us some level of confidence in fitting these data in photon-space. Here we also make the same assumption with regards to both the BATSE and OSSE spectra. Given the steepness of the measured spectra, this assumption is likely to be a safe one, at least to first-order. This approach greatly simplifies the analysis using multiple observations from multiple experiments.
The initial fits to the combined datasets were performed over the 50 keV to 5 MeV energy range. A significant amount of scatter in the OSSE and BATSE data at energies below 200 keV led to rather poor chi-square values for the resulting fits. In order to obtain improved chi-square values, the final model fits were derived from the data between 200 keV and 5 MeV. In addition to the reduced scatter of the BATSE and OSSE data, the general agreement of these two datasets was also much improved at energies above 200 keV. Below 200 keV, there exists discrepancies of up to 25% between the BATSE and OSSE data, whereas at energies above 200 keV the spectra were found to agree quite closely, with differences typically less than 5%. This range is also above the range where many models predict the presence of a backscatter component from photons scattering off a cooler optically thick accretion disk. Spectral fits were derived by allowing both the BATSE and OSSE normalizations to be free parameters. We found that constraining the COMPTEL normalization to that of either BATSE or OSSE gave consistent results for the physical parameters of interest. The final fit parameters were derived by constraining the COMPTEL normalization to that of OSSE.
A summary of the broad-band fit results is given in Table 3. The broad-band photon spectrum, along with the various spectral fits, is shown in Figure 5. The same data is also shown in Figure 6, but now plotted in terms of $`E^2`$ times the photon flux; this plot shows that the power peaks around 100 keV and also accentuates the differences between the spectra and the various model fits. In addition, we also plot (in both Figures 4 and 5) an estimated upper limit from EGRET data collected during Cycles 1–4 of the CGRO mission, an exposure which is similar to that for the other data included in this analysis. This upper limit is based on data from Hartman et al. (1999) and assumes an $`E^3`$ source spectrum. These data do not provide any evidence for a cutoff to the high energy power-law.
## 6 Discussion
We have assembled a broad-band $`\gamma `$-ray spectrum of Cygnus X-1 using data collected from COMPTEL, BATSE and OSSE during the first three years of the CGRO mission. The data were collected contemporaneously and selected so as to have a common level of hard X-ray flux. The hard X-ray flux during these observations varied between the $`\gamma _1`$ and $`\gamma _2`$ levels of Ling et al. (1987, 1997). Although there is poor coverage at soft X-ray energies during this time frame, data from both BATSE (Ling et al., 1997) and OSSE (Phlips et al., 1996) indicate that, for the included observations (Table 1), the hard X-ray / $`\gamma `$-ray spectrum remained fairly stable. In particular, the spectral form is that of the “breaking $`\gamma `$-ray state”, as defined by Grove et al. (1998), which is generally associated with the low X-ray state. (McConnell et al. (2000) compare the spectrum presented here with that obtained during a high X-ray state observation in 1996, showing evidence for spectral variability at MeV energies.)
The COMPTEL data provide evidence of significant emission out to at least 2 MeV. There is additional, but perhaps less compelling, evidence for emission in the 2–5 MeV energy band, with no evidence for emssion above 5 MeV. An analysis of the COMPTEL data alone does not adequately constrain the nature of the broad-band $`\gamma `$-ray spectrum. A more complete interpretation of the COMPTEL data requires the consideration of the spectrum at lower energies. Here we have used contemporaneous data from both BATSE and OSSE to define the spectrum at lower energies.
As noted previously (McConnell et al., 1994), these data do not provide any evidence for an “MeV bump”. The spectrum above 750 keV can be described as a power-law with a photon spectral index of $`\alpha `$ = -3.2. There is no evidence for a break in this power-law. An estimated upper limit from EGRET data does not constrain the extrapolated power-law. An important goal of future measurements will therefore be the determination of the break energy of this power-law spectrum.
Broad-band spectral fits, at energies above 200 keV, show that both the single-component Comptonization model and the exponentiated power-law model do not provide a good fit to the data at energies above 1 MeV, although they do provide very good fits to the BATSE-OSSE data alone. The COMPTEL data dictate the need for an additional spectral component. The high energy fit is significantly improved using a two-component Comptonization model. Unfortunately, the analytical model (T94) begins to fail in this case, because the derived electron temperature of the high-temperature component exceeds the range that is generally allowed for this model (Skibo et al., 1995). Nonetheless, the improvements that result from the two-component Comptonization model suggest that perhaps a stratified Comptonization region (providing a range of both electron temperatures and optical depths) would be more appropriate for modeling the spectrum. Ling et al. (1997) reached the same conclusions based on Monte Carlo modeling of BATSE spectra combined with (non-contemporaneous) COMPTEL data (McConnell et al., 1994). Physically, this is a much more appealing concept than that of a multi-component Comptonization model in which there is no interaction between isothermal Comptonization regions. Thermal gradients are incorporated into several models, all of which lead to the generation of a high energy tail (e.g., Skibo & Dermer, 1995; Chakrabarti & Titarchuk, 1995; Misra & Melia, 1996; Ling et al., 1997).
Hybrid thermal/nonthermal models may also be a viable possibility for explaining the observed $`\gamma `$-ray spectrum. The existence of such distributions is clearly established in the case of solar flares (e.g., Coppi, 1999) and it is therefore natural to expect that similar distributions exist elsewhere in the universe. Several such models have been discussed in the literature (e.g., Crider et al., 1997; Gierlinski et al., 1997; Poutanen & Svensson, 1996; Poutanen, 1998; Poutanen & Coppi, 1998; Coppi, 1999). Fits to our broad-band spectrum using the model of Poutanen & Svensson (1996) provide a quantitative estimate of the electron distribution. In this case, the spectral data can be explained by a thermal Maxwellian distribution with an electron temperature of $`kT_e=86`$ keV, coupled to a power-law electron spectrum with a power-law index of $`p_e=4.5`$. The transition between the Maxwellian and the power-law occurs at an electron kinetic energy of $`570`$ keV ($`\gamma _{min}=2.12`$). The electron population is embedded in an accretion disk corona with an optical depth of $`\tau =1.63`$. The derived spectral index for the power-law tail is somewhat harder (4.5 versus 3.2) but still consistent with that derived by Crider et al. (1997).
The role of $`\pi ^o`$-production can also not be ruled out as a contribution to the measured spectrum. As initially proposed by Jordain & Roques (1994), this mechanism was used to explain the inadequacy of the ST80 model at energies below 1 MeV. It has not been applied to Cygnus X-1 at energies above 1 MeV, which would be required to test this model with the present results. Models for advection-dominated accretion flows (ADAF) predict ion temperatures of $`10^{12}`$ K in the inner part of the accretion flow (e.g., Chakrabarti & Titarchuk, 1995), suggesting the possibility of pion production (although it is not clear whether an ADAF is present during the low X-ray state of Cygnus X-1).
The spectrum presented here clearly indicates the need for a non-Maxwellian electron energy distribution. In particular, it strongly suggests the presence of a high energy tail to that distribution. Whether this results from thermal gradients or from a more isothermal system coupled with a non-thermal component remains an open question. The shape of the electron distribution and its high energy tail can only be determined by measurements that extend into the MeV energy region.
There have been several studies of the broad-band hard X-ray emission from Cygnus X-1 based, in part, on OSSE data. These studies have concentrated on joint observations with lower energy experiments, such as GINGA, ASCA or RXTE. Since these observations are of relatively short duration, these spectra typically do not exhibit a well-defined hard tail due to limited statistics near 1 MeV. Nonetheless, the presence of a hard tail in the continuum emission near 1 MeV now appears to be well-established. An excess above the standard Comptonization models has been discussed several times in the literature for both Cygnus X-1 (e.g., McConnell et al., 1994; Phlips et al., 1996; Ling et al., 1997) and for GRO J0422+32 (van Dijk et al., 1995). The data presented here provide the best quantitative measurement of the spectrum of Cygnus X-1 at these energies and thus provide the best opportunity to study the nature of this hard tail emission. The IBIS and SPI instruments on INTEGRAL will provide only a marginal improvement in the continuum sensitivity over that of CGRO (Lichti et al., 1996; Ubertini et al., 1996). In addition, with their much smaller FoV, the INTEGRAL instruments, unlike COMPTEL, may be severely limited in terms of their total exposure to Cygnus X-1. COMPTEL may therefore provide the best available data on the MeV spectrum of Cygnus X-1 for many years to come. The continued analysis of data from CGRO (only a fraction of the total COMPTEL data are used here) may help to further clarify the nature of this high energy emission and perhaps enable us to more fully elucidate the physics of the accretion process around stellar-mass black holes.
The authors would like to acknowledge Juri Poutanen and Andrzej Zdziarski for comments on the original manuscript. Juri Poutanen kindly provided the XSPEC version of his hybrid thermal/non-thermal model (Poutanen & Svensson, 1996) that was used in this anlaysis. The COMPTEL project is supported by NASA under contract NAS5-26645, by the German government through DLR grant 50 Q 9096 8 and by the Netherlands Organization for Scientific Research NWO. This work has also been supported at UNH by the CGRO Guest Investigator Program under NASA grant NAG5-7745. |
warning/0001/math0001167.html | ar5iv | text | # Some cyclic covers of complements of arrangements
## 1. Introduction
Let $`𝒜`$ be a hyperplane arrangement in $`^{\mathrm{}}`$, with complement $`M(𝒜)=^{\mathrm{}}_{H𝒜}H`$. The cohomology of $`M(𝒜)`$ with coefficients in a local system arises in a number of applications, both outside and inside arrangement theory. Included among the former are the Aomoto-Gelfand theory of multivariable hypergeometric integrals \[AK, Ge\], and the representation theory of Lie algebras and quantum groups and solutions of the Knizhnik-Zamolodchikov differential equation in conformal field theory \[Va\]. This note is motivated by one of the latter applications, the cohomology of the Milnor fiber of a central arrangement.
Let $`𝒞`$ be a central arrangement of hyperplanes in $`^{\mathrm{}+1}`$, an arrangement for which each hyperplane $`H𝒞`$ contains the origin. For each such hyperplane, let $`\alpha _H`$ be a linear form with kernel $`H`$. Then $`Q=Q(𝒞)=_{H𝒞}\alpha _H`$ is a defining polynomial for the arrangement $`𝒞`$, and is homogeneous of degree equal to the cardinality of $`𝒞`$. The complement $`M(𝒞)=^{\mathrm{}+1}Q^1(0)`$ may be realized as the total space of the global Milnor fibration
$$F(𝒞)\stackrel{}{}M(𝒞)\stackrel{𝑄}{}^{},$$
where $`F(𝒞)=Q^1(1)`$ is the Milnor fiber of $`Q`$, see \[Mi\]. We shall refer to $`F(𝒞)`$ as the Milnor fiber of $`𝒞`$, and write $`F=F(𝒞)`$ when the arrangement $`𝒞`$ is understood.
Suppose that the cardinality of $`𝒞`$ is $`n`$ and let $`(x_0,x_1,\mathrm{},x_{\mathrm{}})`$ be a choice of coordinates on $`^{\mathrm{}+1}`$. The geometric monodromy, $`h:FF`$ of the Milnor fibration is given by $`h(x_0,\mathrm{},x_{\mathrm{}})=(\xi _nx_0,\mathrm{},\xi _nx_{\mathrm{}})`$, where $`\xi _n=\mathrm{exp}(2\pi \mathrm{i}/n)`$. Since $`h`$ has finite order $`n`$, the algebraic monodromy $`h^{}:H^q(F;)H^q(F;)`$ is diagonalizable and the eigenvalues of $`h^{}`$ belong to the set of $`n`$-th roots of unity. Denote the cohomology eigenspace of $`\xi _n^k`$ by $`H^q(F;)_k`$ and write $`b_q(F){}_{k}{}^{}=dim_{}H^q(F;)_k`$.
It is known \[CS1\] that these cohomology eigenspaces are isomorphic to the cohomology of the complement of a decone of $`𝒞`$ with coefficients in certain complex rank one local systems. See Section 2 for a summary of these results, and see \[OT\] as a general reference on arrangements. Let $`𝒜`$ be a decone of $`𝒞`$, an affine arrangement in $`^{\mathrm{}}`$, and denote the rank one local systems arising in the context of the Milnor fiber by $`_k`$, $`1kn`$. Then $`H^{}(F(𝒞);)_kH^{}(M(𝒜);_k)`$ for each $`k`$. Furthermore, the local systems $`_k`$ are rational in the sense of \[CO\]. The results of this work, in the context of the Milnor fiber problem, yield combinatorial bounds
$$dim_{}H^q(A_{}(𝒜),a)dim_{}H^q(F;){}_{k}{}^{}\mathrm{rank}__rH^q(A__r(𝒜),\overline{a}),$$
where $`r=n/(k,n)`$ and $`A_R(𝒜)`$ is the Orlik-Solomon algebra of $`𝒜`$ with coefficients in the ring $`R`$ equipped with appropriate differential. See Section 5 for details. The lower bounds are well known. The upper bounds are new.
The local system $`_n`$ is trivial, and thus corresponds to the constant coefficient cohomology of $`M(𝒜)`$. This is well understood in terms of the Orlik-Solomon algebra. While pursuing the remaining cases, we were led to a family of cyclic covers of $`M(𝒜)`$, which includes the Milnor fiber $`F(𝒞)`$. In this note, we show how a number of known results on the Milnor fiber extend naturally to all members of this family of covers, and give an explicit and elementary proof of the polynomial periodicity of the Betti numbers of the members of this family.
## 2. Milnor Fibration
Recall from the Introduction that $`𝒞`$ is a central arrangement of $`n`$ hyperplanes in $`^{\mathrm{}+1}`$, with coordinates $`(x_0,x_1,\mathrm{},x_{\mathrm{}})`$. Associated with $`𝒞`$, we have the defining polynomial $`Q=Q(𝒞)`$, the complement $`M(𝒞)=^{\mathrm{}+1}Q^1(0)`$, and the Milnor fiber $`F=F(𝒞)=Q^1(1)`$. The geometric monodromy $`h:FF`$ of the Milnor fibration has order $`n=|𝒞|`$. The cyclic group $`_n`$ generated by the geometric monodromy $`h`$ acts freely on $`F`$. This free action gives rise to a regular $`n`$-fold covering $`FF/(_n)`$.
Consider the Hopf bundle $`^{\mathrm{}+1}\{0\}^{\mathrm{}}`$ with projection map $`(x_0,x_1,\mathrm{},x_{\mathrm{}})(x_0:x_1:\mathrm{}:x_{\mathrm{}})`$ and fiber $`^{}`$. Let $`p:M(𝒞)M^{}`$ denote the restriction of this projection to $`M`$, where $`M^{}=p(M)`$. The restriction $`p__F:FM^{}`$ of the Hopf bundle to the Milnor fiber is the orbit map of the free action of the geometric monodromy $`h`$ on $`F`$ and we therefore have $`F/(_n)M^{}`$. These spaces and maps fit together with the Milnor fibration in the following diagram:
$$\begin{array}{ccccc}_n& & F& \stackrel{p__F}{}& F/_n\\ & & & & & & \\ ^{}& & M(𝒞)& \stackrel{p}{}& M^{}\\ & & Q& & \\ ^{}/_n& \stackrel{}{}& ^{}\end{array}$$
Note that $`M^{}`$ is the complement of the projective hypersurface defined by the homogeneous polynomial $`Q`$. Thus it is the complement of the projective quotient of the arrangement $`𝒞`$. If we designate one of its hyperplanes, $`H_{\mathrm{}}`$, the hyperplane at infinity, the remaining arrangement is called the decone of $`𝒞`$ with respect to $`H_{\mathrm{}}`$. We call this $`\mathrm{}`$-arrangement $`𝒜`$ here and observe that $`M(𝒜)=M^{}`$ is independent of the choice of $`H_{\mathrm{}}`$ and $`|𝒜|=n1`$. We shall assume that $`𝒜`$ contains $`\mathrm{}`$ linearly independent hyperplanes.
The cohomology groups of the Milnor fiber have been studied extensively, see for instance \[OR, CS1, Ma, CS2, De\]. We summarize some known results from \[CS1\]. Since $`h`$ has finite order $`n`$, the algebraic monodromy $`h^{}:H^q(F;)H^q(F;)`$ is diagonalizable and the eigenvalues of $`h^{}`$ belong to the set of $`n`$-th roots of unity. Denote the cohomology eigenspace of $`\xi _n^k=\mathrm{exp}(2\pi \mathrm{i}k/n)`$ by $`H^q(F;)_k`$, and denote the characteristic polynomial of $`h^{}:H^q(F;)H^q(F;)`$ by $`\mathrm{\Delta }_q(t)=det(th^{}\mathrm{id})`$.
###### Proposition 2.1 (\[CS1, 1.1\]).
Let $`\xi _n^j`$ and $`\xi _n^k`$ be two $`n`$-th roots of unity which generate the same cyclic subgroup of $`_n=\xi _n`$. Then, for each $`q`$, the cohomology eigenspaces $`H^q(F;)_j`$ and $`H^q(F;)_k`$ are isomorphic.
###### Corollary 2.2.
For each $`q`$, $`0q\mathrm{}`$, there are nonnegative integers $`d_{k,q}`$ so that
$$\mathrm{\Delta }_q(t)=\underset{k|n}{}\mathrm{\Phi }_k(t)^{d_{k,q}},$$
where $`\mathrm{\Phi }_k(t)`$ denotes the $`k`$-th cyclotomic polynomial.
###### Theorem 2.3 (\[CS1, 1.6\]).
Define the rank one local system $`_k`$ on $`M(𝒜)`$ by the representation $`\tau _k:\pi _1(M(𝒜))^{}`$ given by $`\gamma _H\xi _n^k`$ for each meridian loop $`\gamma _H`$ about the hyperplane $`H𝒜`$. Then, for each $`k`$, $`1kn`$, we have
$$H^{}(F(𝒞);)_kH^{}(M(𝒜);_k).$$
In light of this result, it is natural to study the local system on $`M(𝒜)`$ induced by the representation given by $`\gamma _H\xi _m^k`$ for arbitrary $`m`$. In subsequent sections, we focus on the context in which these local systems arise.
## 3. Cyclic Covers
The realization of the Milnor fiber of a central arrangement as a cover of the complement of a decone in the previous section motivates the following construction.
Let $`𝒜`$ be an affine arrangement in $`^{\mathrm{}}`$, with coordinates $`𝐱=(x_1,\mathrm{},x_{\mathrm{}})`$. Associated with $`𝒜`$, we have the defining polynomial $`f=Q(𝒜)`$, and the complement $`M(𝒜)=^{\mathrm{}}f^1(0)`$. For each positive integer $`m`$, let $`g_m:^{}^{}`$ denote the cyclic $`m`$-fold covering defined by $`g_m(z)=z^m`$, and let $`p_m:X_m(𝒜)M(𝒜)`$ denote the pullback of $`g_m`$ along the map $`f:M(𝒜)^{}`$, where
$$X_m(𝒜)=\{(𝐱,z)M(𝒜)\times ^{}f(𝐱)=z^m\}.$$
This family of cyclic covers of $`M(𝒜)`$ generalizes the Milnor fiber in a number of ways which we now pursue. First, we have the following.
###### Proposition 3.1.
Let $`𝒞`$ be a central arrangement of $`n`$ hyperplanes in $`^{\mathrm{}+1}`$, with defining polynomial $`Q=Q(𝒞)`$ and Milnor fiber $`F(𝒞)`$. If $`𝒜`$ is a decone of $`𝒞`$, then the covering spaces $`p__F:F(𝒞)M(𝒜)`$ and $`p_n:X_n(𝒜)M(𝒜)`$ are equivalent.
###### Proof.
The relation between the defining polynomials $`f`$ of $`𝒜`$ and $`Q`$ of $`𝒞`$ is given by $`Q=x_0^nf(x_1/x_0,\mathrm{},x_n/x_0)`$, and $`F(𝒞)=Q^1(1)`$. Using this, it is readily checked that the map $`X_n(𝒜)F(𝒞)`$ defined by $`(x_1,\mathrm{},x_{\mathrm{}},z)(1/z,x_1/z,\mathrm{},x_{\mathrm{}}/z)`$ is a homeomorphism inducing an equivalence of covering spaces. ∎
The characteristic homomorphism $`\mathrm{\Phi }:\pi _1(M(𝒜))_n`$ of the covering $`p__F:F(𝒞)M(𝒜)`$ was identified in \[CS1, 1.2\]. It is given by $`\mathrm{\Phi }(\gamma _H)=g_n`$, where $`\gamma _H`$ is a meridian loop about the hyperplane $`H𝒜`$ and $`g_n`$ is a fixed generator of $`_n`$. A straightforward generalization of the proof of this fact from \[CS1, 1.2\] yields
###### Proposition 3.2.
The characteristic homomorphism $`\mathrm{\Phi }_m:\pi _1(M(𝒜))_m`$ of the covering $`p_m:X_m(𝒜)M(𝒜)`$ is given by $`\mathrm{\Phi }_m(\gamma _H)=g_m`$ for a fixed generator $`g_m`$ of $`_m`$ and meridian loops $`\gamma _H`$ about the hyperplanes $`H`$ of $`𝒜`$.
The covering spaces $`X_m(𝒜)`$ fit together nicely in the sense of the following.
###### Proposition 3.3.
If $`m=kr`$, then the map $`^{\mathrm{}}\times ^{}^{\mathrm{}}\times ^{}`$ defined by $`(𝐱,z)(𝐱,z^r)`$ induces a cyclic $`r`$-fold covering $`p_{m,k}:X_m(𝒜)X_k(𝒜)`$.
###### Proof.
Let $`X_{k,r}(𝒜)X_k(𝒜)`$ denote the pullback of $`g_r:^{}^{}`$ along the map $`X_k(𝒜)^{}`$ defined by $`(x_1,\mathrm{},x_{\mathrm{}},z)z`$, with
$$X_{k,r}(𝒜)=\{(𝐱,z,w)M(𝒜)\times ^{}\times ^{}f(𝐱)=z^k\text{and}z=w^r\}.$$
It is then readily checked that the map $`X_m(𝒜)X_{k,r}(𝒜)`$ defined by $`(𝐱,z)(𝐱,z^r,z)`$ is a homeomorphism compatible with the projection maps. ∎
###### Remark 3.4.
The space $`X_m(𝒜)`$ admits a self-map $`h_m:X_m(𝒜)X_m(𝒜)`$ defined by $`h_m(𝐱,z)=(𝐱,\xi _m^1z)`$, where $`\xi _m=\mathrm{exp}(2\pi \mathrm{i}/m)`$. In the case $`m=n`$ of the Milnor fiber, the map $`h_n:X_n(𝒜)X_n(𝒜)`$ corresponds to the geometric monodromy $`h:F(𝒞)F(𝒞)`$ under the equivalence of covering spaces exhibited in the proof of Proposition 3.1. For arbitrary $`m`$, the “monodromy” map $`h_m`$ generates a cyclic group of order $`m`$, which acts freely on $`X_m(𝒜)`$. The resulting regular $`m`$-fold covering $`X_m(𝒜)X_m(𝒜)/h_m`$ clearly coincides with $`p_m:X_m(𝒜)M(𝒜)`$. More generally, for $`m=kr`$ composite, the map $`h_m^k`$ generates a cyclic group of order $`r`$, which also acts freely on $`X_m(𝒜)`$, and the covers $`X_m(𝒜)X_m(𝒜)/h_m^k`$ and $`p_{m,k}:X_m(𝒜)X_k(𝒜)`$ coincide.
## 4. Cohomology
We now study the cohomology of the covering spaces $`X_m(𝒜)`$. Fix a basepoint $`𝐱_0M(𝒜)`$. From the Leray-Serre spectral sequence of the fibration $`p_m:X_m(𝒜)M(𝒜)`$, we obtain $`H^{}(X_m(𝒜);)=H^{}(M(𝒜);^m)`$, the cohomology of $`X_m(𝒜)`$ with trivial $``$-coefficients is isomorphic to the cohomology of the base $`M(𝒜)`$ with coefficients in the rank $`m`$ local system $`^m`$ with stalk $`_𝐱^m=H^0(p_m^1(𝐱);)^m`$. The results presented in this section are natural generalizations in the context of arrangements of those of \[CS1, Section 1\].
###### Proposition 4.1 (cf. \[CS1, 1.3–1.5\]).
Let $`T\mathrm{GL}(m,)`$ be the cyclic permutation matrix of order $`m`$ defined by $`T(\stackrel{}{e}_i)=\stackrel{}{e}_{i+1}`$ for $`1in1`$ and $`T(\stackrel{}{e}_n)=\stackrel{}{e}_1`$, where $`\{\stackrel{}{e}_i\}`$ is the standard basis for $`^m`$. Note that $`T`$ is diagonalizable with eigenvalues $`\xi _m^k`$, $`1km`$.
1. The local system $`^m`$ is induced by the representation $`\tau ^m:\pi _1(M(𝒜),𝐱_0)\mathrm{GL}(m,)`$ given by $`\tau ^m(\gamma _H)=T`$ for each meridian $`\gamma _H`$.
2. The local system $`^m`$ decomposes into a direct sum, $`^m=_{k=1}^m_k^m`$ of rank one local systems. For each $`k`$, the local system $`_k^m`$ is induced by the representation $`\tau _k^m:\pi _1(M(𝒜),𝐱_0)^{}`$ defined by $`\tau _k^m(\gamma _H)=\xi _m^k`$.
3. We have $`H^{}(X_m(𝒜);)=_{k=1}^mH^{}(M(𝒜);_k^m)`$.
The above result provides one decomposition of the cohomology $`H^{}(X_m(𝒜);)`$. Another is given by the monodromy maps $`h_m:X_m(𝒜)X_m(𝒜)`$ of Remark 3.4. Since $`h_m`$ has finite order $`m`$, the induced map $`h_m^{}:H^q(X_m(𝒜);)H^q(X_m(𝒜);)`$ is diagonalizable, with eigenvalues among the $`m`$-th roots of unity. Denote the cohomology eigenspace of $`\xi _m^k`$ by $`H^q(X_m(𝒜);)_k`$, and let $`\mathrm{\Delta }_q^{(m)}(t)=det(th_m^{}\mathrm{id})`$ denote the characteristic polynomial of $`h_m^{}:H^q(X_m(𝒜);)H^q(X_m(𝒜);)`$. We then have the following generalizations of Theorem 2.3, Proposition 2.1, and Corollary 2.2.
###### Proposition 4.2 (cf. \[CS1, 1.6\]).
For each $`k`$, $`1km`$, we have
$$H^{}(X_m(𝒜);)_kH^{}(M(𝒜);_k^m)$$
###### Proposition 4.3 (cf. \[CS1, 1.1\]).
Let $`\xi _m^j`$ and $`\xi _m^k`$ be two $`m`$-th roots of unity which generate the same cyclic subgroup of $`_m=\xi _m`$. Then the cohomology eigenspaces $`H^{}(X_m(𝒜);)_j`$ and $`H^{}(X_m(𝒜);)_k`$ are isomorphic.
###### Corollary 4.4.
If $`\xi _m^j`$ and $`\xi _m^k`$ are $`m`$-th roots of unity which generate the same cyclic subgroup of $`_m=\xi _m`$, then $`H^{}(M(𝒜);_j^m)H^{}(M(𝒜);_k^m)`$.
###### Corollary 4.5.
For each $`q`$, $`0q\mathrm{}`$, there are nonnegative integers $`d_{k,q}^{(m)}`$ so that
$$\mathrm{\Delta }_q^{(m)}(t)=\underset{k|n}{}\mathrm{\Phi }_k(t)^{d_{k,q}^{(m)}},$$
where $`\mathrm{\Phi }_k(t)`$ denotes the $`k`$-th cyclotomic polynomial.
We relate the cohomology of the spaces $`X_m(𝒜)`$ for various $`m`$ using these results.
###### Theorem 4.6.
If $`k`$ divides $`m`$, then the cohomology $`H^{}(X_k(𝒜);)`$ is a direct summand of $`H^{}(X_m(𝒜);)`$.
###### Proof.
From Proposition 4.1, we have $`H^{}(X_m(𝒜);)=_{q=1}^mH^{}(M(𝒜);_q^m)`$, where the local system $`_q^m`$ is induced by the representation $`\tau _q^m`$ given by $`\gamma _H\xi _m^q`$. Writing $`m=kr`$, we see that the representations $`\tau _p^k`$, $`1pk`$, are among the $`m`$ representations $`\tau _q^m`$, $`1qm`$. In other words, $`_{p=1}^kH^{}(M(𝒜);_p^k)=H^{}(X_k(𝒜);)`$ is a direct summand of $`H^{}(X_m(𝒜);)`$. ∎
###### Remark 4.7.
This result may be interpreted in terms of the intermediate coverings $`p_{m,k}:X_m(𝒜)X_k(𝒜)`$ of Proposition 3.3 as follows. One can check that the projection $`p_{m,k}`$ commutes with the monodromy maps, $`p_{m,k}h_m=h_kp_{m,k}`$. Consequently, the induced map $`p_{m,k}^{}:H^{}(X_k(𝒜);)H^{}(X_m(𝒜);)`$ preserves the eigenspaces of $`h_k^{}`$. Using Proposition 4.2 and the above theorem, one can show that $`p_{m,k}^{}`$ maps $`H^{}(X_k(𝒜);)=_{p=1}^kH^{}(M(𝒜);_p^k)`$ isomorphically to the summand $`_{p=1}^kH^{}(M(𝒜);_{pr}^{kr})`$ of $`H^{}(X_m(𝒜);)`$.
These results also show that, to determine the cohomology of $`X_m(𝒜)`$, it suffices to compute $`H^{}(M(𝒜);_1^k)`$ for divisors $`k`$ of $`m`$. Proposition 4.2 and Corollary 4.4 yield
###### Theorem 4.8.
The Betti numbers of the space $`X_m(𝒜)`$ are given by
$$b_q(X_m(𝒜))=dim_{}H^q(X_m(𝒜);)=\underset{k|m}{}\varphi (k)b_q(_1^k),$$
where $`\varphi `$ is the Euler phi function and $`b_q(_1^k)=dim_{}H^q(M(𝒜);_1^k)`$.
The summand $`H^{}(M(𝒜);_1^1)`$ of $`H^{}(X_m(𝒜);)`$ corresponds to the constant coefficient cohomology of $`M(𝒜)`$. This is well understood in terms of the Orlik-Solomon algebra defined next.
## 5. Orlik-Solomon Algebra
Let $`A=A(𝒜)`$ be the Orlik-Solomon algebra of $`𝒜`$ generated by the 1-dimensional classes $`a_H`$, $`H𝒜`$. It is the quotient of the exterior algebra generated by these classes by a homogeneous ideal, hence a finite dimensional graded $``$-algebra. There is an isomorphism of graded algebras $`H^{}(M(𝒜);)A(𝒜)`$. In particular, $`dimA^q(𝒜)=b_q(𝒜)`$ where $`b_q(𝒜)=dimH^q(M(𝒜);)`$ denotes the $`q`$-th Betti number of $`M(𝒜)`$ with trivial local coefficients $``$. The absolute value of the Euler characteristic of the complement is a combinatorial invariant:
(1)
$$\beta (𝒜)=(1)^{\mathrm{}}\underset{q=0}{\overset{\mathrm{}}{}}(1)^qb_q(𝒜)=|\chi (M(𝒜))|.$$
Let $`𝝀=\{\lambda _HH𝒜\}`$ be a collection of complex weights. Define a differential $`A^qA^{q+1}`$ by multiplication by $`a_𝝀=_{H𝒜}\lambda _Ha_H`$. This provides a complex $`(A^{},a_𝝀)`$. Associated to $`𝝀`$, we have a rank one representation $`\rho :\pi _1(M(𝒜))^{}`$ given by $`\gamma _Ht_H=\mathrm{exp}(2\pi \mathrm{i}\lambda _H)`$ for any meridian loop $`\gamma _H`$ about the hyperplane $`H𝒜`$, and a corresponding rank one local system $``$ on $`M(𝒜)`$. Note that $`\rho `$ and $``$ are unchanged if we replace the weights $`𝝀`$ with $`𝝀+𝐦`$, where $`𝐦=\{m_HH𝒜\}`$ is a collection of integers. The following inequalities are well known, see \[CO\].
###### Proposition 5.1.
For all $`𝛌`$ and all $`q`$ we have
$$\underset{𝐦^{|𝒜|}}{sup}dim_{}H^q(A^{},a_{𝝀+𝐦})dim_{}H^q(M(𝒜);)dim_{}H^q(M(𝒜);).$$
For the local systems arising in the context of the covers $`X_m(𝒜)`$, the weights are rational. Suppose $`\lambda _H=k_H/N`$ for all $`H`$, with integers $`k_H`$ and $`N`$, and assume without loss that the g.c.d. of the $`k_H`$ is prime to $`N`$. In this case there are better upper bounds. The Orlik-Solomon ideal is defined by integral linear combinations of the generators, hence the algebra may be defined over any commutative ring $`R`$, denoted $`A_R(𝒜)`$. Write $`A_{}=A_{}(𝒜)`$. Left-multiplication by the element $`a_𝝀=\lambda _Ha_HA_{}^1`$ induces a differential on the Orlik-Solomon algebra, and we denote the resulting complex by $`(A_{}^{},a_𝝀)`$. Similarly, associated to the element $`a_𝐤=Na_𝝀=k_Ha_H`$, we have the complex $`(A_{}^{},a_𝐤)`$. We showed in \[CO\] that the complexes $`(A_{}^{},a_𝝀)`$ and $`(A_{}^{},a_𝐤)`$ are chain equivalent. The coefficients of $`a_𝐤`$ are integers, so we consider the Orlik-Solomon algebra with integer coefficients and the associated complex $`(A_{}^{},a_𝐤)`$. Let $`(A_N^{},\overline{a}_𝐤)`$ be the reduction of $`(A_{}^{},a_𝐤)modN`$, where $`A_N=A__N`$ denotes the Orlik-Solomon algebra with coefficients in the ring $`_N`$ and $`\overline{a}_𝐤=a_𝐤modN`$.
###### Theorem 5.2 (\[CO, 4.5\]).
Let $`𝛌=𝐤/N`$ be a system of rational weights, and let $``$ be the associated rational local system on the complement $`M`$ of $`𝒜`$. Then, for each $`q`$,
$$dim_{}H^q(M(𝒜);)\mathrm{rank}__NH^q(A_N^{},\overline{a}_𝐤).$$
There are examples in \[CO\] which show that the inequality can be strict.
## 6. Polynomial Periodicity
We continue the study of the cohomology of the spaces $`X_m(𝒜)`$. We first investigate the implications of a well known vanishing theorem of Schechtman, Terao, and Varchenko in this context. We then establish the polynomial periodicity of the Betti numbers of this family of spaces.
An edge of $`𝒜`$ is a nonempty intersection of hyperplanes and $`L(𝒜)`$ is the set of edges. Given $`YL(𝒜)`$, let $`𝒜_Y=\{H𝒜YH\}`$. Define the weight of $`Y`$ by $`\lambda _Y=_{H𝒜_Y}\lambda _H`$. Call $`YL(𝒜)`$ dense if $`\beta ((𝒜_Y)_0)>0`$ where $`(𝒜_Y)_0`$ is a decone of $`𝒜_Y`$. The projective closure, $`𝒜_{\mathrm{}}`$, of $`𝒜`$ adds the infinite hyperplane, $`H_{\mathrm{}}`$, with weight $`_{H𝒜}\lambda _H`$. Recall that $`𝒜^{\mathrm{}}`$ contains $`\mathrm{}`$ linearly independent hyperplanes.
###### Theorem 6.1 (\[STV, 4.3\]).
Call the local system $``$ nonresonant if $`\lambda _Y_0`$ for every dense edge $`YL(𝒜_{\mathrm{}})`$. In this case
$$H^q(M(𝒜);)=0\text{ for }q\mathrm{},\text{and}dim_{}H^{\mathrm{}}(M(𝒜);)=\beta (𝒜).$$
Recall the decomposition $`H^{}(X_m(𝒜);)=_{q=1}^mH^{}(M(𝒜);_q^m)`$ of the cohomology of $`X_m(𝒜)`$ from Proposition 4.1. By Theorem 4.8, it suffices to consider the case $`q=1`$. In the notation of the previous section, the local system $`_1^m`$ arises from the rational and equal weights $`\lambda _H=1/m`$ for all $`H𝒜`$. The projective closure, $`𝒜_{\mathrm{}}`$, of $`𝒜`$ is the projective quotient of the cone $`𝒞`$ of $`𝒜`$. Assign the weight $`|𝒜|/m`$ to $`H_{\mathrm{}}`$.
###### Proposition 6.2.
If either (i) $`m>|𝒜|`$, or (ii) $`|(𝒜_{\mathrm{}})_Y|`$ is relatively prime to $`m`$ for every dense edge $`YL(𝒜_{\mathrm{}})`$, then the local system $`_1^m`$ is nonresonant.
###### Proof.
If $`YH_{\mathrm{}}`$, then $`\lambda _Y<0`$. Otherwise, we have $`\lambda _Y=|𝒜_Y|/m`$, which cannot be a positive integer in either case (i) or (ii). ∎
Call a positive integer $`k`$ nonresonant if $`k`$ satisfies either of the hypotheses of Proposition 6.2. Recall the factorization, $`\mathrm{\Delta }_q^{(m)}(t)=_{k|m}\mathrm{\Phi }_k(t)^{d_{k,q}^{(m)}}`$, of the characteristic polynomial of the monodromy $`h_m^{}:H^q(X_m(𝒜);)H^q(X_m(𝒜);)`$ provided by Corollary 4.5. The results of Section 4 and this section provide the following information concerning the exponents $`d_{k,q}^{(m)}`$ arising in this factorization.
###### Proposition 6.3.
For every $`m`$, we have $`d_{1,q}^{(m)}=b_q(𝒜)`$ for all $`q`$. If $`k`$ is nonresonant, then $`d_{k,q}^{(m)}=0`$ if $`q<\mathrm{}`$ and $`d_{k,\mathrm{}}^{(m)}=\beta (𝒜)`$.
###### Proof.
The first statement follows from Proposition 4.2 and the fact that $`_1^1`$ is the trivial local system. The second statement follows from Proposition 6.2. ∎
Proposition 6.2 also facilitates an elementary and explicit proof of the polynomial periodicity of the Betti numbers of the family of covering spaces $`X_m(𝒜)`$. We refer to Sarnak and Adams \[SA, Ad\] for results along these lines in greater generality, and to Hironaka \[H1, H2\] and Sakuma \[Sk\] for related results on branched covers of surfaces and links. A sequence $`\{a_m\}_m`$ is said to be polynomial periodic if there are polynomials $`p_1(x),\mathrm{},p_N(x)[x]`$ so that $`a_m=p_i(m)`$ whenever $`mimodN`$.
###### Theorem 6.4.
For each $`q`$, $`0q\mathrm{}`$, the sequence, $`\{b_q(X_m(𝒜))\}_m`$, of Betti numbers of the cyclic covers $`X_m(𝒜)`$ of the complement $`M(𝒜)`$ is polynomial periodic.
###### Proof.
First note that $`b_0(X_m(𝒜))=1`$ for all $`m`$.
Let $`N=_{p\text{prime}}p^e`$ be the product of all prime powers $`p^e`$ for which the exponent $`e`$ is maximal so that $`p^e|𝒜|`$. Evidently, $`N`$ is the smallest positive integer for which $`k|N`$ for all $`k|𝒜|`$. Note also that if $`mimodN`$, then
(2)
$$\{k1k|𝒜|\text{and}k|m\}=\{k1k|𝒜|\text{and}k|i\}.$$
For $`1q\mathrm{}1`$ and $`1iN`$, define constant polynomials $`p_{q,i}=b_q(X_i(𝒜))`$. From Theorem 4.8, we have $`p_{q,i}=_{k|i}\varphi (k)b_q(_1^k)`$. Similarly, if $`mimodN`$, then $`b_q(X_m(𝒜))=_{k|m}\varphi (k)b_q(_1^k)`$, and by Proposition 6.2, the sum is over all $`k|𝒜|`$. Thus, polynomial periodicity of the Betti numbers $`b_q(X_m(𝒜))`$ for $`1q\mathrm{}1`$ follows from the relation between such divisors of $`m`$ and $`i`$ noted in (2) above.
The polynomial periodicity of the top Betti number $`b_{\mathrm{}}(X_m(𝒜))`$ may be established by an Euler characteristic argument as follows. We have $`\chi (X_m(𝒜))=m\chi (M(𝒜))`$. This, together with (1), yields
$$b_{\mathrm{}}(X_m(𝒜))=m\beta (𝒜)+(1)^{\mathrm{}+1}\left[1+\underset{q=1}{\overset{\mathrm{}1}{}}(1)^qb_q(X_m(𝒜))\right].$$
Defining linear polynomials $`p_{\mathrm{},i}(x)=\beta (𝒜)x+(1)^{\mathrm{}+1}\left[1+_{q=1}^\mathrm{}1(1)^qp_{q,i}\right]`$ for each $`i`$, $`1iN`$, we have $`b_{\mathrm{}}(X_m(𝒜))=p_{\mathrm{},i}(m)`$ if $`mimodN`$. ∎
###### Remark 6.5.
The polynomial periodicity of the Betti numbers of more general classes of covers of a finite CW-complex are established in \[SA, Ad\]. Noteworthy in the above proof are the explicit identifications of the “period” $`N`$ and the polynomials $`p_{q,i}(x)`$ for the cyclic covers $`X_m(𝒜)`$, see the concluding remarks in \[H1\].
A generating function for the Betti numbers $`b_q(X_m(𝒜))`$ is given by the following zeta function, suggested by A. Adem. For each $`q`$, $`0q\mathrm{}`$, define
(3)
$$\zeta _{𝒜,q}(s)=\underset{m=1}{\overset{\mathrm{}}{}}\frac{b_q(X_m(𝒜))}{m^s}.$$
###### Theorem 6.6.
We have
$$\zeta _{𝒜,q}(s)=\zeta (s)\left[\underset{k|𝒜|}{}\frac{\varphi (k)b_q(_1^k)}{k^s}+\delta _{q,\mathrm{}}\beta (𝒜)\underset{k>|𝒜|}{}\frac{\varphi (k)}{k^s}\right]$$
where $`\zeta (s)`$ is the classical Riemann zeta function and $`\delta _{q,\mathrm{}}`$ is the Kronecker delta.
###### Proof.
From Theorem 4.8, we have $`b_q(X_m(𝒜))=_{k|m}\varphi (k)b_q(_1^k)`$. If $`k>|𝒜|`$, we have $`b_q(_1^k)=0`$ for $`q<\mathrm{}`$ and $`b_{\mathrm{}}(_1^k)=\beta (𝒜)`$ by Proposition 6.2. A calculation using these observations yields the result. ∎
## 7. Bounds and Examples
The results of Sections 4 and 6 show that, to determine the cohomology of $`X_m(𝒜)`$, it suffices to compute $`H^{}(M(𝒜);_1^k)`$ for those divisors $`k`$ of $`m`$ for which $`k<|𝒜|`$ and hypothesis (ii) of Proposition 6.2 fails. Combining the results of Proposition 5.1 and Theorem 5.2, we have combinatorial bounds on the local system Betti numbers,
(4)
$$\underset{𝐦^{|𝒜|}}{sup}dim_{}H^q(A^{},a_{𝝀+𝐦})b_q(_1^k)\mathrm{rank}__kH^q(A__k^{},\overline{a}_\mathrm{𝟏}),$$
where $`k𝝀=\mathrm{𝟏}`$ and $`a_\mathrm{𝟏}=_{H𝒜}a_H`$. Evidently, if the two extreme non-negative integers in the above inequalities are equal, then $`b_q(_1^k)`$ is determined. In particular, if $`\mathrm{rank}__kH^q(A__k^{},\overline{a}_\mathrm{𝟏})=0`$, we have $`b_q(_1^k)=0`$ as well. We conclude with several examples which illustrate the utility of these bounds.
###### Example 7.1.
Let $`𝒞`$ be a realization of the MacLane ($`8_3`$) configuration, with defining polynomial $`Q(𝒞)=xy(yx)z(zx\xi _3^2y)(z+\xi _3y)(zx)(z+\xi _3^2x+\xi _3y)`$, and let $`𝒜`$ be a decone of $`𝒞`$. The Poincaré polynomial of $`𝒜`$ is $`P(𝒜,t)=1+7t+13t^2`$, and $`\beta (𝒜)=7`$. A calculation in the Orlik-Solomon algebra of $`𝒜`$ reveals that $`H^q(A__k^{},\overline{a}_\mathrm{𝟏})=0`$ for $`q2`$ and all $`k>1`$. Thus, $`P(X_m(𝒜),t)=1+7t+(6+7m)t^2`$ for all $`m`$. In particular, for $`m=8=|𝒞|`$, the Poincaré polynomial of the Milnor fiber of the MacLane arrangement is $`P(F(𝒞),t)=1+7t+62t^2`$.
###### Example 7.2.
Let $`𝒜`$ be the Selberg arrangement in $`^2`$, with defining polynomial $`Q(𝒜)=xy(xy)(x1)(y1)`$. The Poincaré polynomial of $`𝒜`$ is given by $`P(𝒜,t)=_{q0}b_q(𝒜)t^q=1+5t+6t^2`$, and we have $`\beta (𝒜)=2`$, see (1). The dense edges of $`𝒜_{\mathrm{}}`$ all have cardinality $`3`$, so by Proposition 6.2, if $`k>5`$ or $`k`$ is prime to $`3`$, the local system $`_1^k`$ on $`M(𝒜)`$ is nonresonant. Thus if $`(m,3)=1`$, the Poincaré polynomial of the cover $`X_m(𝒜)`$ is $`P(X_m(𝒜),t)=1+5t+(4+2m)t^2`$.
If $`k=3`$, one can check that $`dim_{}H^q(A^{},a_{𝝀+𝐦})=\mathrm{rank}__3H^q(A__3^{},\overline{a}_\mathrm{𝟏})=1`$, where $`3𝝀=\mathrm{𝟏}`$ and $`𝐦=(\mathrm{}m_H\mathrm{})^5`$ satisfies $`m_H=1`$ if $`H=\{xy=0\}`$ and $`m_H=0`$ otherwise. Consequently, $`b_1(_1^3)=1`$ as well, and if $`3`$ divides $`m`$, we have $`P(X_m(𝒜),t)=1+7t+(6+2m)t^2`$. It follows that the zeta function $`\zeta _{𝒜,1}(s)`$ of (3) is given by $`\zeta _{𝒜,1}(s)=\zeta (s)[5+23^s]`$.
Since the Selberg arrangement is a decone of the braid arrangement $``$ of rank three, the cover $`X_6(𝒜)`$ is homeomorphic to the Milnor fiber $`F()`$, and $`P(F(),t)=1+7t+18t^2`$, as is well known. For further calculations along these lines, see \[CS2, De\].
###### Example 7.3.
Let $`𝒞`$ be the Hessian configuration, with defining polynomial $`Q(𝒞)=x_1x_2x_3_{i,j=0,1,2}(x_1+\xi _3^ix_2+\xi _3^jx_3)`$, and let $`𝒜`$ be a decone of $`𝒞`$ with Orlik-Solomon algebra $`A`$. The Poincaré polynomial of $`𝒜`$ is $`P(𝒜,t)=1+11t+28t^2`$, and $`\beta (𝒜)=18`$. One can check that $`H^q(A__k^{},\overline{a}_\mathrm{𝟏})=0`$ for $`q2`$ if $`k2,4`$, and that, if $`k=2,4`$,
$$\mathrm{rank}__kH^q(A__k^{},\overline{a}_\mathrm{𝟏})=\{\begin{array}{cc}2\hfill & \text{if }q=1\text{,}\hfill \\ 20\hfill & \text{if }q=2\text{,}\hfill \\ 0\hfill & \text{otherwise.}\hfill \end{array}$$
So $`b_1(_1^k)=0`$ and $`b_2(_1^k)=18`$ if $`k2,4`$, while $`b_1(_1^k)2`$ and $`b_2(_1^k)20`$ if $`k=2,4`$.
Concerning the lower bound of (4), it is known that the resonance variety $`_1(𝒞)`$ of the Hessian arrangement has a non-local three-dimensional component $`S^{12}`$, see \[CS3, 5.8\], \[Li, 3.3\]. For $`𝝀S`$, we have $`dim_{}H^1(A^{}(𝒞),a_𝝀)=2`$, see \[Fa, 3.12\]. For an appropriate ordering of the hyperplanes of $`𝒞`$, this component has basis
$$\stackrel{}{e}_5+\stackrel{}{e}_7+\stackrel{}{e}_{12}\stackrel{}{e}_1\stackrel{}{e}_2\stackrel{}{e}_3,\stackrel{}{e}_4+\stackrel{}{e}_9+\stackrel{}{e}_{11}\stackrel{}{e}_1\stackrel{}{e}_2\stackrel{}{e}_3,\stackrel{}{e}_6+\stackrel{}{e}_8+\stackrel{}{e}_{10}\stackrel{}{e}_1\stackrel{}{e}_2\stackrel{}{e}_3.$$
Using this basis, one can show that, for $`a_𝝀=\frac{1}{k}a_\mathrm{𝟏}^{12}`$, there exists $`𝐦^{12}`$ so that $`𝝀+𝐦S`$ if and only if $`k=2,4`$. Thus, $`sup_{𝐦^{12}}dim_{}H^1(A^{}(𝒞),a_{𝝀+𝐦})=2`$ for these $`k`$. A standard argument then shows that $`sup_{𝐦^{11}}dim_{}H^1(A^{}(𝒜),a_{𝝀+𝐦})=2`$ for $`a_𝝀=\frac{1}{k}a_\mathrm{𝟏}^{11}`$ as well. Consequently, the inequalities in the upper bounds on $`b_q(_1^k)`$ noted above are in fact equalities. It follows that the zeta function $`\zeta _{𝒜,1}(s)`$ is given by $`\zeta _{𝒜,1}(s)=\zeta (s)[11+22^s+44^s]`$.
These calculations determine the Betti numbers of the cover $`X_m(𝒜)`$ for any $`m`$, as well as the dimensions of the eigenspaces of the maps $`h_m^{}:H^{}(X_m(𝒜);)H^{}(X_m(𝒜);)`$. In particular, for $`m=12=|𝒞|`$, the Poincaré polynomial of the Milnor fiber of the Hessian arrangement is $`P(F(𝒞),t)=1+17t+232t^2`$, and the characteristic polynomials $`\mathrm{\Delta }_q(t)=\mathrm{\Delta }_q^{(12)}`$ of the algebraic monodromy are
$$\mathrm{\Delta }_0(t)=t1,\mathrm{\Delta }_1(t)=(t1)^9(t^41)^2,\text{and}\mathrm{\Delta }_2(t)=(t1)^8(t^41)^2(t^{12}1)^{18}.$$
###### Example 7.4.
Let $`𝒜`$ be the Ceva(3) arrangement in $`^3`$, with defining polynomial $`Q(𝒜)=(x^3y^3)(x^3z^3)(y^3z^3)`$. It is known that $`dim_{}H^1(A^{},a_𝝀)1`$ for all $`𝝀`$. In the case $`k=3`$, is is also known that $`b_1(_1^3)=\mathrm{rank}__3H^1(A__3^{},\overline{a}_\mathrm{𝟏})=2`$. This example is discussed in detail in \[Fa, 4.5\], \[Li, 3.3\], \[CS3, 6.2\], and \[CO, 3.5 and 4.7\].
As illustrated by this arrangement, the lower bound of (4) may be strict. On the other hand, we know of no example where the upper bound of (4) is strict. |
warning/0001/cond-mat0001037.html | ar5iv | text | # Determining the Wiedemann-Franz ratio from the thermal Hall conductivity: Application to Cu and YBa₂Cu₃O_6.95.
## Abstract
The Wiedemann-Franz ($`WF`$) ratio compares the thermal and electrical conductivities in a metal. We describe a new way to determine its value, based on the thermal Hall conductivity. The technique is applied to copper and to untwinned YBaCuO. In the latter, we uncover a $`T`$-linear dependence and suppression of the Hall-channel $`WF`$ ratio. We discuss the implications of this suppression. The general suppression of the $`WF`$ ratio in systems with predominant electron-electron scattering is discussed.
74.72.-h,72.15.Eb,72.15.Gd,72.15.Lh
The electron fluid in a metal is an excellent conductor of both electric charge and entropy. In familiar metals such as Au, Pb and Cu, the thermal conductivity of the electrons is so large that the phonons account for less than one percent of the heat current at any temperature $`T`$. Wiedemann and Franz ($`WF`$) observed that the ratio of the thermal and electrical conductivities is very nearly the same in many metals . This ratio is expressed as the Lorenz number, which is conveniently written in dimensionless form as
$$=\frac{\kappa _e}{T\sigma }\left(\frac{e}{k_B}\right)^2,$$
(1)
where $`\kappa _e`$ is the electronic thermal conductivity, $`\sigma `$ the electrical conductivity, $`e`$ the elementary charge, and $`k_B`$ is Boltzmann’s constant. In standard transport theory, $``$ equals the Sommerfeld value $`\pi ^2/3`$ if the mean-free-paths ($`mfp`$) are assumed identical for charge and heat transport. Indeed, in all conventional metals, the observed $``$ is rather close to this number above 273 K . However, below $``$273 K, the observed $``$ falls significantly below this value, implying that the heat current is more strongly scattered relative to the charge current. Because $``$ compares directly the charge and entropy currents, it has contributed strongly to our current understanding of how charge and entropy currents are affected by distinct scattering processes in conventional metals .
In solids, the observed thermal conductivity $`\kappa `$ is the sum of the electronic and phonon conductivity $`\kappa _{ph}`$, viz. $`\kappa =\kappa _e+\kappa _{ph}`$. In conventional metals where $`\kappa _{ph}\kappa _e`$, one may use $`\kappa `$ instead of $`\kappa _e`$ to evaluate $``$ with negligible error. However, in conductors with relatively small carrier densities, $`\kappa _{ph}`$ is often much larger than $`\kappa _e`$. These conductors (with resistivities $`\rho `$ exceeding 100 $`\mu \mathrm{\Omega }`$cm) include many interesting conductors such as the cuprates, doped fullerenes, quasicrystals, and newer materials such as $`\mathrm{CaB}_6`$ . Recent theoretical discussions of charge versus entropy currents are most applicable to these conductors. Ironically, their $`WF`$ ratio seems experimentally inaccessible using conventional techniques. Hence, a fresh experimental approach is desirable.
The thermal Hall effect provides a rather efficient way to screen out the phonon heat current (even when it is dominant). In zero field, the observed thermal current $`𝐉_Q`$ is parallel to the temperature gradient $`T`$ (applied $`\widehat{𝐱}`$). A field $`𝐇\widehat{𝐳}`$ generates a Lorentz force that acts only on the electronic component of $`𝐉_Q`$. This produces a transverse Hall current parallel to $`\pm \widehat{𝐲}`$ (depending on the sign of the charge carriers). Thus, the transverse Hall conductivity $`\kappa _{xy}`$ involves only the electrons (the phonon current is strictly unaffected by $`𝐇`$). By forming the ratio in Eq. 1 with $`\kappa _{xy}`$ and the electrical Hall conductivity $`\sigma _{xy}`$, we may measure directly the ‘Hall’ Lorenz number $`_{xy}[\kappa _{xy}/T\sigma _{xy}](e/k_B)^2`$. At sufficiently high $`T`$, $`_{xy}`$ approaches the Sommerfeld value .
We have applied the technique to a detwinned crystal of optimally-doped cuprate $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_{6.95}`$ (YBCO) with a critical temperature $`T_c`$ = 93.3 K, and to elemental Cu. A longitudinal gradient $`_xT`$ ($`2`$ K) is applied $`𝐚`$ of the YBCO crystal, and the transverse gradient $`_yT`$ is measured using a pair of thermocouples (chromel-constantan) versus field ($`𝐇\pm 𝐜`$). The Hall signal $`_HT`$ is the field-odd component of $`_yT`$ . $`\kappa _{xy}`$ is computed as $`[_HT/_xT](\kappa _{xx}\kappa _b/\kappa _a)`$. We measured $`\kappa _a=\kappa _{xx}`$, and used the anisotropy $`\kappa _b/\kappa _a`$ previously measured in another crystal from the same batch (see Ref. ). The copper sample was cut from oxygen-free high-conductivity (OFHC) stock, but was not vacuum annealed to further reduce lattice disorder.
First, we discuss the $`T`$-dependence of $`\kappa _{xy}`$ in Cu, and compare $`_{xy}`$ with the standard Lorenz number (see Ref. for early data on $`\kappa _{xy}`$). Figure 1 displays traces of $`\kappa _{xy}`$ versus $`H`$ between 60 and 350 K measured in Sample 1. In order to calculate $``$ and $`_{xy}`$, we have also measured its electrical conductivity $`\sigma `$ and Hall conductivity $`\sigma _{xy}`$. (Below 50 K, where $`\kappa _{xy}`$ and $`\sigma _{xy}`$ display increasing curvature vs. $`H`$, $`_{xy}`$ is obtained from their values in the limit of weak-fields.)
The values of $`\kappa `$ vs. $`T`$ in Sample 1 are in close agreement with a compilation by Powell et al. . However, in samples that have been carefully annealed, the peak value (at 20 K) may be twice as large \[see Berman and MacDonald (BM)\].
Figure 2 shows that $``$ and $`_{xy}`$ in our sample start at values slightly higher than the Sommerfeld value at 350 K, and decrease with falling $`T`$, with $`_{xy}`$ decreasing slightly faster. In Cu, the Debye temperature $`\theta _D`$ is 343 K. Both Lorenz numbers reach a minimum near 60 K, and increase again at lower $`T`$. For comparison, we show (broken line) $``$ reported by BM. The minimum is deeper and occurs at a lower $`T`$ (20 K). As discussed later, this is consistent with its higher purity.
In conventional metals, electronic currents are limited by scattering of electrons by phonons (at finite $`T`$). Large-angle scattering involving phonons with large wavevectors $`𝐪`$ ($`q>k_F`$, the Fermi wavevector) are equally disruptive of the charge and heat currents. By contrast, small-angle scattering ($`qk_F`$) relaxes only the heat current, leaving the charge current relatively unaffected . In terms of $`\mathrm{}_S`$ and $`\mathrm{}_e`$ (the mean-free-paths for entropy and charge transport, respectively) we may express $``$ as $`(\pi ^2/3)\mathrm{}_S/\mathrm{}_e`$, where $`\mathrm{}`$ denotes averaging over the Fermi Surface $`FS`$. At high temperature ($`T\mathrm{\Theta }_D`$), we have $`\mathrm{}_S\mathrm{}_e`$ because large-angle scattering is dominant. As $`T`$ decreases below $`\mathrm{\Theta }_D`$, the phonon population is increasingly skewed towards the small-$`q`$ limit, so that $`\mathrm{}_e`$ increases relatively faster than $`\mathrm{}_S`$.
This accounts for the decrease in $``$ in Fig. 2. At very low $`T`$, when elastic scattering from impurities dominates, $`\mathrm{}_e`$ and $`\mathrm{}_S`$ are again equal, and $``$ returns to the Sommerfeld value. In our sample, this turnaround occurs near 60 K, whereas in the cleaner sample of BM, it occurs at 20 K.
In comparing $``$ with $`_{xy}`$, we find that both display a minimum at nearly the same $`T`$. However, on the high-$`T`$ side, $`_{xy}`$ falls faster with decreasing $`T`$. The two rates are compared via $`a_L^2/_{xy}`$. The near constancy of $`a_L`$ is consistent with the expectation that $`\mathrm{}_S/\mathrm{}_e`$, whereas $`_{xy}\mathrm{}_S^2/\mathrm{}_e^2`$. Hence, $`_{xy}`$ falls faster because it involves the squares of the $`mfp`$’s. In Cu, the Hall-Lorenz number approaches the same $`WF`$ ratio at high $`T`$, so that it provides information comparable to $``$. Moreover, additional information on scattering processes may be derived from its low $`T`$ behavior.
Thermal conductivity measurements have produced a wealth of information on the superconducting state of the cuprates. In YBCO, they provided early evidence favoring an unconventional pairing symmetry , and long lifetimes for the quasiparticles at low temperatures . The contributions of the chains to $`\kappa `$ , and the effect of hole doping also have been investigated. Above $`T_c`$, however, the problem of estimating $`\kappa _{ph}`$ has been a serious obstacle to the extraction of $`\kappa _e`$. In fact, neither the magnitude nor the $`T`$ dependence of $`\kappa _e`$ in the normal state may be regarded as experimentally established. Hence, the present technique is especially appropriate. While the behavior of $`\kappa _{xy}`$ in YBCO was investigated previously in the superconducting state, the rapid attenuation of the thermal Hall signal precluded measurements above $``$100 K. The improved resolution now allows $`\kappa _{xy}`$ to be determined reliably up to 320 K.
Figure 3 displays the $`H`$-dependence of $`\kappa _{xy}`$ at selected temperatures from 95 to 320 K. It is worth a second remark that, as $`\kappa _{xy}`$ derives no contribution from the phonons, the raw experimental curves directly mirror the electronic heat current. We proceed to compare it with the charge current. In the normal state of YBCO, $`\sigma _{xy}`$ is known to vary as $`1/T^3`$. This produces the anomalous $`T`$ dependence of the Hall coefficient $`R_H`$ that is so characteristic of the cuprates . It is interesting to ask if the same dependence is observed in $`\kappa _{xy}/T`$ (dividing by $`T`$ to remove the heat capacity contribution). As shown in Fig. 4 (open symbols), the $`T`$ depedence of the $`\kappa _{xy}/B`$ is well-fitted to $`T^{1.2}`$. Hence, we find that $`\kappa _{xy}/T`$ actually has a weaker dependence than $`\sigma _{xy}`$.
When we calculate $`_{xy}`$, we find that it varies linearly with $`T`$ (solid circles in Fig. 4). In addition to this unusual $`T`$ dependence, its value from 95 to 320 K is significantly smaller than the Sommerfeld value.
In some previous reports, values of $``$ (often close to the Sommerfeld value) were obtained using a variety of arguments to subtract $`\kappa _{ph}`$ from $`\kappa `$. However, these arguments are suspect for the following reasons. In the cuprates, $`\kappa _{ph}`$ is strongly decreased by lattice disorder associated with dopants ($`\kappa `$ is 2 to 4 times larger in the pristine parent compound $`\mathrm{La}_2\mathrm{CuO}_4`$ compared with doped compounds ). Moreover, within the plane, $`\kappa _{ph}`$ is strongly anisotropic. In both $`\mathrm{YBa}_2\mathrm{Cu}_3\mathrm{O}_7`$ and its two-chain variant $`\mathrm{YBa}_2\mathrm{Cu}_4\mathrm{O}_8`$, $`\kappa _{ph}`$ is significantly larger along $`𝐛`$ (the chain axis) compared with $`𝐚`$, because of the anisotropy in elastic constants . These factors make comparisons between doped and undoped crystals, or between $`\kappa `$ measured along the axes $`𝐚`$ and $`𝐛`$ in the same sample, quite unreliable. A different approach is to extrapolate the value of $`\kappa _e`$ of quasiparticles in the superconducting state to temperatures just above $`T_c`$. A recent experiment reports extrapolated values of $`\kappa _e`$ in the range 0.8-1.0 W/mK at $`T_c`$. With $`\rho _a`$ = 100 $`\mu \mathrm{\Omega }\mathrm{cm}`$, this gives $``$ = 1.1-1.4, which is closer to our $`_{xy}`$ than to the Sommerfeld value.
Aside from the overall suppressed scale, the nominally $`T`$-linear dependence of $`_{xy}`$ is also unusual. An extrapolation (broken line in Fig. 4) shows that it intersects the Sommerfeld value near 500 K. We note that, the acoustic phonons in YBCO have a maximum energy of 20 mv. Above 200 K, the dominant acoustic phonons should have sufficient momenta to span the full $`FS`$ ($`\mathrm{\Theta }_D`$ = 420 K is higher because of the large unit cell). Certainly, by 320 K, we should have $`\mathrm{}_S=\mathrm{}_e`$ if electron-phonon scattering is the dominant mechanism for relaxing the two currents. Hence, the small values of $`_{xy}`$ and the high temperature scale at which it attains the Sommerfeld value (500 K) seem incompatible with dominant electron-phonon scattering.
In contrast, a suppressed $`WF`$ ratio may be expected in systems with dominant electron-electron ($`ee`$) scattering. A discussion of this point illustrates how normal ($`N`$) and Umklapp ($`U`$) scattering processes influence the $`WF`$ ratio. As in the case of lattice thermal conduction , $`N`$ processes leave the total momentum of the electron gas unchanged, so that the charge current cannot relax without $`U`$ processes. However, (unlike lattice conduction) $`N`$-process $`ee`$ scattering does relax the heat current because it causes a redistribution of energy between hot and cold electrons . This distinction implies that systems in which $`ee`$ scattering is dominant have a strongly reduced Lorenz number. Moreover, because the relative weights of $`N`$ versus $`U`$ processes are not determined by $`\mathrm{\Theta }_D`$ in $`ee`$ scattering, this reduction could prevail to very high $`T`$. These issues are intimately related to behavior of the $`WF`$ ratio. Hence, the $`T`$ dependences of $`\kappa _{xy}`$ and $`_{xy}`$ displayed in Fig. 4 should place strong constraints on transport models for YBCO.
In summary, we have described measurements of the Hall-Lorenz number, obtained as the ratio of $`\kappa _{xy}/T`$ to $`\sigma _{xy}`$. In Cu, this method gives results comparable to direct measurements of $``$. However, in systems in which lattice conduction is not small, the Hall-Lorenz experiment provides a direct comparison of the heat and charge currents of the charge carriers. A quantitative understanding of the information derived from $`_{xy}`$ awaits comparison with microscopic calculations. In view of the increased interest in systems with dominant $`ee`$ scattering, we expect high-temperature $`WF`$ ratio to play an increasingly prominent role.
We acknowledge useful discussions with P. W. Anderson, B. Keimer, K. Damle, and S. Uchida. N.P.O. and K.K. acknowledge support by the U.S. Office of Naval Research. N.P.O. acknowledges an award from the International Joint-Research grants from the New Energy and Industrial Tech. Develop. Org. (NEDO). Y.Z. and Z.X. were supported by funds from a DMR-MRSEC award (NSF DMR 9809483) from the U.S. National Science Foundation.
Permanent address of Z.A.X.: Department of Physics, Zhejiang University, Hangzhou, China
Present address of K.K.: Division of Engineering and Applied Sciences, Harvard Universiy, Cambridge MA 02138 |
warning/0001/cond-mat0001230.html | ar5iv | text | # Clusters and Fluctuations at Mean-Field Critical Points and Spinodals
\[
## Abstract
We show that the structure of the fluctuations close to spinodals and mean-field critical points is qualitatively different than the structure close to non-mean-field critical points. This difference has important implications for many areas including the formation of glasses in supercooled liquids. In particular, the divergence of the measured static structure function in near-mean-field systems close to the glass transition is suppressed relative to the mean-field prediction in systems for which a spatial symmetry is broken.
\]
The structure of the fluctuations near critical points and spinodals and their relation to the behavior of quantities observed in experiments and simulations is important for understanding the properties of many materials. For critical phenomena in non-mean-field Ising systems this relation was found by mapping the thermal critical point onto a percolation transition . In these systems the properties of the clusters at the percolation threshold are identical to those of the fluctuations at the thermal critical point. In particular, the mean cluster diameter is the correlation length, the density of the spanning cluster scales as the order parameter, and the mean number of sites in the spanning cluster is the susceptibility .
For mean-field Ising models there is a line of spinodal critical points as well as the usual critical point. These mean-field thermal singularities also can be mapped onto percolation transitions, but the relation between percolation clusters and critical fluctuations is qualitatively different. This difference has important consequences for supercooled mean-field liquids, which also exhibit spinodals, and for “near-mean-field” systems, which are characterized by long, but finite range interactions and which exhibit many of the characteristics of mean-field systems including pseudospinodals. Many dense systems with short-range interactions exhibit near-mean-field behavior under certain circumstances.
In the following we give scaling arguments that relate the structure of clusters at mean-field critical points and spinodals to measurements of thermal quantities such as the static structure function $`S(k)`$ and discuss the implications of these results for understanding mean-field and near-mean-field supercooled liquids. One motivation for this analysis is that the behavior of a supercooled two-component, two-dimensional ($`d=2`$) Lennard-Jones system has been interpreted as due to the influence of a pseudospinodal (relative to the stable solid). A consequence of this pseudospinodal is that the first peak of $`S(k)`$ is predicted to exhibit a power-law divergence, a prediction consistent with simulations in $`d=2`$. However, experiments and simulations have failed to find similar behavior for $`S(k)`$ in $`d=3`$. This failure is puzzling because $`d=3`$ Lennard-Jones systems should be more mean-field-like than $`d=2`$ systems. This work addresses this apparent contradiction.
We first review how the cluster structure relates to thermal critical phenomena in non-mean-field systems. In mapping the percolation transition onto the thermal critical point, the quantity that is isomorphic to the singular part of the free energy is the mean number of clusters in a correlation region of volume $`\xi ^d`$, where $`\xi `$ is the correlation length . The singular part of the free energy scales as $`\xi ^dϵ^{2\alpha }=ϵ^{2\alpha d\nu }`$, where $`ϵ=(TT_c)/T_c`$, and $`\alpha `$ and $`\nu `$ are the specific heat and correlation length exponents respectively. Because hyperscaling is valid, we have $`d\nu =2\alpha `$, and hence the mean number of clusters in a volume $`\xi ^d`$ is order unity.
To obtain the scaling of the isothermal compressibility or susceptibility $`\chi _T`$, we use the relation of $`\chi _T`$ to the spatial integral of $`\mathrm{\Gamma }(r)`$, the order parameter correlation function. We take $`\mathrm{\Gamma }(r)`$ to be the square of the order parameter over a volume $`\xi ^d`$ and zero outside, and obtain $`\chi _Tϵ^{2\beta }\xi ^d=ϵ^{2\alpha d\nu \gamma }`$ using $`\alpha +2\beta +\gamma =2`$. Hyperscaling gives $`2\alpha d\nu =0`$ and hence $`\chi _Tϵ^\gamma `$ as expected. This simple argument is well known, but fails in mean-field systems where hyperscaling is not valid.
In contrast to systems that obey hyperscaling, we need to distinguish between critical phenomena fluctuations and clusters in the mean-field limit. To discuss the latter we consider a weak long-range (Kac) interaction of the form $`\gamma ^du(\gamma r)`$, where $`u(\gamma r)`$ is integrable, a function only of the distance between spins, and ferromagnetic. If the system size is first taken to infinity and then the range of interaction $`R=\gamma ^1\mathrm{}`$, we obtain the Curie-Weiss description of the thermodynamics and the Ornstein-Zernicke form for $`\mathrm{\Gamma }(r)`$. The exponents for the mean-field critical point are $`\nu =1/2,\alpha =0,\beta =1/2`$ and $`\gamma =1`$. The singular part of the free energy at the mean-field critical point scales as
$$\mathrm{\Delta }f\xi ^dϵ^{2\alpha }=R^dϵ^{2d/2}.$$
(1)
The quantity $`R^d`$ appears because all lengths are in units of the interaction range.
From the percolation mapping, Eq. (1) implies that the mean number of clusters in a volume $`\xi ^d`$ is $`R^dϵ^{2d/2}`$. We can estimate the magnitude of this number from the Ginsburg criterion, which states that a system is well approximated by mean-field theory if $`\chi _T/\varphi ^2\xi ^d<<1`$, where $`\varphi ϵ^\beta `$ is the order parameter. If we use mean-field values for the critical exponents, we obtain the condition $`R^dϵ^{2d/2}>>1`$ for mean-field theory to be valid. This inequality, together with Eq. (1) and the percolation mapping, implies that the number of clusters in a volume $`\xi ^d`$ is much greater than unity for mean-field systems. What does this result imply for the nature of clusters in mean-field systems?
In non-mean-field systems which have one cluster in a volume $`\xi ^d`$, the free energy cost of this cluster is $`\xi ^dϵ^{2\alpha }`$. As discussed above, this cost is order unity in non-mean-field systems. In contrast, because $`\mathrm{\Delta }f>>1`$ in mean-field systems (see Eq. (1)), $`ϵ^{2\alpha }`$ is not the free energy density of one cluster in mean-field. Rather, there are $`R^dϵ^{2d/2}`$ clusters in a volume $`\xi ^d`$, and each cluster has a free energy cost of order unity. This free energy implies that the probability of a cluster is order unity.
It also is clear that $`ϵ^\beta `$ is not the density of one cluster, because if it were, the spin density in a volume of order $`\xi ^d`$ would be $`R^dϵ^{2d/2}ϵ^\beta `$. For a fixed value of $`ϵ`$, $`R`$ can be made arbitrarily large and hence this density can be arbitrarily large, an absurd result. This argument and the one for $`\mathrm{\Delta }f`$ implies that $`ϵ^\beta `$ is the density of all the spins in a volume $`\xi ^d`$ regardless of the cluster to which they belong. It also implies that the density of spins in one cluster near the mean-field critical point is
$$\varphi _{\mathrm{cl}}\frac{ϵ^{1/2}}{R^dϵ^{2d/2}}.$$
(2)
This prediction has been verified numerically.
Are these clusters the critical phenomena fluctuations? To answer this question we calculate the density of mean-field fluctuations. The partition function for the Gaussian approximation of the $`\varphi ^4`$ model is given by
$$Z=\delta \varphi \mathrm{exp}\left\{\beta 𝑑\stackrel{}{x}R^d[(\varphi (\stackrel{}{x}))^2+ϵ\varphi ^2(\stackrel{}{x})]\right\},$$
(3)
where $`\beta =1/k_BT`$ and $`k_B`$ is Boltzmann’s constant. Because we are interested only in the scaling properties of the fluctuations, we take the order parameter $`\varphi (\stackrel{}{x})`$ to be a constant $`\varphi `$ over a volume $`\xi ^d`$ and zero outside, and obtain $`Z\delta \varphi e^{\beta R^dϵ^{1d/2}\varphi ^2}`$. If we integrate $`\varphi `$ until the argument of the exponent becomes order unity, we find that the density of the critical phenomena fluctuations scales as
$$\varphi \frac{ϵ^{1/2}}{(R^dϵ^{2d/2})^{1/2}}.$$
(4)
That is, the structures with the density in Eq. (4) have a free energy cost of one. As in the non-mean-field case, we expect these objects to be the critical fluctuations.
From Eqs. (2) and (4) we see that the critical phenomena fluctuations are not the clusters, but are considerably denser in the mean-field limit $`R^dϵ^{2d/2}>>1`$. The structure of the critical phenomena fluctuations in mean-field systems is that the “vacuum” is not featureless but contains a very large number ($`R^dϵ^{2d/2}`$ in a volume $`\xi ^d`$) of clusters. At the critical point (magnetic field $`h=0`$ in Ising models), the mean number of up and down clusters is equal giving a zero mean magnetization. Because these clusters are independent, the fluctuations in the number of clusters is order $`(R^dϵ^{2d/2})^{1/2}`$. These fluctuations make up the zero magnetization background and are the critical phenomena fluctuations. If we use Eq. (4), $`\chi _T\varphi ^2\xi ^dϵR^dϵ^{d/2}/R^dϵ^{2d/2}=ϵ^1`$ as expected.
These arguments can be extended to spinodals in mean-field Ising models. If the spinodal is approached by varying $`h`$ and keeping $`T`$ fixed, the density of the clusters is $`\mathrm{\Delta }h^{1/2}/R^d\mathrm{\Delta }h^{3/2d/4}`$ with $`\mathrm{\Delta }h=h_sh`$. The numerator represents the magnetic field scaling of the order parameter. The denominator is very large in mean-field and is the Ginsburg criterion for spinodals. As at the critical point, the mean number of up and down clusters is equal when the infinite cluster, which is related to the metastable magnetization, is subtracted. All scaling arguments for mean-field critical points apply at spinodals with appropriate changes in the values of the exponents ($`\nu =1/4,\alpha =1/2,\beta =1/2`$, and $`\gamma =1/2`$ ).
We can determine the decay time of the clusters and the critical phenomena fluctuations by constructing an action from the linearized Langevin equation describing the dynamics of an Ising model with long-range interactions. If we assume a random Gaussian noise, the probability measure for the order parameter $`\varphi (\stackrel{}{x},t)`$ is
$`\mathrm{exp}\{\beta {\displaystyle }d\stackrel{}{x}dt[ϵ\left({\displaystyle \frac{\stackrel{~}{\varphi }(\stackrel{}{x},t)}{t}}\right)^2+(MR^dϵ^{2d/2})^2`$ (5)
$`\{^2\stackrel{~}{\varphi }(\stackrel{}{x},t)+ϵ\stackrel{~}{\varphi }(\stackrel{}{x},t)\}^2+H(\overline{\psi },\psi )]\},`$ (6)
where $`M`$ is a (constant) mobility, $`\stackrel{~}{\varphi }=ϵ^{1/2}\varphi `$, and $`H(\overline{\psi },\psi )=\overline{\psi }(\stackrel{}{x},t)[\frac{}{t}+MR^dϵ^{2d/2}(^2+ϵ\})\psi (\stackrel{}{x},t)`$. The variables $`\overline{\psi }(\stackrel{}{x},t)`$ and $`\psi (\stackrel{}{x},t)`$ obey a Grassmann algebra and can be used to convert the average over the Gaussian noise to an average over the order parameter. In the mean-field limit, the Langevin equation describing the dynamics is linear, and the Grassmann fields and the order parameter are independent. Hence, the Grassmann fields are irrelevant to a calculation of averages of functions of $`\varphi `$. Using the action in Eq. (5), we can average a function of $`\varphi `$ by functionally integrating $`\varphi `$ up to its value at the critical point where we expect that the argument of the exponential in Eq. (5) is order one. Because all of the terms in the action in Eq. (5) are real and positive, each term in the exponential must be order one, which implies that $`ϵ\xi ^d\stackrel{~}{\varphi }^2/t1`$. Hence, the time scale that an object lives depends on $`\stackrel{~}{\varphi }`$.
For critical fluctuations, $`\stackrel{~}{\varphi }(R^dϵ^{2d/2})^{1/2}`$, which leads to $`\tau ϵ^1`$, the scaling for the decorrelation time in a mean-field system with a non-conserved order parameter. For clusters, $`\stackrel{~}{\varphi }_{\mathrm{cl}}(R^dϵ^{2d/2})^1`$, which leads to $`\tau _{\mathrm{cl}}ϵ^1/R^dϵ^{2d/2}`$. We stress that these arguments are valid only in the mean-field limit and are a good approximation in the near-mean-field case.
These considerations imply that the clusters have a lifetime that is considerably shorter than the lifetime of the critical phenomena fluctuations. In the mean-field limit, $`R^dϵ^{2d/2}\mathrm{}`$, the lifetime of the clusters is zero.
We now apply these ideas to a mean-field model of a supercooled fluid. We consider the step potential $`u(\gamma r)=0`$ for $`\gamma r>1`$ and $`u(\gamma r)=1`$ for $`\gamma r1`$. For $`\gamma 0`$, $`S(k)=1/[1+\rho \beta u(k)]`$ in the fluid phase , where $`\rho `$ is the density and $`u(k)`$ is the Fourier transform of $`\gamma ^du(x)`$ with $`x`$ and $`k^1`$ scaled by $`\gamma `$. This mean-field form of $`S(k)`$ indicates that the system is unstable for $`\rho \beta `$ such that $`1+\rho \beta _su(k_0)=0`$, where $`u(k_0)<0`$ is the minimum of $`u(k)`$. As $`TT_s=(k_B\beta _s)^1`$, $`S(k_0)ϵ^1`$, where $`ϵ=|TT_s|/T_s`$. The divergence of $`S(k_0)`$ is unchanged if a short-ranged reference potential is added.
The scaling laws for the spinodal can be rewritten with $`ϵ`$ as the scaling variable with $`\nu =1/2`$, $`2\alpha =3`$, $`\beta =1`$, and $`\gamma =1`$ rather than the values quoted previously for the $`\mathrm{\Delta }h`$ scaling field. The number of clusters in a volume $`\xi ^d`$ scales as $`R^dϵ^{3d/2}`$, the density of clusters as $`ϵ/R^dϵ^{3d/2}`$, and the density of critical phenomena fluctuations as $`ϵ/(R^dϵ^{3d/2})^{1/2}`$. The Ginsburg criterion is $`\mathrm{\Lambda }=R^dϵ^{3d/2}>>1`$, and the time scale for the critical fluctuations and the clusters is $`ϵ^1`$ and $`ϵ^1/R^dϵ^{3d/2}`$, respectively.
The primary difference between the Ising and fluid spinodals is that $`S(k)`$ diverges at $`k_00`$ for the fluid. The Fourier transform of $`S(k)`$ yields a correlation function $`\mathrm{\Gamma }(r)`$ that scales as $`r^1e^{r/\xi }\mathrm{sin}k_0r`$ in $`d=3`$, implying that the critical fluctuations near the spinodal have a characteristic length $`k_0^1<<\xi `$ on which there is a periodic spatial variation. Since the critical fluctuations are an incoherent superposition of $`\mathrm{\Lambda }`$ overlapping clusters, the spatial symmetry breaking reflected in $`\mathrm{\Gamma }(r)`$ and $`S(k)`$ also occurs for the clusters. We can show from an analysis of the Langevin equation that the clusters have a triangular structure in $`d=2`$ and a bcc or layered triangular structure in $`d=3`$. (Nucleation near the spinodal is a coalescence of clusters and the nucleation droplets have the symmetries discussed above.)
To understand the consequences of our interpretation of mean-field and near-mean-field fluctuations, we obtain the $`T`$-dependence of $`S(k_0)`$ for a fluid using scaling arguments similar to the ones used above. The structure function is related to an integral over the particles in a cluster times the probability that two particles belong to the same cluster within a critical fluctuation. The integral involves a phase factor $`_j\mathrm{exp}[i\stackrel{}{k}_{0,j}\stackrel{}{r}_j]`$, where $`|\stackrel{}{k}_{0,j}|=k_0`$ and the index $`j`$ labels the basis reciprocal lattice directions for the indicated symmetry. This phase factor is multiplied by the probability that a site belongs to a critical phenomena fluctuation, $`ϵ/(R^dϵ^{3d/2})^{1/2}`$, times the probability that the second particle belongs to the same cluster, $`ϵ/(R^dϵ^{3d/2})`$, times the number of clusters, $`(R^dϵ^{3d/2})^{1/2}`$. In addition, we average over the angles corresponding to the random orientations of the clusters. These considerations yield the scaling form:
$`S(k_0)`$ $``$ $`{\displaystyle 𝑑\stackrel{}{k}_0^{}𝑑\stackrel{}{r}\frac{ϵ^2}{R^dϵ^{3d/2}}e^{i\stackrel{}{k}_0\stackrel{}{r}}}`$ (7)
$``$ $`{\displaystyle \frac{ϵ^2}{R^dϵ^{3d/2}}}\xi R^{1d}ϵ^{(3d)/2}.`$ (8)
The spatial integral is over a region the size of $`\xi ^d`$ and $`d\stackrel{}{k}_0^{}`$ denotes an integral over angles.
Eq. (7) predicts that $`S(k_0)ϵ^{\stackrel{~}{\gamma }}`$ with $`\stackrel{~}{\gamma }=1`$ in $`d=1`$, $`\stackrel{~}{\gamma }=1/2`$ in $`d=2`$, and $`\stackrel{~}{\gamma }=0`$ in $`d3`$ in contrast to $`\gamma =1`$ for all $`d`$ as predicted by mean-field theory. We stress that the weakening of the divergence of $`S(k_0)`$ in $`d=2`$ and its supression for $`d>2`$ is a consequence of the quenched periodic structure of the clusters. If this structure is modified by finite size effects or defects, then the suppression might not be as strong.
The difference between the two calculations for the scaling behavior of $`S(k_0)`$ is the limiting procedure. If the limit $`\mathrm{\Lambda }\mathrm{}`$ ia taken before the calculation of $`S(k_0)`$, the lifetime of the clusters is zero in comparison to a measurement time. In this limit the clusters appear rotationally symmetric and a calculation similar to Eq. (7) would yield $`S(k_0)ϵ^1`$ for all $`d`$, the same result as in Ref. . However, if we assume $`\mathrm{\Lambda }`$ to be arbitrarily large but finite, the lifetime of the clusters is nonzero and the measurement time is smaller than the cluster lifetime. This assumption is consistent with the way measurements are done in experiments and simulations.
It is difficult to estimate critical exponents by fitting data directly to the suggested asymptotic form. However, because the spinodal is well defined only in the mean-field limit and simulations can be done only for finite $`R`$, we must estimate $`\stackrel{~}{\gamma }`$ by approaching the pseudospinodal. But we cannot approach it too closely because we will reach the Becker-Döring limit where the system nucleates very quickly. Because the Metropolis algorithm becomes very inefficient in near-mean-field because almost all single particle moves are accepted, we included a small hard core of diameter $`\sigma `$ to decrease the acceptance probability.
Our Monte Carlo results (see Fig. 1) in $`d=1`$ were fit to $`AT^{\stackrel{~}{\gamma }}ϵ^{\stackrel{~}{\gamma }}`$ with $`\stackrel{~}{\gamma }=1.0\pm 0.05`$ and $`T_s0.83`$. The $`T^{\stackrel{~}{\gamma }}`$ term accounts in part for corrections to scaling which typically depend on $`d`$. The errors were estimated from different possible fits. The $`d=2`$ results were fit to $`Aϵ^{\stackrel{~}{\gamma }}`$ with $`\stackrel{~}{\gamma }=0.4\pm 0.05`$ and $`T_s0.81`$. The quality of the fits is not as good as in $`d=1`$, but the estimate for $`\stackrel{~}{\gamma }`$ is consistent with our prediction of $`\stackrel{~}{\gamma }=0.5`$. The $`d=3`$ results are consistent with $`Aϵ^{\stackrel{~}{\gamma }}`$ with $`\stackrel{~}{\gamma }=0.16\pm 0.02`$ and $`T_s0.71`$ and with $`A\mathrm{log}ϵ+B`$ with $`T_s0.86`$. Details of the simulations will be discussed in a longer paper.
In summary, our theoretical and simulation results imply that the structure of the fluctuations in mean-field and near-mean-field systems differs qualitatively from that of non-mean-field systems. This structure leads to a suppression of the divergence of the measured static structure function near a pseudospinodal relative to the mean-field prediction in systems for which a spatial symmetry is broken. The dependence on $`d`$ of this suppression is such that there is no apparent divergence for $`d3`$ subject to logarthmic corrections. In these systems there is a growing correlation length as the pseudospinodal is approached that cannot be obtained from a measurement of $`S(k)`$. We note that a divergent dynamical length has been found above the glass transition in the spherical $`p`$ spin model . This work, coupled with our results, raises the question of how the spinodal influences the dynamical length. Our predictions have important implications for the understanding of processes such as nucleation and glass formation in supercooled fluids. In particular, we expect that the existence of clusters should help us understand the universal scaling behavior found for the dielectric response of organic glass formers.
We thank R. C. Brower, Bulbul Chakraborty, and John Rundle for useful discussions. This work was supported by NSF grant DMR-9633385 (WK and HG) and DOE DE-FG02-95ER 14498 (WK and MA). Acknowledgment is made to the donors of the Petroleum Research Foundation, administrated by the American Chemical Society, for partial support of the research at Kalamazoo College. Work on LA-UR 00-117 at Los Alamos National Laboratory was supported by the U.S. DOE LDRD-DR 98605. |
warning/0001/hep-th0001084.html | ar5iv | text | # Untitled Document
IFT-P.006/2000
The Tachyon Potential in Open Neveu-Schwarz String Field Theory
Nathan Berkovits<sup>1</sup> e-mail: nberkovi@ift.unesp.br
Instituto de Física Teórica, Universidade Estadual Paulista
Rua Pamplona 145, 01405-900, São Paulo, SP, Brasil
A classical action for open superstring field theory has been proposed which does not suffer from contact term problems. After generalizing this action to include the non-GSO projected states of the Neveu-Schwarz string, the pure tachyon contribution to the tachyon potential is explicitly computed. The potential has a minimum of $`V=\frac{1}{32g^2}`$ which is $`60\%`$ of the predicted exact minimum of $`V=\frac{1}{2\pi ^2g^2}`$ from D-brane arguments.
January 2000
1. Introduction
The Neveu-Schwarz (NS) sector of the non-GSO projected superstring has recently been reconsidered as part of a sensible physical theory . Although this sector contains a tachyon, there have been proposals for removing the undesired properties of the the tachyon by assuming a tachyon potential which is bounded from below.
The most efficient method for computing the tachyon potential uses open string field theory , however the cubic action of for open superstring field theory contains contact term problems which spoil gauge invariance. Recently, a new action for open superstring field theory has been constructed which does not suffer from contact term problems. This action resembles a Wess-Zumino-Witten action and can be naturally obtained by embedding the N=1 description of the superstring into an N=2 string .
In this paper, the pure tachyon contribution to the tachyon potential will be explicitly computed using this new action. The pure tachyon contribution is
$$V(T)=\frac{1}{4g^2}T^2+\frac{1}{2g^2}T^4,$$
which has a minimum of $`V(T_0)=\frac{1}{32g^2}`$ when $`T_0=\pm \frac{1}{2}.`$ This value of the minimum is $`60\%`$ of the predicted exact minimum of $`V(T_0)=\frac{1}{2\pi ^2g^2}`$ using D-brane arguments<sup>2</sup> In the original version of this paper, the mass of the brane-antibrane was incorrectly stated to be $`\frac{1}{\pi ^2g^2}`$. This value of the mass is only correct if one doubles the number of states in the string field theory action to allow for strings which end on the brane or antibrane . where the mass of the brane-antibrane is $`\frac{1}{2\pi ^2g^2}`$ . It would be interesting to check if the remaining $`40\%`$ comes from including contributions to the effective tachyon potential from non-tachyon fields, as was found for the bosonic string tachyon potential in .
2. Neveu-Schwarz String Field Theory Action
Using the superstring field theory action of , the GSO-projected NS contribution is given by
$$S=\frac{1}{2g^2}Tr(e^\mathrm{\Phi }Qe^\mathrm{\Phi })(e^\mathrm{\Phi }\eta _0e^\mathrm{\Phi })_0^1𝑑t(e^{t\mathrm{\Phi }}_te^{t\mathrm{\Phi }})\{e^{t\mathrm{\Phi }}Qe^{t\mathrm{\Phi }},e^{t\mathrm{\Phi }}\eta _0e^{t\mathrm{\Phi }}\}$$
where $`\eta _0=𝑑z\eta (z)`$ is defined by fermionizing the super-reparameterization ghosts as $`\gamma =\eta e^\varphi `$ and $`\beta =\xi e^\varphi `$,
$$Q=𝑑z[c(T_{matter}\eta \xi \frac{1}{2}\varphi \varphi ^2\varphi bc)+\eta e^\varphi G_{matter}\eta \eta e^{2\varphi }b],$$
and $``$ signifies the two-dimensional correlation function in the “large” RNS Hilbert space where $`\xi cc^2ce^{2\varphi }=2`$. The normalization of (2.1) has been fixed by requiring that the quadratic Yang-Mills contribution to the action is $`S=\frac{1}{4g^2}Trd^{10}xF_{mn}F^{mn}`$, which is the correct sign for the $`(+\mathrm{}+)`$ metric that is being used. String fields are multiplied using the midpoint interaction of and $`\mathrm{\Phi }`$ is related to the NS string field $`V`$ of by $`\mathrm{\Phi }=\xi _0V`$ or $`V=\eta _0\mathrm{\Phi }`$.
In the GSO-projected sector, the NS string field $`\mathrm{\Phi }`$ is bosonic. Since the unprojected NS states are fermionic with respect to the projected NS states, it will be convenient to define $`\widehat{\mathrm{\Phi }}=\mathrm{\Phi }_B\times I+\mathrm{\Phi }_F\times \sigma _1`$ where $`\mathrm{\Phi }_B`$ described the projected states, $`\mathrm{\Phi }_F`$ describes the unprojected states, $`I`$ is the $`2\times 2`$ identity matrix, and $`(\sigma _1,\sigma _2,\sigma _3)`$ are the Pauli matrices . Furthermore, it will be convenient to define
$$\widehat{Q}Q\times \sigma _3,\widehat{\eta }_0\eta _0\times \sigma _3,$$
which satisfy $`\widehat{Q}(\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2)=(\widehat{Q}\widehat{\mathrm{\Phi }}_1)\widehat{\mathrm{\Phi }}_2+\widehat{\mathrm{\Phi }}_1(\widehat{Q}\widehat{\mathrm{\Phi }}_2)`$ and $`\widehat{\eta }_0(\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2)=(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_1)\widehat{\mathrm{\Phi }}_2+\widehat{\mathrm{\Phi }}_1(\widehat{\eta }_0\widehat{\mathrm{\Phi }}_2)`$.
The complete non-GSO projected NS string field theory action is defined by
$$S=\frac{1}{4g^2}Tr(e^{\widehat{\mathrm{\Phi }}}\widehat{Q}e^{\widehat{\mathrm{\Phi }}})(e^{\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{\widehat{\mathrm{\Phi }}})_0^1𝑑t(e^{t\widehat{\mathrm{\Phi }}}_te^{t\widehat{\mathrm{\Phi }}})\{e^{t\widehat{\mathrm{\Phi }}}\widehat{Q}e^{t\widehat{\mathrm{\Phi }}},e^{t\widehat{\mathrm{\Phi }}}\widehat{\eta }_0e^{t\widehat{\mathrm{\Phi }}}\}$$
where the trace is over the $`2\times 2`$ matrices as well as the Chan-Paton matrices.
One can check that (2.1) is invariant under the WZW-like gauge transformation
$$\delta e^{\widehat{\mathrm{\Phi }}}=(\widehat{Q}\widehat{\mathrm{\Omega }})e^{\widehat{\mathrm{\Phi }}}+e^{\widehat{\mathrm{\Phi }}}(\widehat{\eta }_0\widehat{\mathrm{\Omega }}^{})$$
where $`\widehat{\mathrm{\Omega }}`$ and $`\widehat{\mathrm{\Omega }}^{}`$ are string fields of the form $`\widehat{\mathrm{\Omega }}=\mathrm{\Omega }_F\times \sigma _3+i\mathrm{\Omega }_B\times \sigma _2`$ with $`\mathrm{\Omega }_F`$ being fermionic and projected while $`\mathrm{\Omega }_B`$ is bosonic and unprojected. One subtle point in proving this gauge invariance is that $`\widehat{\mathrm{\Phi }}_1\widehat{\mathrm{\Phi }}_2=\widehat{\mathrm{\Phi }}_2\widehat{\mathrm{\Phi }}_1`$ since when $`A`$ and $`B`$ are unprojected states, $`AB=BA`$ where the minus sign is if they are bosons and the plus sign is if they are fermions. This reversal of the usual statistics comes from square-root factors produced by the $`\frac{1}{2}`$-integer conformal weight of unprojected NS states. Note that a similar subtlety occurs with unprojected states using the action of .
3. Computation of Tachyon Potential
Expanding the action of (2.1) in powers of $`\widehat{\mathrm{\Phi }}`$, one obtains
$$S=\frac{1}{2g^2}Tr\frac{1}{2}(\widehat{Q}\widehat{\mathrm{\Phi }})(\widehat{\eta }_0\widehat{\mathrm{\Phi }})\frac{1}{6}\widehat{\mathrm{\Phi }}\{\widehat{Q}\widehat{\mathrm{\Phi }},\widehat{\eta }_0\widehat{\mathrm{\Phi }}\}\frac{1}{24}[\widehat{\mathrm{\Phi }},\widehat{Q}\widehat{\mathrm{\Phi }}][\widehat{\mathrm{\Phi }},\widehat{\eta }_0\widehat{\mathrm{\Phi }}]+\mathrm{}.$$
To compute the term with $`N`$ $`\widehat{\mathrm{\Phi }}`$’s, one uses the map
$$w(z)=\left(\frac{1iz}{1+iz}\right)^{\frac{2}{N}}$$
from the disc to a $`2\pi /N`$ wedge of the complex plane. Rotating this map by a factor $`e^{\frac{2\pi i}{N}}`$ allows each successive string field to get mapped to a different $`2\pi /N`$ wedge. The center of the $`J^{th}`$ disc gets mapped to the point $`e^{\frac{2\pi i(J1)}{N}}`$ and, to obtain an SL(2,R)-invariant expression, the $`J^{th}`$ string field gets multiplied by a factor $`(e^{\frac{2\pi i(J1)}{N}}\frac{4}{Ni})^h`$ where $`h`$ is the conformal weight of the string field and $`\frac{4}{Ni}`$ is $`\frac{dw}{dz}|_{z=0}`$ .
The tachyon field $`T(x)`$ appears in the string field $`\widehat{\mathrm{\Phi }}`$ as
$$\widehat{\mathrm{\Phi }}=i\xi ce^\varphi T(x)\times \sigma _1$$
where the factor of $`i`$ is needed to get the right sign for the kinetic term. One can easily compute that at zero momentum,
$$\widehat{Q}\widehat{\mathrm{\Phi }}=T(\frac{1}{2}cc\xi e^\varphi +\eta e^\varphi )\times \sigma _2,\widehat{\eta }_0\widehat{\mathrm{\Phi }}=Tce^\varphi \times \sigma _2.$$
Since $`\xi cc^2ce^{2\varphi }=2`$, the only pure tachyon contribution to the action of (2.1) comes from the quadratic and quartic terms of (3.1). The quadratic contribution to the action (which is minus the tachyon potential) is given by
$$V_2=\frac{1}{2g^2}T^2(\frac{1}{2}cc\xi e^\varphi (1))(ce^\varphi (1))(2i)^{\frac{1}{2}}(2i)^{\frac{1}{2}}=\frac{1}{4g^2}T^2.$$
The quartic contribution to the action is given by
$$V_4=\frac{T^4}{24g^2}(i)^{\frac{1}{2}}(1)^{\frac{1}{2}}(i)^{\frac{1}{2}}(1)^{\frac{1}{2}}(\xi ce^\varphi (1))(\eta e^\varphi (i))(\xi ce^\varphi (1))(ce^\varphi (i))$$
$$+(\eta e^\varphi (1))(\xi ce^\varphi (i))(\xi ce^\varphi (1))(ce^\varphi (i))+(\xi ce^\varphi (1))(\eta e^\varphi (i))(ce^\varphi (1))(\xi ce^\varphi (i))$$
$$+(\eta e^\varphi (1))(\xi ce^\varphi (i))(ce^\varphi (1))(\xi ce^\varphi (i))=\frac{T^4}{2g^2}.$$
So $`V(T)=V_2+V_4=\frac{1}{4g^2}T^2+\frac{1}{2g^2}T^4`$ which has a minimum of $`V(T_0)=\frac{1}{32g^2}`$ when $`T_0=\pm \frac{1}{2}`$.
Acknowledgements: I would like to thank Oren Bergman, Ashoke Sen, Ion Vancea and Barton Zwiebach for useful discussions, Caltech for their hospitality, and CNPq grant 300256/94-9 for partial financial support.
References
relax A. Sen, Stable Non-BPS States in String Theory JHEP 9806 (1998) 007, hep-th/9803194; A. Sen, Tachyon Condensation on the Brane Antibrane System JHEP 9808 (1998) 012, hep-th/9805170. relax O. Bergman and M.R. Gaberdiel, A Non-Supersymmetric Open String Theory and S-Duality, Nucl. Phys. B499 (1997) 183, hep-th/9701137; T. Yoneya, Spontaneously Broken Space-Time Supersymmetry in Open String Theory without GSO Projection, hep-th/9912255. relax V.A. Kostelecky and S. Samuel, The Static Tachyon Potential in the Open Bosonic String Theory, Phys. Lett. B207 (1988) 169. relax A. Sen, Universality of the Tachyon Potential, hep-th/9911116. relax A. Sen and B. Zwiebach, Tachyon Condensation in String Field Theory, hep-th/9912249. relax E. Witten, Interacting Field Theory of Open Superstrings, Nucl. Phys. B276 (1986) 291. relax C. Wendt, Scattering Amplitudes and Contact Interactions in Witten’s Superstring Field Theory, Nucl. Phys. B314 (1989) 209; J. Greensite and F.R. Klinkhamer, Superstring Amplitudes and Contact Interactions, Nucl. Phys. B304 (1988) 108. relax N. Berkovits, Super-Poincare Invariant Superstring Field Theory, Nucl. Phys. B450 (1995) 90, hep-th 9503099; N. Berkovits, A New Approach to Superstring Field Theory, proceedings to the $`32^{nd}`$ International Symposium Ahrenshoop on the Theory of Elementary Particles, Fortschritte der Physik (Progress of Physics) 48 (2000) 31, hep-th/9912121; N. Berkovits and C.T. Echevarria, Four-Point Amplitude from Open Superstring Field Theory, hep-th/9912120. relax N. Berkovits and C. Vafa, N=4 Topological Strings, Nucl. Phys. B433 (1995) 123, hep-th/9407190. relax A. Sen, private communication. relax D. Friedan, E. Martinec, and S. Shenker, Conformal Invariance, Supersymmetry, and String Theory, Nucl. Phys. B271 (1986) 93. relax E. Witten, Noncommutative Geometry and String Field Theory, Nucl. Phys. B268 (1986) 253. relax E. Cremmer, A. Schwimmer and C. Thorn, The Vertex Function in Witten’s Formulation of String Field Theory, Phys. Lett. B179 (1986) 57; D.J. Gross and A. Jevicki, Operator Formulation of Interacting String Field Theory, Nucl. Phys. B283 (1987) 1; A. LeClair, M.E. Peskin and C.R. Preitschopf, String Field Theory on the Conformal Plane (I), Nucl. Phys. B317 (1989) 411. |
warning/0001/astro-ph0001163.html | ar5iv | text | # Correlated X-ray Spectral and Timing Behavior of the Black Hole Candidate XTE J1550–564: A New Interpretation of Black Hole States
## 1 Introduction
The X-ray transient XTE J1550–564 was discovered on 1998 September 7 (Smith, 1998) with the All Sky Monitor (ASM) on board the Rossi X-ray Timing Explorer (RXTE). Soon after optical (Orosz, Bailyn and Jain, 1998) and radio (Campbell-Wilson et al., 1998) counterparts were identified. Observations with the RXTE Proportional Counter Array (PCA) were performed on an almost daily basis between 1998 September 7 and 1999 May 20.
The 1998/1999 outburst of XTE J1550–564 consisted of two parts, separated by a minimum that occurred around 1998 December 3 (MJD 51150; see Figure 1). The first part of the outburst (until MJD 51150) has been the subject of work by Cui, et al. (1999, timing behavior during the rise), Sobczak et al. (1999, spectral behavior), and Remillard et al. (1999a, timing behavior). During the initial 10 day rise a quasi-periodic oscillation (QPO) was found, together with a second harmonic (Cui, et al., 1999). It had a frequency between 82 mHz and 4 Hz, that smoothly increased with the X-ray flux. During the strong flare (reaching 6.8 Crab) that occurred around MJD 51075 (Remillard et al., 1998), a QPO was found with a frequency of 185 Hz (Remillard et al., 1999a). High frequency oscillations were also found on three other occasions, with frequencies between 161 and 238 Hz. Low frequency QPOs (3–13 Hz) were also observed, during the strong flare and the decay of the first part of the outburst. Correlations between spectral parameters and the low frequency QPOs (for the entire outburst) have been presented by Sobczak et al. (2000a).
Traditionally the behavior of black hole X-ray transients has been described in terms of transitions between four canonical black hole states (for overviews of black hole spectra and power spectra we refer to Tanaka and Lewin (1995) and van der Klis (1995b)). The classification of an observation into one of these four states is based on luminosity, spectral and timing properties, and on the order in which they occur.
The spectra of black hole X-ray binaries are often described in terms of a disk black body component, believed to be coming from an accretion disk, and a power law tail at high energies, which is thought to arise in a Comptonizing region (e.g. Sunyaev and Titarchuk (1980)). The power spectra can be described by a (broken) power law, with sometimes one or more quasi-periodic oscillation peaks superimposed. The parameter usually thought to determine the state of the black hole is the mass accretion rate. The definitions of the different states are rather loose and have shown some variation between authors; we therefore only give a general overview of the four canonical states in order of inferred increasing mass accretion rate:
* Low State (LS): The 2–20 keV X-ray spectrum can be described by a single power law, with a photon index ($`\mathrm{\Gamma }`$) of $``$1.5 plus sometimes a weak disk black body component ($`kT<`$1 keV; less than a few percent of the total luminosity). The power spectrum shows strong band-limited noise with a typical strength of 20–50% rms and a break frequency ($`\nu _b`$) below 1 Hz.
* Intermediate State (IS): In the X-ray spectrum both the power law ($`\mathrm{\Gamma }`$2.5) and disk black body components ($`kT`$1 keV) are present. The noise in the power spectrum is weaker (typically 5–20% rms) and the break frequency is higher ($`\nu _b`$1–10 Hz) than in the LS. QPOs between 1 and 10 Hz are sometimes observed.
* High State (HS): The X-ray spectrum is dominated by the disk black body component ($`kT`$1 keV), and the power law component ($`\mathrm{\Gamma }23`$) is weak or absent. The noise in the power spectrum is power law like and very weak, with a strength of less than 2–3 % rms.
* Very High State (VHS): Like in the IS the X-ray spectrum is a combination of a disk black body ($`kT`$1–2 keV) and a power law ($`\mathrm{\Gamma }2.5`$). The power spectrum shows noise, that can either be HS-like (power law) or LS-like (band limited; 1-15% rms, $`\nu _b`$1–10 Hz). QPOs are often seen in the VHS with frequencies between 1 and 10 Hz.
Note that there is little difference between the power spectral and X-ray spectral properties of the VHS and IS, although conventionally the total flux in the VHS is described as much higher than that in the IS (Belloni et al., 1996; Méndez and van der Klis, 1997; Belloni et al., 1997). The reason that the IS was introduced as a separate state (Belloni et al., 1997) basically is that it occurred in GS 1124–68 at epochs when the source appeared to be evolving gradually from HS to LS and had a flux that was only 10% of that during the peak of the VHS (Ebisawa, et al., 1994). It could not, therefore, in the one-dimensional classification outlined above, be identified with the VHS which by definition occurs at the upper end of the inferred mass accretion range, above the HS. All three states, LS, IS, and VHS are characterized by the presence of strong band limited noise and a hard power law component, and are in that sense much more similar to each other than to the HS, which is characterized by these features being very weak or absent.
According to Sobczak et al. (1999) XTE J1550–564 went through the VHS, HS, and IS during the first part of the outburst. In this paper we present a study of the correlated spectral and timing behavior of XTE J1550–564 during the second part of its outburst. We will discuss the results for XTE J1550–564 using some of the canonical terminology, in order to compare our results with those of other transients. However, we will also discuss the discrepancies of the canonical one-dimensional model with the results obtained for XTE J1550–564; these discrepancies concern in particular the way in which the various states relate to each other. We find the source moved through all the four black hole states, in a way that is highly suggestive of a new two dimensional interpretation of the black hole states.
In Section 2 we explain our analysis methods. Our results are presented in Sections 3, 4, and 5. These sections are intended for the specialized reader, the level of detail is rather high. A separate summary of the most important results is therefore given at the beginning of the discussion (Section 6). We end by summarizing our main conclusions in Section 7.
## 2 Observations and analysis
For our analysis we used all Public Target Of Opportunity RXTE/PCA (Bradt, Rothschild and Swank, 1993; Jahoda et al., 1996) data for XTE J1550–564 taken between 1998 November 22 23:38 UTC (MJD 51139)and 1999 May 20 19:37 UTC (MJD 51318). This adds up to 171 single observations, corresponding to a total observing time of $``$400 ks. When we refer to a single observation, we mean a part of the data with its own unique RXTE observation ID; all observations will be referred to by their Modified Julian Date (MJD) at the start of the observation. In all light curves and color-color diagrams each point represents one single observation.
The PCA data were obtained in several different modes (see Table 1), some of which were active simultaneously. On 1999 March 22 (MJD 51259), the high voltages of the PCA instrument were changed, resulting in a different energy gain. The count rates and the colors obtained after this change (gain Epoch 4) can not be directly compared with the data obtained before the change (gain Epoch 3). Data obtained during satellite slews, Earth occultations, and South Atlantic Anomaly passages were removed from our data set.
The Standard 2 data (see Table 1) were used to create light curves, color curves, color-color diagrams (CDs), and a hardness-intensity diagram (HID). Only data of proportional counter units (PCUs) 0 and 2 were used for this, since these were the only two that were active during all the observations. All PCA count rates and colors given in this paper are only for those two PCUs combined. The data were background subtracted, but dead time corrections ($`<`$6%) were not applied. For the light curves and colors, the photon energy channel boundaries were chosen in such a way that the corresponding energies for Epoch 3 and 4 matched as well as possible. A color is the ratio of the count rates in two energy bands. We define the soft color (SC) as the ratio of the count rates in the 6.2–15.7 keV and 2–6.2 keV bands (Epoch 3), or 6.1–15.8 keV and 2–6.1 keV bands (Epoch 4); hard color (HC) is defined as the ratio of the count rates in the 15.7–19.4 keV and 2–6.2 keV bands (Epoch 3), or 15.8–19.4 keV and 2–6.1 keV bands (Epoch 4). This definition of colors, with the same band in the denominator for both the hard and soft color, has the advantage that comparison with a two-component model is straightforward. Namely, if the source spectrum is dominated by the contribution from two spectral components (in our case a disk black body and a power law), then a color data point will lie on the line connecting the color points of the individual spectral components (Wade, 1982; van Teeseling and Verbunt, 1994) and the ratio of distances from the data point to the points representing the components is the inverse ratio of their contribution in the 2–6.2 keV (Epoch 3) or 2–6.1 keV (Epoch 4) bands. For the HID we used the 2–60 keV count rate (representing the full energy range covered by the PCA) as intensity, and the hard color (see above) as hardness. The observation starting at MJD 51298.18 was included in the light curves, but not in the CD, HID, and color curves, since at high energies the source could not be detected above the background.
The Standard 2 data were also used to perform a number of spectral fits. The spectra were background subtracted, and fitted in the 2.5–25.0 keV (Epoch 3) or 3.1–25.0 keV (Epoch 4) band, using a systematic error of 2%. Fits were made with XSPEC 10.00, using a fit function that consisted of a disk black body, a power law, a Gaussian line with a fixed energy of 6.5 keV and a width of 1–1.5 keV, and an edge around 7 keV. Interstellar absorption was modelled using the Wisconsin cross sections (Morrison and McCammon, 1983), with $`N_H`$ fixed to a value of 2$`\times 10^{22}`$ atoms/cm<sup>2</sup> (Sobczak et al., 1999). We found that the results of the spectral fits were very sensitive to the version of the PCA response matrix we used. The response matrices were initially created using FTOOLS version 4.2 and later with the updated version 5.0. Our initial fits showed that the inner disk radius and color temperature of the disk black body component were, respectively, correlated and anti-correlated with the hardness of the spectrum. However, both correlations were found to be reversed when using the updated version of the response matrices. In view of this we decided to omit the spectral fits from the current paper, and only discuss the spectral behavior using the color-color diagrams (which are matrix-independent). For a complete spectral analysis of XTE J1550–564 we refer to Sobczak et al. (1999, 2000c).
The three high time resolution modes (with time resolution $`2^{13}`$ s, see Table 1) were used to produce 1/16–512 Hz power spectra in their respective energy bands and in the combined 2–60 keV band; the same was done for the $`2^{20}`$ s mode (MJD 51297–51318), which was split into three energy bands with similar energy ranges. In order to study the variability at lower frequencies, 1/128–128 Hz power spectra were created in the 2–60 keV band, and for some observations also in 8 energy bands covering 2–60 keV. No background or dead time corrections were applied to the data before the power spectra were created; the effect of dead time on the Poisson noise was accounted for in the power spectral fits. The power spectra were selected on time, count rate, color, or a combination of these, before they were averaged and fitted. Although most of the power spectra presented in this paper are normalized according to Leahy et al. (1983), the actual power spectral fits were made to power spectra that were rms normalized (van der Klis, 1995a).
The power spectra were fitted with a combination of several functions. A constant was used to represent the Poisson level. The noise at low frequencies was fitted with a power law ($`P\nu ^\alpha `$), or with a broken power law ($`P\nu ^{\alpha _1}`$ for $`\nu <\nu _b`$; $`P\nu ^{\alpha _2}`$ for $`\nu >\nu _b`$). In practice, the low frequency noise component in the power spectra of most observations could be fitted with a single power law. However, when combining several observations, due to the smaller uncertainties, it became apparent that a single power law did not yield acceptable fits, especially around 1 Hz; using a broken power law for those combined observations resulted in much better fits. For the single observations we continued using a single power law, since for a broken power law the break frequency was poorly constrained, and $`\chi ^2`$ did not differ significantly between the two fit functions. Most QPOs were fitted with a Lorentzian, $`P[(\nu \nu _c)^2+(FWHM/2)^2)]^1`$, where $`\nu _c`$ is the central frequency and FWHM the full-width-at-half-maximum. In some cases narrow QPOs were found for which a Lorentzian provided inadequate fits; in those cases a Gaussian was used, $`Pe^{(\nu \nu _c)^2/\sigma ^2}`$, where $`\nu _c`$ is the central frequency and $`\sigma `$ the width of the Gaussian. Furthermore, we sometimes used an exponentially cutoff power law ($`P\nu ^\alpha e^{\nu /\nu _{cutoff}}`$) to fit an extra noise component at low frequencies. The errors on the fit parameters were determined using $`\mathrm{\Delta }\chi ^2=1`$. The energy dependence of the power spectral features was in general obtained by fixing all parameters, except the amplitude, to their values obtained in a specific band. However, in some cases, when the shape of the QPOs was found to change between energy bands, the FWHM and/or frequency were not fixed. Upper limits on the strength of the power spectral features were determined by fixing all their parameters, except the amplitude, to values obtained in another energy band or observation, and using $`\mathrm{\Delta }\chi ^2=2.71`$ (95% confidence).
Unless otherwise stated, all the power spectral parameters are those in the 2–60 keV band, and the noise rms amplitude is that in the 0.01–1 Hz range.
## 3 Light curves, color–color diagrams
Figure 1 shows the one day averaged ASM 2–12 keV light curve of the 1998/1999 outburst of XTE J1550–564, with the dashed line marking the start of our PCA data set. It shows a broad local minimum around MJD 51150, which naturally divides the outburst into two parts. Following this minimum ($``$8.5 ASM counts s<sup>-1</sup>) the count rate increased by a factor of 10 within 20–25 days, and then rose to about 200 counts s<sup>-1</sup> in 40 days. After a relatively flat top XTE J1550–564 showed an initially slow decline ($``$55 days) to about 100 ASM counts $`s^1`$, which was followed by a decrease by a factor of 100 in $``$45 days. On one day during during the first part of the decline, MJD 51250, the ASM lightcurve showed a strong dip; the count rate was found to be $``$40% lower than in the two adjacent observations (see also Section 4.3.3). At the end of the outburst, around MJD 51310, two small flares occurred, reaching a few ASM counts $`s^1`$. Although XTE J1550–564 never again reached the level of the MJD 51074–51076 flare (6.8 Crab or $``$490 ASM counts s<sup>-1</sup>), it was bright (in the 2–12 keV band) for a longer period of time during the second part than during the first part of the outburst; during the first part it was observed above 150 counts s<sup>-1</sup> on only six days, during the second part it was above this level on more than 70 days.
The PCA data set used in our analysis started on MJD 51138 (dashed line in Figure 1), just before the minimum between the two parts of the outburst was reached. The PCA light curves in different energy bands are shown in Figures 2a-f. The time of the PCA gain change is indicated by the dotted line. As expected, the overall 2–60 keV light curve, dominated by the contribution of the low energy bands, has a shape that is similar to that of the ASM light curve. Note that as the photon energy increases, the local minimum near MJD 51150 occurs later and seems to become broader. While the light curves in the low energy bands have more or less the same profile as in the ASM light curve, in the high energy bands they look strikingly different. The light curve in Figure 2e shows several strong flares on top of the overall outburst profile, and above 17.5 keV (Fig. 2f) the light curve is dominated by these flares. For later use the relatively small flares were numbered 1 to 5 (Fig. 2e), and the bigger (and broader) ones, that clearly showed up as branches in the CD (see below) were numbered I to V.
Figures 2g and 2h show the evolution of soft and hard color with time. The change from Epoch 3 to 4 is again indicated by a dotted line. Until the end of flare/branch II the colors were quite well correlated with the count rate in all energy bands, the only clear exception being the drop in hard color during the rise (MJD 51150–51160). After that, while the count rates dropped, the colors increased (indicating a hardening of the spectrum) and were only correlated with the count rates during flares/branches III-V (see Fig. 2e).
Combining the two color curves, two color-color diagrams (CDs) were produced. The CDs for Epoch 3 and 4 are shown in Figure 3a and 3b, respectively. In both CDs we also plotted the expected colors for a disk black body (DBB) spectrum at different temperatures (squares), and the expected colors for a power law spectrum with different indices (triangles), both for an assumed $`N_H`$ of $`2\times 10^{22}`$ atoms cm<sup>-2</sup> (Sobczak et al., 1999). Note that the values of the expected colors (unlike the observed colors) do depend on the version of the PCA response matrix that is used; the lines in Figure 3 were produced using FTOOLS version 4.2. As explained in Section 2, if the energy spectrum were a combination of only a DBB and a power law, then the corresponding point in the CD would lie on the straight line that connects the appropriate point on the DBB curve with the one on the power law curve. Although fits show that more spectral components are needed to obtain a good $`\chi ^2`$ (see e.g. Sobczak et al. (1999)), we do gain some insight into how the parameters (temperature of the DBB, index of the power law) and relative strength of the two most important spectral components behave.
The pattern traced out in the CD before the gain change (Figure 3a) shows roughly three branches. A spectrally soft branch lies very close and nearly parallel to the DBB curve between 0.8 and 1.05 keV, and two spectrally hard branches (I and II, corresponding to big flares I and II in Figure 2) that lie more or less parallel to the power law curve. It should be noted that the values for the temperature are read from the CD; values obtained from spectral fits yield a maximum temperature of $``$1.1 keV (instead of 1.05 keV), see Sobczak et al. (2000c). In the following, we give a short description of how the source moved through the CD. The spectrum of the first observation, on MJD 51139 (marked BEGIN in Figure 3a), can be described as a combination of a DBB and a power law spectrum; it neither lies close to the DBB curve nor to the power law curve. As time progressed the source moved to the left in the CD along branch I, i.e. towards a pure DBB spectrum. Around MJD 51157 the source was located close to the DBB curve, with a temperature of $``$0.8 keV. Subsequently the temperature increased to $``$1.0 keV, while the source stayed close to the DBB curve. Around MJD 51179, at SC$``$0.13 in the CD, the source suddenly left the DBB curve in the direction of the power law curve (flare 1), indicating a relative increase in the strength of the power law component. On MJD 51182 the source had returned close to the DBB curve, with a somewhat higher temperature than before. After that, the temperature increased to a maximum value of $``$1.05 keV (on MJD 51204), and then decreased to $``$1.0 keV (on MJD 51231). During this period four more flares (4–5) occurred, around MJDs 51200, 51208, 51215, and 51233 (see also Figure 2e). Similar to flare 1, these four flares also pointed away from the DBB curve. After each flare, except flare 5, the source returned to the DBB curve at a similar temperature as before. During the decay of flare 5 a new (big) flare started, which developed into branch/flare II. This branch lay relatively close to the power law curve, indicating that the power law became the dominant spectral component. The spectrally hardest observation on branch II was the one on MJD 51250. After MJD 51250, the source moved back into the direction of the DBB curve again, but at a lower temperature. The time it took the source to move down branch II was shorter than for it to move up that branch, $``$4.5 days and $``$9 days, respectively. The transition down from HC=0.02 to HC=0.008 even occurred in less than one day. After the gain change, on MJD 51260 (see inset in Figure 3b, indicated by BEGIN) the source was found relatively close to the DBB curve. A small branch (III) was traced out around MJD 51272, almost parallel to the power law curve. From MJD 51283 to 51298 another branch (IV) was traced out which also lay parallel to the power law curve. Finally, during the last observations (MJDs 51299–51318) a branch (V) was traced out, that pointed towards the regions of the power law curve with indices lower than 2.5.
The HID (2–60 keV count rate vs. hard color) is shown in Figure 4. It clearly shows that the five flares (1–5) occurred at the highest count rates, and that the five branches (I–V) are well separated from each other in count rate. On each hard branch the count rate never varied more than by a factor 2–3. Note that at the lowest count rates the observations between hard branches tend to have harder spectra than those at higher count rates.
## 4 Power spectra
The power spectra will be presented in order of time, but when it seems more appropriate also according to their position in the CD. We start by giving a short overview of the broad band power spectral behavior in the next paragraph. Then, a more detailed study follows of the power spectra during the start of the second part of the outburst, the broad maximum and flares, branch II, and the decay.
In Figure 5b we show the total power in the 1/128–128 Hz power spectra as a function of time. One can immediately see that increases in source variability occur whenever the source is on one of the five branches. Between MJDs 51150 and 51240, when the X-ray spectrum was soft, the total power had a strength of 1% to 3% rms, which is typical for the high state. Similar weak noise was also found between branches II, III, and IV. On branches I, II, and III the power spectra had strengths between 5% and 15% rms, suggesting that the source was is the intermediate and/or very high state. On the last two branches the power increased from around 5% to almost 60% rms. Noise rms amplitudes of several tens of percent are usually only found in the low state.
### 4.1 Start of the second part of the outburst
During our first PCA observations the source was still in the decay of the first part of the outburst. The power spectra of the first three observations (MJDs 51139–51143) were very similar. Their combined 2–60 keV power spectrum (Figure 6a) showed band limited noise at low frequencies, a peaked feature around 2.5 Hz and a QPO around 9 Hz, which were fitted with, respectively, a broken power law, and two Lorentzians. The strength of the noise was 2.49$`\pm `$0.07% rms, with $`\nu _b=4.2_{0.6}^{+0.3}`$ Hz, $`\alpha _1=0.1_{0.6}^{+0.3}`$, and $`\alpha _2=1.2\pm 0.1`$. The peaked component or QPO close to the break had a frequency of 2.60$`\pm `$0.06 Hz, a FWHM of 1.0$`\pm `$0.5 Hz, and an rms amplitude of 2.4$`\pm `$0.6%. The QPO at 9.17$`\pm `$0.14 Hz had a FWHM of 3.6$`\pm `$0.8 Hz, and an rms amplitude of 3.8$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$%. The power spectrum depended strongly on energy, as can be seen from Figure 6. At low energies (Fig. 6b) it was dominated by the noise component with a peaked feature around the break, whereas at higher energies (Fig. 6c) both this peak and the noise component were replaced by a broad peak around 9 Hz. The moderate noise strength and presence of QPOs classify these observations as IS/VHS. Since the count rates in these observations are considerably lower than during earlier (and later) observations where QPOs and moderate noise strengths were found, the source was probably in the IS rather than in the VHS. These observations were also classified as IS by Sobczak et al. (1999). It should be noted that by combining several power spectra, narrow features in the individual power spectra may be smoothed out and form broad bumps like the one seen in Figure 6c. For instance, Figure 7 shows the 6.5–13.1 keV power spectrum of the first of these three observations, on MJD 51139. It could be fitted with QPOs at 3$`\pm `$0.2, 6.5$`\pm `$0.3, 9.63$`\pm `$0.14, and 13.3$`\pm `$0.3 Hz and no broken power law needed (a weak, single power law at low frequencies was used instead). This is reminiscent of the complex of harmonically related QPOs that was found during later observations (see Section 4.3.2). A small frequency shift between different observations would in this already be enough to smooth out all but the most significant peaks.
On MJD 51145 the source had moved considerably down branch I, towards the DBB curve in the CD (HC$``$0.01). The count rate had dropped from $``$530 counts s<sup>-1</sup>, in the previous observations, to $``$340 counts s<sup>-1</sup>. The power spectrum did not show any QPO and was fitted with a power law with a strength of less than 1% ($`\alpha `$ fixed to 1). Although the count rate actually went down, the weak noise, the absence of QPOs, and the softer X-ray spectrum indicate that the source had changed from the IS to the HS.
On MJD 51147 the location in the CD was close to that of the first three observations, on MJD 51139–51143, and the count rate had increased again to $``$450 counts s<sup>-1</sup>. The power spectrum looked similar to that of MJDs 51139–51143, but the more complex shape of the noise made it necessary to use a fit function comprised of a power law, a power law with an exponential cutoff (for the low frequency noise), and a Lorentzian (for a QPO around 7 Hz). The power law component had an rms amplitude of 1.31$`{}_{0.26}{}^{}{}_{}{}^{+0.42}`$% and an index of 1.2$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$. The cutoff power law had an rms amplitude of 2.6$`\pm `$0.2% (0.01–1 Hz), $`\nu _{cutoff}=7\pm 1`$ Hz, and an index of 0.0$`\pm `$0.1. The QPO at 7.4$`\pm `$0.3 Hz had a FWHM of 2.3$`{}_{0.5}{}^{}{}_{}{}^{+0.7}`$ Hz and an rms amplitude of 2.6$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$%. Like in the MJD 51139–51143 observations, the power spectrum showed a strong energy dependence. The source had probably returned to the IS.
On MJD 51150 the source was very close to the MJD 51145 observation in the CD, although the count rate was somewhat higher ($``$420 counts s<sup>-1</sup>). The power spectrum showed no QPOs, and the power law noise was weak (less than 1%), which is typical for a HS.
After MJD 51150 the source moved closer to the DBB curve, along branch I. From MJDs 51152 to 51178 (HC$``$0, SC$``$0.5–0.13) the individual power spectra could be described by a single power law with an index of $``$1 and a strength that increased from $``$0.5% to $``$2% rms. During this period the count rate increased from $``$470 to $``$7950 counts s<sup>-1</sup>.
### 4.2 The broad maximum and the flares
During the broad maximum of the source five flares were observed (see Figure 2). We compared the averaged power spectrum of these flares with that of the parts between the flares (hereafter ‘interflares’). For the interflare observations we took all observations between MJDs 51170 and 51237 that were neither in a flare, nor within one observation from a flare. Although their hard color (Fig. 2h) and 0.01–0.1 Hz noise (Fig. 5a) showed behavior similar to that of the flares, the observations on MJD 51220 were not included in either category, since they did not show up as a flare in the light curves and the CD.
Figure 8 shows the 1/128–128 Hz power spectra (2–60 keV) of the combined flare observations and the combined interflare observations. For reasons explained in Section 2, we did not use a single power law to fit the noise (as we did for the individual observations). Instead, we used a broken power law, with a Lorentzian for a QPO around 15–18 Hz. The power spectral fit parameters for the flare and interflare observations are given in Table 2. Both power spectra show a clear break around 3 Hz. The noise component in the power spectrum of the flares is steeper than that in the interflare power spectrum, both below and above the break. The rms normalized power spectra of the flares and interflares cross each other around 1 Hz (see Figure 8), with the 0.01–1 Hz noise being stronger in the flares, and the 1–10 Hz noise being stronger in the interflares (see Table 2).
Figure 5a shows the strength of the 0.01–0.1 Hz noise as a function of time. When comparing this figure with Figure 2, it is evident that increases in the strength of the 0.01–0.1 Hz noise occurred at the times of the hard flares. Note that we used the 0.01–0.1 Hz noise instead of the 0.01–1 Hz noise, since the effect is more pronounced in the 0.01–0.1 Hz range. Figure 10 shows the energy dependence of the 0.01–0.1 Hz (Fig. 10a) and 1–10 Hz (Fig. 10b) noise, for both the flares (bullets) and the interflares (diamonds). From this figure it is apparent that the fractional rms energy spectrum of the 0.01–0.1 Hz noise in the flares was softer than that in the interflares (Fig. 10c), as opposed to the spectrum of the source itself, which was harder in the flares than in the interflares (Fig. 2g,h). Apart from it being stronger in the interflares below 3 keV, the 1–10 Hz noise showed (Fig. 10d) no clear spectral change between the flares and interflares. Although the detections of the noise are very significant, we note that the amplitudes are compatible with the HS observations of other sources. The large amount of data, and the relatively high count rates of XTE J1550–564 made it possible to study the HS power spectra in much higher detail than was possible in other sources before.
Apart from the difference in the noise, some of the individual power spectra during the flares also showed a QPO around 18 Hz. The only flare in which the QPO was not significantly detected was flare 3, the softest flare. Figure 8a shows the 2–60 keV power spectrum of the combined flares (including flare 3), with the QPO at a frequency of 17.87$`\pm `$0.17 Hz. The energy spectrum of the QPO in the flares is shown in Figure 10a. In the highest energy band (13.1–60 keV) the probable second harmonic of the QPO was detected at 35.5$`\pm `$2.0 Hz, with a FWHM of 10$`{}_{3}{}^{}{}_{}{}^{+4}`$ Hz and an rms amplitude of 4.6$`{}_{0.6}{}^{}{}_{}{}^{+1.2}`$%. In the two lower bands only upper limits could be determined to the rms amplitude of the harmonic: 0.17% (2–6.5 keV) and 0.8% (6.5–13.1 keV). We also searched for a QPO around 18 Hz in the combined interflare power spectra; a QPO was found at 15.6$`{}_{0.3}{}^{}{}_{}{}^{+0.2}`$ Hz. Its energy spectrum is shown in Figure 10b. The energy spectra of the 17.87 Hz QPO, its harmonic, and the 15.6 Hz QPO are consistent with each other.
### 4.3 Branch II - The Very High State
As mentioned in Section 3, the source did not return to the DBB curve after the fifth flare. On MJD 51237 (HC$``$0.001, SC$``$0.13) the source reached the position closest to the DBB curve. After that, on MJDs 51239 and 51240, it started to move away from the DBB curve, along branch II. On both these days the 18 Hz QPO was found again. On MJD 51241 the source had moved further away from the DBB curve than during the previous flares: HC$``$0.008, SC$``$0.24. The power spectrum of this observation (see Figures 11a, 12, and 20) was rather different from those seen during the HS and the flares. A broad peak around 6 Hz was found, and also a peak around 280 Hz. The presence of the 6 Hz and 280 Hz peaks, the reappearance of the hard component in the energy spectrum, and the much higher count rate compared to branch I suggest that the source had entered the VHS, as was already reported by Homan, Wijnands and van der Klis (1999). This VHS lasted from MJD 51241 until MJD 51259 (the entire branch II in the CD).
In the power spectra of nearly all the VHS observations one or more QPOs are present around 6 Hz. Although the frequency of this QPO varied between 5 and 9 Hz, for reasons of clarity this QPO will be referred to as ‘the 6 Hz’ QPO. Some observations also show a single QPO with a frequency between 100 and 300 Hz. Based on the Q-value (the QPO frequency divided by the QPO FWHM) of the 6 Hz QPO, and its harmonic structure, Wijnands, Homan and van der Klis (1999) distinguished two types of VHS power spectra: one type with a relatively broad ($`Q<3`$) 6 Hz QPO and sometimes a harmonic at 12 Hz (type A low-frequency QPOs; see Sec. 4.3.1), and one with a relatively narrow ($`Q>6`$) 6 Hz QPO, with harmonics at 3 and 12 Hz (type B low-frequency QPOs; see Sec. 4.3.2). We decided to divide type A into two subclasses; one in which the 6 Hz QPO is strong (rms $`>`$ 2%) and the harmonic at 12 Hz was detected (type A-I), and one in which the 6 Hz QPO was weak (rms $`<`$ 2%) and no harmonic was detected (type A-II). In addition to type A and B a third type, type C, was introduced by Sobczak et al. (2000b), which mainly occurred during the first part of the outburst. Its harmonic structure is similar to that of type B, but the 6 Hz QPO has a higher Q-value (Q$``$10) and its time lag behavior is different. The only type C observation, found by Sobczak et al. (2000b), during the second part of the outburst (MJD 51250) was already classified as an odd type B observation by Wijnands, Homan and van der Klis (1999). Since this observation and the one on MJD 51254 showed odd behavior, compared to the other VHS observations, they will be discussed separately in Section 4.3.3. Figure 11 shows representative 1/128–128 Hz 2–60 keV power spectra of type A-I (Fig. 11a, MJD 51241), A-II (Fig. 11c, MJD 51244), B (Fig. 11b, MJD 51245), and C (Fig. 11d, MJD 51250).
In Table 3 the types of all observations on branch II can be found. In the remainder of this section we describe the three types of low frequency power spectra and the high frequency QPOs in more detail.
#### 4.3.1 Type A observations
On MJD 51241 (HC$``$0.008) the first observation in the very high state was made. The power spectrum was of type A-I, and it can be regarded as representative for the other type A-I power spectra. It is therefore the only type A-I power spectrum that will be discussed in detail in this paper. Figure 12 shows the power spectrum of MJD 51241 in four energy bands. Apart from a QPO around 6 Hz, there was some excess around 11 Hz, which in the high energy power spectra showed up as a QPO. In the 13.1–60 keV band, the 6 Hz peak had almost disappeared. The two peaks seemed to be harmonically related, but fitting them simultaneously in the 2–60 keV band gave frequencies that were not consistent with the two being harmonically related: 5.93$`\pm `$0.03 Hz and 10.41$`\pm `$0.13 Hz. However, when comparing the frequency of the lower frequency QPO in the 2–6.5 keV band (5.85$`\pm `$0.03 Hz) with that of the higher frequency QPO in the 13.1–60 keV band (11.52$`\pm `$0.19 Hz) we found a ratio of 1.97$`\pm `$0.03, which suggests an harmonic relation (see also Wijnands, Homan and van der Klis (1999)). The reason that the frequencies differed so much between the energy bands might be that the chosen fit function was not appropriate, or that an extra component was present between the two QPOs. The FWHM of the 5.85 Hz QPO in the 2–6.5 keV band was 2.41$`\pm `$0.01 Hz, and that of the 11.52 Hz QPO in the 13.1–60 keV band 7.23$`\pm `$0.7 Hz. Their rms amplitudes in the 2–60 keV band were 3.00$`\pm `$0.05% and 2.64$`{}_{0.08}{}^{}{}_{}{}^{+0.09}`$%. At low frequencies a noise component was present. It could be fitted with a single power law with an index of 0.90$`\pm `$0.03, and a strength of 1.31$`\pm `$0.03% rms. Figure 13 shows the photon energy spectra of the two low frequency QPOs and the noise component. The frequencies of the low frequency QPOs were not fixed, for reasons explained above. The 5.8 Hz QPO first increased in strength with energy, but above 10 keV it dropped by a factor of $``$2.5. Its harmonic showed a strong increase with photon energy, from $``$1% rms in the lowest energy bands to more than 11% rms in the highest band. The 0.01–1 Hz noise had a relatively flat energy spectrum, with strengths between 1% and 2% rms, although it became slightly weaker ($``$0.9%) above 10 keV. Selections were made on time, color and count rate, but no significant dependencies were found.
The next observation, on MJD 51242, was located close to the previous observation in the CD, and its power spectrum was also very similar. Again two low frequency QPOs were found, with frequencies of 5.71$`\pm `$0.01 Hz (2–6.5 keV) and 11.2$`\pm `$0.3 Hz (13.1–60 keV), FWHM of 3.1$`\pm `$0.4 Hz and 8$`\pm `$1 Hz, and rms amplitudes (2–60 keV) of 2.53$`{}_{0.12}{}^{}{}_{}{}^{+0.11}`$% and 2.51$`{}_{0.14}{}^{}{}_{}{}^{+0.17}`$% respectively.
In the power spectrum of MJD 51244 (see Figure 11c) a QPO was found with a frequency of 8.5$`\pm `$0.3 Hz, a FWHM of 3.8$`{}_{0.7}{}^{}{}_{}{}^{+0.9}`$ Hz, and an rms amplitude of 1.34$`{}_{0.12}{}^{}{}_{}{}^{+0.13}`$% (2–60 keV). The QPO was considerably weaker than the 6 Hz QPOs in the two previous observations ($``$3% and $``$2.5%, 2–60 keV), and no sub- or second harmonics were found. The energy dependence of the QPO was rather steep, but in the highest band only an upper limit could be determined: 0.8$`\pm `$0.2% (2–6.5 keV), 2.9$`\pm `$0.3% (6.5–13.1 keV), and $`<`$3.7% (13.1–60 keV). Although the Q-value is similar to that of the previous two observations the above characteristics set this observation apart, and we therefore defined it to be of type A-II. In the CD the observation was located further along branch II, away from the DBB line and towards a stronger power law spectral component (HC$``$0.13).
The next type A observations occurred on MJD 51246, after a type B observation on MJD 51245 (see Sec. 4.3.2). In the CD it was located close to the MJD 51244 observation. A QPO was found at 7.72$`\pm `$0.13 Hz, with a FWHM of 3.1$`\pm `$0.4 Hz. Since the QPO was rather weak (1.20$`\pm `$0.06% rms) and no harmonics were found, it was classified as type A-II.
During MJDs 51247–51254 the source moved further up branch II in the CD, and the power spectra only showed type B and C QPOs (see Secs. 4.3.2 and 4.3.3). Type A QPOs reappeared on MJD 51255, when the source returned to a location in the CD close to the other type A observations (see Figure 3: HC$``$0.01, SC$``$0.2). From MJD 51255 to 51258.5 the source showed both type A-I and A–II QPOs, with similar properties as those in the beginning of the VHS. The power spectrum of MJD 51258.9 showed a broad feature around 6 Hz that had a width larger than 6 Hz, and the power spectrum of MJD 51259 showed a QPO at $``$3 Hz with a FWHM of $``$1 Hz, and a broad (FWHM$``$12 Hz) peak around 9 Hz. We classified these two power spectra as type A, but it was not clear of what sub-type they are; both broad peaks might have been unresolved pairs of harmonics.
Figure 14 shows the frequency of all VHS (branch II) QPOs as a function of the hard color, which is a good measure of the position along branch II, as can be seen from Figure 3. The frequencies of the broad peaks in the MJD 51258.9 and MJD 51259 power spectra are not included. Note that type A-II (triangles) observations were located both at lower and at higher hard colors with respect to the type A-I (circles) observations in Figure 14.
#### 4.3.2 Type B observations
The first type B observation occurred on MJD 51245. The source had moved further up branch II (HC$``$0.016), compared to the previous (type A) observations. Figure 15 shows the power spectrum of this observation in four energy bands. A very sharp QPO was present around 6 Hz, with harmonics around 12 and 18 Hz, and a sub-harmonic around 3 Hz. There were also indications for a peak around 24 Hz, but our fits showed it was not significant. Fitting the power spectrum with a power law and Lorentzians (six in total) gave a poor result ($`\chi _{red}^2=2.6`$, $`d.o.f.=230`$). We tried using a power law and Gaussian functions (again six) instead, which improved the quality of the fit ($`\chi _{red}^2=1.4`$, $`d.o.f.=230`$, see also Wijnands, Homan and van der Klis (1999)). The two extra peaks, at $``$0.2 Hz and at $``$1.25 times the frequency of the 6 Hz QPO, were added to the fit function to account for a low frequency component, and for the shoulder of the 6 Hz QPO, respectively. The fit results (2–60 keV) for the four QPOs and the shoulder component are given in Table 4. The QPOs are not perfectly related harmonically, most likely because the fit function did not describe the data well enough. The 0.01–1 Hz noise was fitted with a single power law, with $`\alpha =1.8\pm 0.1`$ and an rms amplitude of 5.6$`{}_{0.2}{}^{}{}_{}{}^{+0.3}`$%. The photon energy spectra of the various power spectral components are shown in Figure 16. Except for the 3 Hz QPO, and the noise component, all QPOs showed a considerable increase in strength with photon energy. The 0.01–1 Hz noise only showed a weak increase, and the 3 Hz QPO behaved similar to the 6 Hz QPO in the type A-I power spectra, in that it seemed to become weaker above 10 keV. Selections were made on color, time and count rate. It was found that the harmonic at 18 Hz was more significant at low hard colors.
Other type B power spectra were found between MJD 51247 and MJD 51253, with QPOs that were similar to those on MJD 51245. Their 6 Hz QPOs had frequencies between 5.3 Hz and 6.1 Hz, rms amplitudes between 3.3% and 3.4%, and Q-values between 6.2 and 7.3. The 3 Hz QPOs had rms amplitudes between 1% and 2.3% and showed a weak trend of an increase with hard color. The MJD 51245 observation remained the only one in which the harmonic around 18 Hz was significantly detected. Figure 14 shows the frequencies of the type B QPOs (squares) as a function of the hard color.
#### 4.3.3 Special cases: Strong noise on MJD 51250 and the transition on MJD 51254
Although, based on the Q-value ($``$8) of the 6 Hz QPO and the harmonic content, the power spectrum of the MJD 51250 observation was classified as type B by Wijnands, Homan and van der Klis (1999), they also found that the time lag behavior of this observation was quite different from that of the other type B observations. Based on this time lag behavior the observation was classified as type C by Sobczak et al. (2000b), making it one of only two (see below for the second) type C observations during the second part of the outburst. More deviations from the type B power spectra were found; in addition to the power law noise (1.6$`\pm `$0.2 %) a strong noise component was present at 0.1–1 Hz (see Figure 11d; also Wijnands, Homan and van der Klis (1999)). This noise component, which we fitted using a zero-centered Lorentzian with a width of $``$3 Hz, was present in all the energy bands, with rms amplitudes of 13.1$`\pm `$0.3% (2–60 keV), 3.3$`\pm `$1.2% (2–6.5 keV), 18.6$`\pm `$0.5% (6.5–13.1 keV), and 28.0$`\pm `$0.8% (13.1–60 keV). Four low frequency QPOs were found in the 2-60 keV band at (rms amplitudes in brackets) 1.7$`\pm `$0.6 Hz (2.74$`{}_{0.45}{}^{}{}_{}{}^{+0.77}`$%), 3.35$`\pm `$0.03 Hz (5.8$`\pm `$2.0%), 6.68$`\pm `$0.15 Hz (6.7$`\pm `$0.2%), and 13.68$`\pm `$0.15 Hz (3.2$`{}_{0.2}{}^{}{}_{}{}^{+0.3}`$%). When comparing these numbers with those of the type B observations, it shows that the rms amplitudes of the 3 Hz and 6 Hz QPOs are, respectively, a factor $``$2.5–6 and $``$2 higher than in the type B observations. The QPO frequencies are shown as stars in Figure 14. In the ASM and PCA light curves the MJD 51250 observation is clearly visible as a dip (see Figures 1 and 2), with a count rate of only $``$5150 counts s<sup>-1</sup>, compared to $``$8200 counts s<sup>-1</sup> on MJD 51249 and $``$7025 counts s<sup>-1</sup> on MJD 51253. This dip is strongest at low energies, causing a hardening of the spectrum a (see Figure 2).
Figure 17 shows the 2–60 keV light curve and the color curves of the MJD 51254 observation. Clearly visible is the jump in count rate that occurred around 1100 s after the start of the observation. The soft color seemed to be unaffected by this change, and though the hard color showed a small change ($``$10%) it was more gradual than the change in count rate. It should be noted that the observed transition is not related to the temporary gain change that was applied to the PCA later during this observation (around t=3100 s; not shown here).
Figure 18 shows the power spectra from before (0–1000 s) and after (1500–3000 s) the jump in the 2–60 and 13.1–60 keV bands. The 2–60 keV power spectrum before the jump showed a broad noise component around a few Hz, that was fitted with a power law with an exponential cutoff. It had a strength (1–100 Hz) of 6.7$`\pm `$0.1% rms, a power law index ($`\alpha `$) of $``$0.7$`\pm `$0.1, and a cutoff frequency of 3.9$`\pm `$0.3 Hz. In the 13.1–60 keV band the strength of this component was 12.5$`{}_{0.7}{}^{}{}_{}{}^{+0.9}`$% rms. In that same band we found a QPO at 9.8$`\pm `$0.1 Hz, with an rms amplitude of 8.6$`{}_{0.5}{}^{}{}_{}{}^{+0.6}`$% and a FWHM of 3.2$`\pm `$0.5 Hz. The post-jump 2–60 keV power spectrum showed a similar noise component as before the jump, though somewhat weaker (4.7$`\pm `$0.1% rms), with two QPOs superimposed on it, at 3.17$`\pm `$0.02 Hz and 6.14$`\pm `$0.03 Hz. These QPOs had rms amplitudes and FWHM of 2.3$`\pm `$0.1% and 0.71$`\pm `$0.06 Hz (3.17 Hz QPO), and 3.08$`\pm `$0.08% and 1.23$`\pm `$0.06 Hz (6.14 Hz QPO), respectively. In the post jump 2–6.5 keV power spectrum an additional QPO at 1.77$`\pm `$0.04 Hz was found (3.7$`\sigma `$) when the high count rate part was selected. Both before and after the transition power law noise was present: respectively, 3.4$`\pm `$0.2% rms with $`\alpha =0.92\pm 0.02`$ (before) and 3.8$`\pm `$0.1% rms with $`\alpha =1.04\pm 0.03`$ (after).
The fast transition in the power spectrum can be seen in Figure 19, which shows the dynamical power spectrum in the 13.1–60 keV band. The time scale for the change in the power spectrum is similar to that of the transition in the 2–60 keV light curve. It was not possible to track the QPO across the transition since it became weaker during the transition.
Wijnands, Homan and van der Klis (1999) and Sobczak et al. (2000b) classified the power spectrum after the jump as type B. Indeed, the strengths of the 3 and 6 Hz QPOs were consistent with those in the other type B observations. On the other hand, the Q-value of the 6 Hz QPO was only 5, and the 5% rms noise component under the QPOs was not seen in other type B observations. The type of the power spectrum before the jump is not clear either. The hardness of that part of the observation suggests type B or C, but the Q-value of the 9.8 Hz QPO was only $``$3. The strength of that QPO was lower than that of the type B and C 6 Hz QPOs in the same energy band ($``$11% rms), but higher than that of the type B and C 12 Hz QPOs (5–6% rms). Since the power spectrum showed a strong noise component, and the 2–60 keV count rate was lower than that of the type B part it was most likely of type C. The QPO frequencies of both parts are shown in Figure 14 (the part before \[HC=0.205\] as type C, the part after \[HC=0.195\] as type B).
The exceptional cases of low frequency QPOs presented in this section clearly demonstrate that the A, B (and C) classification. which works well for the majority of the observations, is not able to unambiguously describe all of them.The characteristics on which this classification is based (Q-value, harmonic content, time lags) show strong correlations, but deviations from the usual correlations do occur.
#### 4.3.4 High frequency QPOs
During several observations in the VHS (branch II), QPOs were found with frequencies between 100 and 300 Hz. An example can be seen in Figure 20, which shows the 284 Hz QPO found in the power spectrum of MJD 51241. The frequencies of the high frequency QPOs ($`\nu _{HF}`$) are given in Table 3, and the locations in the CD of the observations in which they were found are indicated in Figure 21. Note that we only report QPOs whose single-trial significance exceeds 3$`\sigma `$. It can be seen from Figure 21 that $`\nu _{HF}`$ is related to the location in the CD. It decreased from 284 Hz to 102 Hz as the source moved up branch II, and increased again to 280 when it moved down this branch. Figure 14 more clearly shows that $`\nu _{HF}`$ decreased as the hard color increased. In the high energy bands, the QPOs tended to be stronger when they had a frequency around 280 Hz, as can be seen from Table 3. The Q-values of the QPOs were not related to those of the low frequency QPOs, and had values between 5.6 and 13.
We measured time lags for the $``$282 Hz QPO in the combined Fourier spectra of the observations on MJDs 51241, 51242, and 51255. Lags were measured between three energy bands, in the frequency range 272–292 Hz. All lags were consistent with being zero: 0.00$`\pm `$0.11 ms (2–6.5 keV and 6.5–13.1 keV), $``$0.08$`\pm `$0.13 ms (2–6.5 keV and 13.1–60 keV), and $``$0.08$`\pm `$0.04 ms (6.5–13.1 keV and 13.1–60 keV), where a positive number means that the photons in the second band lag those in the first one. A time lag analysis of the low frequency QPOs in the VHS (branch II) can be found in Wijnands, Homan and van der Klis (1999), Sobczak et al. (2000b), and Cui, Zhang and Chen (2000).
Figure 22 shows the frequency of several low frequency QPOs ($`\nu _{LF}`$) plotted against $`\nu _{HF}`$, for those observations where they were detected simultaneously. We also included the values for the high frequency QPO that was observed on branch III (represented by the diamond; see Section 4.4). The 123 Hz QPO on MJD 51254 was only found in the data taken after the count rate jump (Section 4.3.3). This data included a $``$600 s interval (with different PCA gain settings) that was not used for the analysis of the low frequency QPOs in Figure 14. The low frequency QPOs plotted at $`\nu _{HF}=123`$ Hz in Figure 22 therefore have a slightly different frequency than those of the same observation in Figure 14. Four lines could be fitted to the data points, with the first, third and fourth line having slopes that were, respectively, 0.50$`\pm `$0.02, 2.17$`\pm `$0.10, and 4.3$`\pm `$0.5 times the slope of the second line. This is consistent with the four lines representing the fundamental and second, fourth and eighth harmonics. Note that the four lines do not pass through the origin and cross each other around $`\nu _{LF}=0`$ Hz and $`\nu _{HF}=75`$ Hz. The only two points that were not fitted by these four lines were the sixth harmonic in the MJD 51245 observation ($`\nu _{HF}`$=178 Hz) and the sixteenth harmonic in MJD 51250 observation ($`\nu _{HF}`$=102 Hz). These components were only observed once, and therefore no fits could be made. The four lines can be used to connect the low frequency QPOs in Figure 14. For example, using the second line in Figure 22, it can be seen that the type A-I 10–12 Hz QPOs are related to the type A-II 8–9 Hz QPOs, the type B $``$6 Hz QPOs, the 3.1 Hz (type B?) QPO on MJD 51254, and the 1.7 Hz type C QPO on MJD 51250. The QPOs that lie on the second line in Figure 22, and those that based on similarities in the power spectrum and hard color are expected to, have been represented by the filled symbols in Figure 14. The filled symbols show that the frequency of the low frequency QPOs decreases with hard color, like that of the high frequency QPO.
### 4.4 The Decay
On MJD 51260 (indicated by BEGIN in Figure 3b) the power spectrum showed no QPOs; it could be fitted with a single power law, with a strength of 0.76$`\pm `$0.05% rms and an index of 1.1$`\pm `$0.01. This weak noise, the absence of QPOs, and the relatively soft colors suggest that the source had returned to the HS.
The power spectrum of the next observation (MJD 51261) showed a noise component with a similar strength (0.56$`\pm `$0.07% rms), but also a QPO at 17.0$`{}_{0.3}{}^{}{}_{}{}^{+0.5}`$ Hz, with a FWHM of 2.5$`{}_{1.1}{}^{}{}_{}{}^{+1.9}`$ and a rms amplitude of 1.02$`{}_{0.16}{}^{}{}_{}{}^{+0.20}`$%. Based on the hardness at which this QPO is found, its frequency, and its FWHM, it may be related to the 15.6/17.9 Hz QPO that was found in the flare/interflare observations during MJD 51170–51237 (see section 4.1). In the next few observations (MJD 51263–51267) no QPO around 17 Hz was found, and the power spectra could be fitted with single power laws, with strengths between 0.4% and 1.2% rms, typical for HS.
On MJD 51269 the source had started to move up branch III. Two QPOs were found in the power spectrum of that observation: at 4.54$`\pm `$0.15 Hz (1.7$`\pm `$0.2% rms, FWHM=1.5$`{}_{0.4}{}^{}{}_{}{}^{+0.5}`$ Hz) and 9.6$`\pm `$0.6 Hz (2.0$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$% rms, FWHM=4.5$`{}_{1.5}{}^{}{}_{}{}^{+2.4}`$ Hz). The noise at low frequencies was fitted with a single power law, with a strength of 1.32$`\pm `$0.07% rms and an index of 0.8$`\pm `$0.1. Based on their Q-values, the QPOs are either of type A-I or A-II; the strength of the QPOs (and their frequency) suggests type A-II, whereas the detection of an harmonic suggest type A-I (see section 4.3.1).
The next two observations (MJDs 51270 and 51271) were located near the top of branch III. Their power spectra were very similar. The MJD 51270 power spectrum showed a QPO at 8.9$`\pm `$0.1 Hz (1.8$`{}_{0.1}{}^{}{}_{}{}^{+0.2}`$% rms, FWHM=2.1$`{}_{0.3}{}^{}{}_{}{}^{+0.5}`$ Hz) and a broad peak around 2 Hz that was fitted with a Gaussian at 1.8$`\pm `$0.1 Hz (3.2$`\pm `$0.2% rms, FWHM=2.9$`\pm `$0.4 Hz) plus an exponentially cutoff power law (4.4$`\pm `$0.3% rms, $`\alpha =`$1.6$`\pm `$0.6, $`\nu _{cutoff}=4\pm `$1 Hz). The power spectrum of MJD 51271 showed a QPO at 9.05$`\pm `$0.12 Hz (1.8$`{}_{0.2}{}^{}{}_{}{}^{+0.3}`$% rms, FWHM=1.4$`{}_{0.4}{}^{}{}_{}{}^{+0.5}`$ Hz) and a broad peak around 2 Hz that was fitted with a Gaussian at 1.7$`\pm `$0.3 Hz (3.8$`{}_{0.3}{}^{}{}_{}{}^{+0.4}`$% rms, FWHM=3.74$`\pm `$0.5 Hz) plus an exponentially cutoff power law (4.5$`\pm `$0.4% rms, $`\alpha =`$2.3$`\pm `$1.0, $`\nu _{cutoff}=3\pm `$1 Hz). The combined 1/128–128 Hz power spectrum of the two observations is shown in Figure 23b. The strength of the 0.01–1 Hz noise, which was fitted with a power law, was $``$1.5% rms, but it should be noted that some of the power in the 0.01–1 Hz range was absorbed by the Gaussian and the exponentially cutoff power law. The two $``$9 Hz QPOs had relatively high Q-values (4.2 and 6.5), which suggests that they were of type B; this seems to be confirmed by the shape of the power spectra at higher energies; Figure 24 shows the 6.5–60 keV power spectrum of MJD 51271, which could be fitted with a power law and QPOs at 3.1$`\pm `$0.1 Hz, 5.7$`\pm `$0.2 Hz, 9.0$`\pm `$0.1 Hz, and 12.5$`\pm `$1.0 Hz. This is reminiscent of the type B QPO found on branch II, and the IS power spectrum shown in Figure 7. In the combined 2–60 keV power spectrum of the two observations a QPO at 251$`\pm `$3 Hz was found. It had an rms amplitude of 2.21$`\pm `$0.15% and a FWHM of 42$`\pm `$6 Hz. Its location in Figure 22 is shown by a diamond. Although the frequency of the QPO lies in the $`\nu _{HF}`$ range found on branch II, the count rate at which is was found was considerably lower ($``$1350 counts s<sup>-1</sup> compared to 4700–8300 counts s<sup>-1</sup> on branch II).
The power spectrum of the next observation (MJD 51273), which in the CD was located close to the MJD 51269 observation, showed a QPO at 4.52$`\pm `$0.13 with an rms amplitude of 1.2$`\pm `$0.02% and a FWHM of 1.1$`{}_{0.3}{}^{}{}_{}{}^{+0.5}`$ Hz. The QPO is most likely of type A, based on its strength and lack of harmonic structure. The low frequency noise had a strength of 1.20$`\pm `$0.05% rms. The combined power spectrum of MJDs 51269 and 51273 is shown in Figure 23a.
The observation on MJD 51274 showed no QPOs, and the 0.01–1 Hz noise had an rms amplitude of less than 0.4%. During MJD 51274–51280 (HC$``$0.007, SC$``$0.13) the source showed similar power spectra that, when combined, were fitted with a single power law with a strength 0.90$`\pm `$0.09% rms and an index of 1.0$`\pm `$0.2, which is typical for the HS.
Between MJD 51283 and MJD 51298 XTE J1550–564 traced out branch IV in the CD. At the top of the branch (SC$`>`$0.35) the power spectra showed a QPO and a peaked noise component. These were not found at the bottom of the branch (SC$`<`$0.35). The combined power spectrum of the bottom of branch IV (Fig. 23c) was fitted with a power law with an rms amplitude of 2.4$`\pm `$0.2% and an index of 0.7$`\pm `$0.1. In the combined power spectrum of the top of branch IV (Fig. 23d) a QPO was found at 11.3$`\pm `$0.5 Hz, with an rms amplitude of 4.1$`{}_{0.5}{}^{}{}_{}{}^{+0.7}`$% and a FWHM of 2.9$`{}_{0.8}{}^{}{}_{}{}^{+1.1}`$ Hz. A peaked noise component was present below 10 Hz. It was fitted by a Lorentzian with a frequency of 3.0$`\pm `$0.3 Hz, an rms amplitude of 10.1$`{}_{0.8}{}^{}{}_{}{}^{+0.9}`$%, and a FWHM of 5.7$`\pm `$0.9 Hz. The 0.01–1 Hz noise was fitted with a power law that had a strength of 2.8$`\pm `$0.4% and an index of 0.3$`\pm `$0.1.
Both branch III and IV showed behavior that was similar to that seen in the VHS (branch II). Since the count rates were lower than on the branch II, these branches were probably IS (at least, when QPOs were seen).
Between MJD 51299 and MJD 51318 branch V was traced out in the CD. It reached much harder colors than before, and the movement up the branch was accompanied by a considerable increase in the strength of the low frequency noise, as can be seen from Figure 5. There was a clear difference between the power spectra at the top and bottom of the branch. The combined MJD 51299–51306 (bottom part of branch V, SC$`<`$0.6) power spectrum was fitted with a power law with an rms amplitude of 8.9$`\pm `$1.0% and an index of 0.7$`\pm `$0.1. A single power law yielded a poor $`\chi _{red}^2`$ (3.2 for d.o.f.=67) for the combined MJD 51307–51318 (top of branch V, SC$`>`$0.6) power spectrum, and a broken power law was used instead (see Figure 25). Its rms amplitude was 15.9$`\pm `$0.3%, with $`\nu _{break}=0.9\pm 0.1`$ Hz, $`\alpha _1=0.33\pm 0.04`$, and $`\alpha _2=1.2\pm 0.1`$ ($`\chi _{red}^2=1.1`$ for d.o.f.=67). The strength and shape of the noise, and the spectral hardness suggests that the bottom of branch V the source was in the IS, and that at the top of branch V the source was in the LS.
## 5 Radio Observation
On 1999 March 11 (MJD 51248) we observed the radio counterpart of XTE J1550–564 (Campbell-Wilson et al., 1998) with the Australia Telescope compact array (ATCA), in a high-resolution 6 km configuration. Observations were made simultaneously at 6.3 and 3.5 cm, and at 21.7 and 12.7 cm, in order to obtain broad band spectral coverage. The observations were interleaved with those of a nearby reference source B1554-64, for phase calibration every 25 min. The source was clearly detected at all four wavelengths; the mean flux densities at 21.7, 12.7, 6.3 & 3.5 cm were, respectively, 5.1, 3.0, 2.8, and 1.9 mJy (errors $``$10%). The four flux densities were fitted with a power law corresponding to a spectral index ($`\alpha =\mathrm{\Delta }logS_\nu /\mathrm{\Delta }log\nu `$) of $``$0.53$`\pm `$0.12. The location of the MJD 51248 RXTE observation is indicated with by ‘ATCA’ in Figure 3.
## 6 Discussion
In this section we present a discussion of our results. We start by briefly summarizing the results. After that the source states and power spectra are discussed.
### 6.1 Summary of Results
In the period of 1998 November 22 to 1999 May 20 XTE J1550–564 showed a wide variety of behavior. To organize the different phenomena and relate them to each other, it is useful to compare the rapid time variability and the energy spectra. An initial division of the observations can be made based on the strength of the broad band (1/128–128 Hz) power, which is shown in Figure 5b. It can be seen that the source alternated between states with low power (a few percent rms) and states with high power (more than a few percent). When comparing this figure with Figure 2 it is obvious that observations with high power were mainly found when the spectrum was hard. These spectrally hard states appeared as branches (I–V) in the color-color (Fig. 3) and hardness-intensity (Fig. 4) diagrams. From the hardness-intensity diagram it is apparent that the hard branches occurred at five distinct count rate levels, and that they were separated by periods that were spectrally soft(er); in the color-color diagram the hard branches lay more or less parallel to the power law curve. The power spectra on the five hard branches often showed QPOs, and in some cases also strong peaked and/or band limited noise (Fig. 11). These properties classify the observations on the branches as VHS, IS, or LS. The power spectra that were not on the hard branches showed noise with strength of a few percent rms and, when combined, a weak QPO around 17 Hz (Fig. 8). These observations can be classified as HS; in the color-color diagram they lay close to the disk blackbody curve. When the source moved up a hard branch the low frequency noise changed from a weak power law to strong band-limited. On branch II this change was accompanied by the QPOs changing from the broad type A, to the narrow types B and C (Fig. 14). On branches II and III we also found high frequency QPOs in the 102–284 Hz range. Their frequencies were anti-correlated with the hardness of the energy spectrum (Fig. 14), and correlated with the frequency of the type A, B, and C low frequency QPOs (Fig. 22).
### 6.2 Source States
In recent years the picture of black hole behavior that emerged from observations was consistent with a one-dimensional scheme, in which four canonical source states were linked by one parameter, usually taken to be the mass accretion rate (see, e.g. Esin, McClintock and Narayan (1997) and Esin et al. (1998) for recent elaborations on this view). In the context of two-component spectral models, often interpreted in terms of emission from an accretion disk and a hot comptonizing medium, this implies that both components contribute to the energy spectrum and power spectrum in amounts that depend strictly on this parameter; if both components were to vary independently, the description of the phenomenology would have to be at least two-dimensional. XTE J1550–564 seems to provide evidence for such two-dimensional behavior. Although the source was observed in all four canonical states, their occurrence was more complex than expected on the basis of a simple relation with the mass accretion rate.
We start by discussing the relation between the source states and the position of the source in the CD and HID (see Figures 3 and 4). The motion of the source through the CD was along branches. One branch (hereafter the soft branch) lay parallel to the DBB curve in the CD, and quite close to it (Fig 3a). Whenever XTE J1550–564 was on or close to this branch (e.g. flares 1-5), it could be classified as being in the HS: it showed soft energy spectra, and the (power law-like) low frequency noise had a strength of only a few percent rms. All the other branches (hereafter hard branches) lay approximately parallel to the power law curve in the CD. In time, these hard branches were traced out one after the other, and all but the last two were clearly separated from each other by intervals that showed HS behavior. When the source was on a hard branch it was in the VHS, IS or LS: the energy spectrum was hard, and the power spectrum showed QPOs and/or strong noise. The hard branches were similar to each other in that the shape of the noise changed from power law like to band limited as the source moved up such a branch (i.e. when it became harder). When the source was on a hard branch still relatively close to the soft branch it would, based on the energy and power spectrum, usually be classified as being in the canonical HS - only further up the hard branches full-fledged VHS, IS and LS behaviour emerged. Canonical LS (variability) behavior was only found at the top of branch V, the branch that reached the hardest colors.
Although the observations on hard branches I, III, and IV were classified as IS, and those on hard branch II as VHS, they had in fact very similar properties, the only difference being the count rate at which they were observed. This was already found for the IS and VHS in other sources, e.g., GS 1124–68 and GX 339-4 (Belloni et al., 1997; Méndez and van der Klis, 1997). We therefore regard the VHS as an instance of the IS, but just the one that happens to be the brightest.
The behavior on the hard branches that were traced out before our observations, during the first part of the outburst, was similar to that during our observations (Remillard et al., 1999a; Cui, et al., 1999). During the first part of the outburst LS behavior (strong band limited noise with a break around 0.1 Hz) was observed only when the hardness was similar to that at the top of branch V (at the start of the outburst, when the count rate was $``$100 times as high as on branch V). Moreover, the source evolved from LS to the HS via a VHS (or IS as we shall henceforth call it), clearly showing the same ordering of states (as a function hardness) as during our observations.
The above shows that as the hardness increased the source evolved from the HS via the IS, to the LS; the hard branches therefore corresponded to HS$``$IS$``$LS transitions, or at least attempts to, since not every branch reached the LS and the source did not always return completely to the HS. Similar conclusions were also drawn by Rutledge et al. (1999), on the basis of a comparative study of 10 black hole candidates. They found that the VHS and IS were spectrally intermediate to the HS and LS, and also that as the hardness increases the noise switches from HS-like to LS-like. They also concluded that transitions between the HS and LS could take place at luminosities/count rates both above and below that of the HS. In our observations transitions between HS and IS were found at around 8000, 1200, 600, and 200 counts s<sup>-1</sup> (see Fig. 4); transitions between IS and LS were found at around 40 (Fig. 4) and 4000 counts s<sup>-1</sup> (during the first part of the outburst). Moreover we observed HS behavior at all count rate between 200 and 10000 count s<sup>-1</sup>. All this argues against a one-dimensional description of the state transitions as a function of the mass accretion rate.
The observations of XTE J1550–564 contradict the old picture of black hole states, in which hard states are only found at the highest and lowest count rates. XTE J1550–564 clearly shows that hard states can be observed at any count rate level. However, some remarks should be made. All the hard states (I–V) were observed during the decay of the source (I during the decay of the first part of the outburst, II–V during the decay of the second part). Also, the intervals between branches II–IV, although they could be classified as HS, had significantly harder spectra than the HS observations during the rise, which were extremely soft. This suggests that the conditions for the presence of the hard spectral component are more favorable during the decay or phases of low count rate. This is supported by the fact that all small flares (1–5) were observed at count rates higher than that of the hard branches. The fact that these flares did not develop into real transitions suggests that the conditions for transitions and development of the hard spectral component are less favorable at the highest count rates. Also, it can be argued that the only LS that was observed during our observations was found at the lowest count rates at the end of the outburst, as expected in the canonical picture of black hole states. On the other hand, a LS was also found during the first part of the outburst when the count rate was at least a factor 100 higher than during the LS at the end of the outburst, clearly showing that LS is not only found at the lowest count rates.
Based on the CD, HID, and power spectral fits one gets the impression that the behavior of XTE J1550–564 is two-dimensional, i.e. at least two (observable) parameters are needed two describe the appearance of the source and to account for the occurrence of the different states. Phenomenologically, the two parameters describing this two-dimensional behavior are the count rate and the spectral hardness. A schematic representation of the behavior of XTE J1550–564 in terms of these parameters is shown in Figure 26. It shows that the states are arranged in a comb-like topology, with the soft (HS) branch being the spine, and the HS$``$IS$``$LS transitions being represented by the teeth. Note that the parameter on the vertical axis is the logarithm of the count rate, and that we decided to show the hard branches as horizontal lines; for reasons of clarity we did not depict the exact movement of the source through the diagram. We also included the location of the flares, which occurred at count rates higher than that of the brightest hard branch.
The most important aspect of XTE J1550–564 is probably the fact that two observable parameters (count rate and hardness) varied too a large extent independently from each other. This suggests that at least two physical parameters underlie this behavior, since this complex behavior is hard to explain within a framework where the appearance of the source is determined by a only single physical parameter (e.g. only mass accretion rate). Physically the two underlying parameters might for example be the mass accretion rate through the disk (roughly increasing with count rate, at least in the HS) and the size of a Comptonizing medium (increasing with spectral hardness). In that case the state of the source is determined by the (relative) size of the Comptonizing medium, with it being small or absent in the HS and growing in size towards the LS. The fact that we see increases in the spectral hardness at many count rate levels suggests that the size of this medium, and therefore also the state of the source, is to a large degree not determined by the accretion rate through the disk. We note that the inner disc radius as derived from variability properties (i.e. QPO frequencies, see Section 6.3.2) correlates well with hardness, suggesting that the Comptonizing medium grows as the inner disk edge moves out. The two physical parameters do not necessarily vary completely independently from each other. For instance, changes in one parameter may be triggered by changes in the other, and certain values of one parameter may restrict the value of the other parameter (e.g. in between the hard branches XTE J1550–564 seemed to become slightly harder towards lower count rates (Fig. 4)).
While previous authors, inspired by the description of black hole spectra in terms of two components, have also discussed black hole phenomenology in terms of two-dimensional diagrams (Miyamoto et al., 1994; Nowak, 1995), the overall picture has in our view been considerably clarified by the clues provided by XTE J1550–564 described in this paper. The basic phenomenology seems to be one where the hard and soft component can vary to a large extent independently from each other. The HS is the name we have given in the past to all cases where the hard component is weak compared to the soft component, and the IS/VHS and LS are unified into cases where the hard component is, respectively, comparable to or dominating the soft component. The difference between the VHS and IS is reduced to a difference in the luminosity of the soft component at which they occur, and the difference between IS/VHS and LS is caused by differences in the relative contributions of the soft and hard components.
This two-dimensional interpretation might very well be applicable to all black hole candidates showing the canonical states. The often observed order of states (VHS$``$HS$``$IS$``$LS) is still consistent with the diagram drawn Figure 26. The fact that black hole state behavior often seems one-dimensional can be explained by considering how the behavior of XTE J1550–564 would have appeared if the data were of lower quality, the source sampling was more infrequent, its distance was larger, and if it had different characteristic time scales for variations in the soft and hard component. Much of its subtle behavior would have been missed or may have been misinterpreted. It is mainly thanks to the combination of source brightness, the quality of the RXTE/PCA data, and the excellent source sampling that we clearly see the two-dimensional nature of its behavior. Since the observations of XTE J1550–564 strongly suggest that mass accretion rate through the disk and state are decoupled, it would be possible to see transients that remain in the same state during a whole outburst. Moreover, one could also see state transitions in sources with a more or less constant mass accretion rate. Suggestions of such behavior have been seen in, respectively, GS 2023+338 (Sunyaev et al., 1991; Terada et al., 1992; Miyamoto et al., 1992) and Cyg X-1 (Zhang et al. (1997), however see Frontera et al. (2000)).
The radio brightness of XTE J1550–564 during the MJD 51248 ATCA observations indicates that an outflow event was going on or had recently occurred. The spectral index suggests that the radio source was optically thin, and observed during the decay of such an outflow event. This event might be associated with the state change on MJD 51237/51239 (the onset of branch II). Although XTE J1550–564 was observed in radio only once (on branch II), observations of other black hole candidates (Fender et al., 1999, 2000) suggest that radio emission is associated with spectrally hard states. The hard branches may therefore correspond to changes in the accretion flow geometry, where an inflow (HS) is gradually accompanied by (or changing into) an outflow (LS). Jet-like outflow models have already been proposed for the VHS in GX 339–4 by Miyamoto and Kitamoto (1991).
### 6.3 Power spectra
Although not always observable in individual observations, QPOs were found on all branches, except for the last one. Several types were found: 1–18 Hz QPOs on the hard branches (type A, B, and C), 15–18 Hz (plus an harmonic) on the soft/HS branch and in the flares, and 100–285 Hz QPOs on the two brightest hard branches.
#### 6.3.1 High frequency QPOs
High frequency QPOs in black hole candidates are a relatively new phenomenon. Previous to XTE J1550–564, they were found in GRS 1915+105 (Morgan, Remillard and Greiner, 1997, 67 Hz) and in GRO J1655-40 (Remillard et al., 1999b, 300 Hz). Remillard et al. (1999a) found high frequency QPOs in XTE J1550–564 in the 161–238 Hz range, during the first part of the outburst. With the observations of the second part of the outburst this range has been expanded to 100–285 Hz. It is obvious that the frequency of the high frequency QPOs in XTE J1550–564 can not be explained by models that predict an approximately constant frequency, e.g. orbital motion at the innermost stable orbit (Morgan, Remillard and Greiner, 1997), Lense-Thirring precession at the innermost stable orbit (Cui, Zhang and Chen, 1998), or trapped-mode disk oscillations (Nowak et al., 1997).
The high frequency QPO was found on two branches; between 100 and 285 Hz on branch II, and at 251 Hz on branch III. The count rates at which it was observed, were much lower on branch III ($``$1350 counts s<sup>-1</sup>) than on branch II ($``$4700-8300 counts s<sup>-1</sup>), which shows that the QPO frequency does not strongly depend on the count rate. A similar effect is also seen for the high frequency QPOs in some neutron star X-ray binaries (Méndez et al., 1999; Ford et al., 2000), where the frequency varies along parallel branches in a frequency–count rate diagram. A certain frequency can therefore be observed at different count rate levels, and a range of frequencies can be found within a relatively small range of count rates. This suggests that if the QPO frequency is related to a certain (variable) radius, this radius varies almost independently from the count rate (and probably mass accretion rate; van der Klis (2000)). Similar two-dimensional behavior as discussed in Section 6.2 may therefore also be present in some of the neutron star sources.
An obvious question to ask is, whether the high frequency QPOs in black hole candidates have the same origin as the kiloHertz QPOs that are observed in the neutron star sources (see van der Klis (2000) for a review). Of course, since the QPOs in the neutron star sources are often observed in pairs, only one (if any) of those two QPOs can have the same origin as the QPOs in black hole sources, which until now have always appeared as single peaks. It is, however, not clear which of the two QPOs that would be; both the lower and upper kHz QPO have Q-values and the rms energy spectra that are consistent with those of the QPO in XTE J1550–564. The frequency ranges in which the high frequency QPOs are observed are 102–284 Hz for the QPO in XTE J1550-564, and 200–1070 Hz for the lower kHz QPO, and 325–1330 for the upper kHz QPO in the neutron stars. Here we combined the kHz QPO data for all neutron star sources in van der Klis (2000). Although the lower and upper kHz QPOs cover a frequency range of a factor of 5.4 and 4.1, respectively, the values for individual neutron star sources are more like that found for XTE J1550–564 ($``$2.8).
Although the high frequency QPO in XTE J1550–564 has parameters that are consistent with those of both the lower and upper kHz QPO (within a simple orbital frequency model), a major difference between XTE J1550–564 and the neutron star sources is the fact that the latter often show two high frequency QPOs. However, this could be explained if one of the two kHz QPOs in the neutron star sources is due to a mechanism that requires the presence of a solid surface.
#### 6.3.2 Low frequency QPOs
On all the hard branches, except branch V, QPOs were found with frequencies between 1 and 18 Hz. Due to the low count rates the quality of the power spectra on branch V was poor, and the presence of QPOs could therefore not be ruled out. The frequencies of the QPOs were not constant, as can be seen from Figure 14 (showing the QPOs found on branch II), but it is not immediately clear from that figure how the frequency of the low frequency QPOs ($`\nu _{LF}`$) depended on the hard color (which is a good measure of the distance along the branch). Although the figure is suggestive of a positive correlation, especially for HC$`<`$0.015, the behavior of the low frequency QPOs during the rise of the first part of the outburst of XTE J1550–564 (anti-correlation with hardness; Cui, et al. (1999)) and the correlation found between the low and high frequency QPOs (see below, and Figure 22) lead us to believe that hard color and frequency were anti-correlated, and that the frequencies of both the high and low frequency QPOs decreased as the source moved up branch II. This contradicts the switch from a correlation into anti-correlation when the hard color passes a certain value, that was reported by Rutledge et al. (1999) for other black hole candidates. We want to stress again that the fact that several harmonics were present and that not always the same harmonic was the strongest one, can easily lead to confusion. On the other hard branches not enough QPOs were observed to confirm the anti-correlation with hardness.
It was usually the QPO that happened to be located between 5 and 8 Hz that was the strongest in the 2–60 keV power spectra of branch II, even though it could be identified with different harmonic components (as can be seen from Figures 14 and 22). Perhaps variations within that frequency range are less prone to damping than outside, or a resonance occurs. The shoulder component that was present at a frequency of 1.25 times that of the 5–8 Hz type B QPOs, has been found before in QPOs in other black hole candidates like GS 1124–68 and GX 339–4 (Belloni et al., 1997), but recently also in the 20-50 Hz QPOs in the neutron star system GX 340+0 (Jonker et al., 2000).
Figures 14 and 22 seem to indicate that the low frequency QPOs evolve from type A via type B into type C, and vice versa. On branch II the type A QPOs were located at the bottom of the branch, and the those of type B and C further along it. This was also seen on branch III for type A and B (Figs. 23 and 24), and to certain extent also on branch I, where indications for type B QPOs were found at a similar hardness as where they were found on branch II (Fig. 7). No direct transitions between type A and B were seen, so it is not known whether such transitions are smooth or abrupt. In the MJD 51254 observation a jump in the count rate was accompanied by a change in the power spectrum that may have been a transition from type C to B. This suggests that the transitions between the different types are quite sudden. Apart from the difference in spectral hardness at which type A and B occurred, it is clear that there are at least two other fundamental differences between the two types. First there is the difference in the Q-value, which is higher in type B. This might also explain why more harmonics are detected in the type B power spectra, since narrower features are easier to identify. Second, there is the difference in time lag spectra (Wijnands, Homan and van der Klis, 1999), which can not be reconciled even when one compares the time lags of similar harmonics (i.e., the type A-I 12 Hz QPO and the type B 6 Hz QPO, which in our view are supposed to be the same harmonic, still have the opposite signs for their time lags). As mentioned already in Section 4.3 the type A QPOs were divided in two types, A-I and A-II. Figure 14 shows that type A-II QPOs (triangles) were found both at higher and lower hardness than type A-I (circles). The A-II subtype should therefore be regarded as a collective rather than a real type. The two filled triangles (HC$``$0.012) can be identified with second harmonics and are for that reason different from the A-I observations, where the fundamental was the dominant harmonic. The nature of the two observations represented by the open triangles in Figure 14 remains uncertain; their frequency is lower than expected on the basis of their hardness and in one case an addiational broad bump was observed in the power spectrum.
We found that the frequencies of the high frequency QPOs and low frequency QPOs that were detected simultaneously are well correlated (Fig. 22); it is mainly on the basis of this that we conclude that the different types (A, B, and C) of low frequency QPOs have the same origin. Similar correlations have also been found for the low and high frequency QPOs in a number of neutron star source (e.g. van der Klis et al. (1996); van der Klis, Wijnands, Horne and Chen (1997); Wijnands et al. (1997); Jonker et al. (1998); Ford and van der Klis (1998); Psaltis et al. (1999); Markwardt, Strohmayer and Swank (1999)) and in neutron star and black hole sources for QPOs and broad noise components (Psaltis, Belloni and van der Klis, 1999). The main correlation that was found for the neutron star and black hole sources extended over a frequency range of 0.1–1200 Hz, and was consistent with the relation $`\nu _{LF}=(42\pm 3)(\nu _{HF}/500Hz)^{0.95\pm 0.16}`$, that was found in the neutron star Z sources (Psaltis et al., 1999; Psaltis, Belloni and van der Klis, 1999). A second correlation was present in Figure 2 of Psaltis, Belloni and van der Klis (1999) that was fitted with $`\nu _{LF}=2.09\times 10^3(\nu _{HF})^{1.46}`$ by di Salvo et al. (2000a) for data of 4U 1728–34. Both relations are plotted in Figure 22 (dashed and dotted lines, respectively), and are apparently not consistent with the data of XTE J1550–564.
The fact that the data in Figure 22 are well fitted with four linear relations that do not pass through the origin excludes models in which $`\nu _{LF}`$ and $`\nu _{HF}`$ are related by a simple power law expression. The four linear fits to the data in Figure 22 cross each other at $`\nu _{LF}`$0 Hz and $`\nu _{HF}`$ 75 Hz. It is not clear what the nature of this $``$75 Hz frequency in a black hole system could be.
Low frequency QPOs were also found on the soft branch (15.6 Hz) and in the high state flares (17.9 Hz). They were much weaker than the low frequency QPOs found on the hard branches and had a rather high Q-value ($``$10). Their frequency was apparently correlated with hard color (assuming that we observed the same harmonical component), unlike that of the hard branch QPOs. The above suggests that the 15.6 Hz and 17.9 Hz QPOs may have a different origin than the type A, B and C QPOs, and it may also explain why the only 18 Hz QPO (MJD 51239) reported by Sobczak et al. (2000a) did not follow the relation between QPO properties and the spectral parameters seen on the hard branches. On the other hand, there some clues that do suggest a relation with the type A, B, and C QPOs: when their values would be plotted in Figure 14 they would lie close to the extrapolation of a line through the filled symbols. Moreover, the 15.6 Hz and the 17.9 Hz QPO both fall on the empirical relation found by Wijnands and van der Klis (1999) for the low frequency QPO and break frequency found in many types of X-ray binaries, including black holes. Hence, it is not clear whether these QPO really have a different origin than the A, B, and C type QPOs, or that they only have appear to be different because the hard spectral component is so much weaker. The properties of the 16–18 Hz QPO are in fact remarkably similar to those of the 14–23 Hz QPO in GRO J1655–564, studied by Sobczak et al. (2000a); the similarity extends to the frequency at which they are found, their amplitude, and their relation with the disk spectral parameters. Note the apparent switch from a correlation of QPO frequency with hardness (15.6 Hz and 17.9 Hz QPOs) to an anticorrelation (A, B and C types) is probably not related the one reported by Rutledge et al. (1999, see above). They did not find QPOs in such spectrally soft states, and the frequencies of the 15.6/17.9 Hz QPOs are still well above that of most A, B, and C type QPOs.
The range over which the fundamental of the low frequency QPOs is observed in XTE J1550–564 is 0.1 Hz to 6 Hz (which includes the QPOs reported by Cui, Zhang and Chen (1998), but not the 16–18 Hz QPOs discussed in the previous paragraph). Using the expression for the lowest line in Figure 22, $`\nu _{HF}=38.1(\nu _{LF}+1.61)`$, and assuming that the high frequency QPO is due to orbital motion at inner disc radius ($`R_{in}`$), $`\nu _{HF}R_{in}^{2/3}`$, one can estimate the corresponding changes in the inner disc radius. For the low frequency QPO changing from 0.1 to 6 Hz we find a decrease in $`R_{in}`$ by a factor 2.6, which is comparable to what was found by di Matteo and Psaltis (1999) for other black hole systems. However, their relation between $`R_{in}`$ and $`\nu _{LF}`$ was based on the empirical relation between $`\nu _{LF}`$ and $`\nu _{HF}`$ found for neutron star sources by Psaltis, Belloni and van der Klis (1999). Using the relation of di Matteo and Psaltis (1999) we find a decrease in $`R_{in}`$ by factor $``$4.2. Both numbers suggest that the inner radius changes are rather small when a black hole changes from a hard (where the 0.1 Hz QPO was observed) to a much softer state (where the 6 Hz QPO was observed). It should be noted that these changes may in fact be somewhat larger if the 15.6/17.9 Hz QPO turns out to be related to the A, B, and C type QPOs, and/or if the lowest peak in the Cui, Zhang and Chen (1998) power spectra is not the fundamental, but a higher harmonic.
There are two types of low frequency QPOs in the neutron star Z sources (Hasinger and van der Klis, 1989) that may be compared with the low frequency QPOs in XTE J1550–564 (and other black hole candidates): these are the horizontal branch QPOs (HBO) and the normal branch QPOs (NBO) (see van der Klis (1995b) for a review). Similar QPOs have also been found in a number of neutron star atoll sources. Of the two QPO types in Z sources it is the HBO that bears most resemblance to the QPOs in XTE J1550–564. Unlike NBOs, HBOs have a strong harmonic content; e.g., in GX 340+0 the HBOs could be fitted with three harmonically related peaks (1st, 2nd and 4th) plus a shoulder component for the second harmonic (Jonker et al., 2000), similar to the type B QPOs in XTE J1550–564. HBOs are found in the 15–60 Hz range, and their frequency changes smoothly; NBOs are found in the 6–20 Hz range, but their frequency changes are strongly related to sudden spectral/state changes. Although the 6–20 Hz range of the NBOs is closer to the values we found for the low frequency QPOs in XTE J1550–564, it should be noted that if one scales the frequency of those low frequency QPOs with a factor $``$5 (which is maximum $`\nu _{upper}`$ in the neutron star sources divided by the maximum $`\nu _{HF}`$ in XTE J1550–564), one gets values that are similar to those found for the HBO. Another indication that the QPOs in XTE J1550–564 are related to the HBOs (at least to the HBOs in the Z source GX 17+2 (di Salvo et al., 2000b; Homan et al., 2001) ) is the fact that they are both observed when a hard power law tail is present in the energy spectrum, and that both show an anti-correlation between their frequency and the strength of this high energy component. Finally we note that the HBO in Z sources is accompanied by a low frequency noise component that becomes stronger when the HBO frequency decreases, which is similar to what is observed in XTE J1550–564, where a strong noise component develops when the QPO frequency drops (see also figures in Cui, et al. (1999) and Remillard et al. (1999a)).
When comparing the properties of the 0.01–0.1 Hz noise in the high state flares with those in between the flares, we find that in the flares the noise was stronger, but had a softer fractional rms spectrum, whereas the overall X-ray spectrum of the source was harder. Though this might at first appear remarkable, it is in perfect agreement with the assumption that the extra noise is associated with the hard power law component. A hard spectral component, with associated noise that remains a constant fraction of it, in combination with a soft spectral component that remains unchanged, will lead to a softer fractional spectral dependence of the noise when the hard component increases.
The transition observed on MJD 51254 showed many similarities to the “dips” and “flip-flops” observed in GX 339–4 and GS 1124–68 (Miyamoto and Kitamoto, 1991; Takizawa et al., 1997), although the time scale of the transition we observed ($``$100 s) is quite long compared to the transitions in these dips and flip-flops. Both showed a QPO in their upper count rate level (but not in their lower level), and a somewhat stronger noise in their lower count rate level. Like in GX 339–4 and GS 1124–68 the transition occurred in a region in the CD where power law noise changes to band limited noise. Also, the count rate differences ($``$10%) were accompanied by relatively subtle spectral differences, showing that on these short time scales spectral hardness and power spectral properties do not correlate as well as they do on longer time scales. The transition probably originated in the accretion disc component; the frequency of the QPOs before and after the transition were different and indicate that the inner disk radius had decreased a few percent. This change did apparently not affect the spectrum of the disk component, since the soft color remained constant. The hard color on the other hand did change, but certainly not as dramatically as the count rate. A slow decrease in the hard color started around the time of the transition maybe reflecting some kind of cooling of the hard spectral component.
## 7 Conclusions
Our main conclusions are summarized as followed:
* XTE J1550–564 was found to change between spectrally hard and soft states on time scales of days to weeks. These transition took place at a more or less constant 2–60 keV count rate level, and were found at count rate levels that differed by up to a factor 1000.
* As the spectral hardness increased, both the spectral and variability properties changed from HS via IS to LS. We regard the VHS as an instance of the IS.
* At least two physical parameters seem to be necessary to account for the behavior of XTE J1550–564. These parameters vary to a large extent independently from each other. One of these parameters is probably the mass accretion rate (through the disk). The other parameter seems to determine the state of the source and may for instance be the (relative) size of a Comptonizing region.
* The inner disc radius, as inferred from variability porperties, increases by a factor of 3–4 as the source moves from the HS to the LS.
* The properties of the QPOs (frequency, coherence and harmonic content) as well as the shape and strength of the broad band noise (weak power law or strong band limited) are well correlated with spectral hardness.
* The frequencies of the low and high frequency QPOs correlated well with each other, but in a way that is inconsistent with empirical relations found for the low and high frequency QPOs in neutron star systems.
This research has made use of data obtained through the High Energy Astrophysics Science Archive Research Center Online Service, provided by the NASA/Goddard Space Flight Center. This work was supported by NWO Spinoza grant 08-0 to E.P.J. van den Heuvel, by the Netherlands Organisation fo Scientific Research (NWO) under contract number 614-51-002, and by the Netherlands Research-school for Astronomy (NOVA). RW was supported by NASA through the Chandra Postdoctoral Fellowship grant number PF9-10010 awarded by the Chandra X-ray Center, which is operated by the Smitsonian Astrophysical Observatory for NASA under contract NAS8-39073. MM is a fellow of the Consejo Nacional de Investigaciones Científicas y Técnicas de la República Argentina. JH would like to thank Peter Jonker for many usefull discussions. |
warning/0001/hep-th0001117.html | ar5iv | text | # I. Introduction.
## I. Introduction.
In the classical theory the space–time is a continuum, where the fundamental elements are points. In continual geometry all the geometrical objects are sets of points. Scalar, vector and tensor quantities are the functions of point coordinates. Mathematical analysis operations (limits, derivations etc.) are defined at a point.
In the Dirac quantum mechanics and the quantum field theory the space–time is also represented as a continuum. In the classical theory of gravitation the space– time is the continuum. But the quantum analysis of the space–time properties provides some arguments counting in favour of the existence of the minimal length, that can be measured by physical methods. The Heizenberg uncertainty relations application to the process of small distances measuring yields the inequality
$$\left(\mathrm{\Delta }L\right)^22l_{pl}^2$$
(1)
where $`l_{pl}=\left(G^1\mathrm{}^1c^3\right)^{\frac{1}{2}}10^{33}sm`$.
The analysis of the space–time properties on the Plank distances leads to the idea of ”space–time foam” \[2-3\]. The postulate about a fundamental length ($`l_f`$) existence is realized in the concept of the quantized (discrete) space–time, that consists of the fundamental elements with the finite sizes \[4-16\]. The lattice space–time with the fixed lattice is most commonly investigated \[9-13\].
But the space–time with the fixed lattice consideration leads to several problems. The first and seemingly the most essential problem is the passage from the lattice space–time to the continuum in the limit $`l_f0`$. The lattice space–time has the power of a countable set. Any subdivision of a lattice yields a set with the same power. Thus in the limit $`l_f0`$ any lattice space–time with any subdivision remains a set with the power of a countable set. The second, the resulting equations in the lattice space and other spaces with determined fundamental element form depend on the form of a fundamental element. The third, the equations in the lattice space are non-invariant under the continual symmetry operations.
Regarding the quantum ideology all physical quantities don’t have any determined values. Thus if the quantum concept is applied to the space–time consistently, then it is possible to operate on only the average sizes of the elements. In this sense all the determined forms of the space–time fundamental element including stereohedra are not consistently quantum description of the space–time. Consistently quantized space–time can be formed by the fundamental elements without any determined forms and sizes only. Nothing but average sizes of each fundamental element in consistently quantized space–time can be determined.
Discrete geometry has been developed in the direction of the formless fundamental element in the last thirty years. The Regge calculus and the space–time foam idea \[2-3\] made the first steps on this way. In refs. the stereohedra space is investigated. This space fundamental element has some set of forms. Random lattice field theory is analyzed in . In ref. develops the topological approach to quantized space–time calculations. Quantum configuration space investigated in is the method of quantized space–time description based on the not–fixed ”floating lattice”.
In the present work the concept of the space of the formless finite fundamental elements is suggested. In this concept the problem of the passage to the continuum in the limit $`l_f0`$ and several other problems of the lattice space–time may be solved.
Notations
$`n`$– the dimensionality of a space or a space–time
$`\eta _{ik}`$– the metric tensor of the plane space or space–time
$`g_{ik}`$–a metric tensor of the Riemannian space or space–time
$`\left\{a\right\}`$– a set of coordinates in the space or the space–time of FFFEs
$`l_f`$– the fundamental length
$`l_{pl}`$– the Plank length
Abbreviations
FFFE – a formless finite fundamental element
FL – the fundamental length
Below $`l_fl_{pl}`$ is supposed.
## II. Geometry of the space of the formless finite fundamental elements.
Introduce the postulates, that differ the geometry with a formless finite fundamental element from the continual geometry.
Postulate 1. The fundamental length ($`l_f`$) is a geometrical quantity of length dimension that means the quantum limit of measurements accuracy in the space (space–time).
Postulate 2. The space (space–time) consists of the fundamental elements that have finite sizes.
Postulate 3. The formless finite fundamental elements have average sizes by order fundamental length at every dimension. All physical and geometrical quantities are described as fields defined on a set of FFFEs.
Postulate 4. All physical and geometrical quantities don’t depend on the form of concrete FFFE. They can depend on average geometrical characteristics of FFFEs only.
In the text below $`n`$ is the dimensionality of the space (or space–time) of FFFEs. This dimensionality is equal to the dimensionality of the continual space. Introduction of the space of FFFEs axiomatically allows to construct the mathematical objects and operations in this space without consideration of the continual space objects and operations.
On the one hand these postulates are the most probable to be obtained consistently from analysis of the space–time quantum properties at the small distances. On the other hand the axiomatical definition of the quantized space–time properties could itself leads to the quantization phenomenon and the quantum field theory as the consequences of the space–time structure.
The postulates 1 and 2 are identical to the postulates of lattice and stereohedra geometries. But the postulates 3 and 4 specifies the geometry with a formless finite fundamental element from other geometrical construction of the discrete space.
The postulate 4 allows to define the space of FFFEs as the set of coverings of the continual space by any number of non-overlapping simply connected regions of any form and arbitrary sizes. This set is provided with the probability measure, i.e. each covering contributes to the space with some probability. This measure enables the calculations based on this coverings set (see the sections III, IV). The average values of sizes of FFFEs are equal to $`l_f`$, and the average number of FFFEs localized in the continual space region by the volume $`V`$ is $`N=\left[V\left(l_f^n\right)^1\right]`$. But the configurations with greatly different from $`l_f`$ FFFEs sizes also have the finite probabilities, for example, the configurations with one fundamental element that expands on all the space–time (or investigated manifold), or the configuration of continual space–time region itself, i.e. covering this region by points. This set of coverings have the power of continuum. Therefore limit passage from the space of FFFEs to the continual space can be carried out correctly.
In the space of FFFEs the coordinates can be introduced in the region, consisting of a number of FFFEs. The space coordinates on one fundamental element don’t have a determined meaning in the space of FFFEs and can be of auxiliary character. The coordinates introducing on all the set of FFFEs is difficult problem due to a number of fundamental elements (regions from coverings) on one manifold is variable. The coordinates can be introduced with sufficient correctness only on the set of configurations in which all sizes of all elements are about equal to $`l_f`$, all $`m`$–dimensional areas are about equal to $`l_f^m`$, and all elements volumes are about equal to $`l_f^n`$.
All geometrical operations in the space of FFFEs are determined with accuracy $`O\left(l_f^k\right)`$. Thus generators of the rotation group in the space of FFFEs are rotations on a finite angle. Evidently the infinite small transformation like the ones in the continual space cannot be the space transformation operations defined in the space of FFFEs, because the infinite small transformation doesn’t cause any modifications in the set of FFFEs. The translation group generators are the translations on a finite distance (by order $`l_f`$).
This accuracy limit of operations determination helps to solve the same problems arising in the lattice space–time consideration. Thus the lattice Dirac equation is relativistic invariant with averaging on the continual rotation group only . In the space–time of FFFEs this problem is solved at the postulate level, because this averaging is the consequence of the postulates 1-4.
The many other constructions of the discrete space–time (i.e. the Regge coverings, the lattice space–time, the random lattice space–time, the stereohedra space–time) are the special cases of FFFE space–time with the special choice of the probability measure. Thus the lattice space–time is the set of coverings with the probability measure that is equal to zero for the configurations differ from the coverings by $`n`$–dimensional cubes with identical sizes. The random lattice space–time is the set of coverings with the probability measure that is not equal to zero for the coverings with the rectangular lattice of the variable step.
The mathematical operations and the physical equations in the space and the space–time of FFFEs could be obtained by two methods. The first one is based on the known operations and equations of the continual space. The operations and equations in the FFFE space are the ones for average values, that are calculated by the method of functional integrals. This method is considered in the section III.
The second method is the postulative introduction of operations and equations in the space of FFFEs that requires the definitions of invariant objects on the set of FFFEs. These objects depend on the place of concrete FFFE among the other elements and average sizes and volume, in the same time they don’t depend on the form and the sizes of concrete FFFE.
Invariant objects defined on each fundamental element must be the invariants of the complete space transformation group. One can note that the plane space of FFFEs has a specific transformation operation that is absent in the continual spaces and in the FFFE curved space. This operation is rearranging of elements. Regarding physics the plane space–time can be free of particles only, when with geometrical consideration the properties of the plane space are identical in all the space. Therefore any number of elements are able to change their localization in any order, and space of FFFE or the manifold of this space is transformed into itself. In the Riemannian space this operation isn’t symmetry operation due to the coordinate dependence of the connection and different values of excitations probabilities on different fundamental elements. Due to this rearranging symmetry the plane space and the Minkowski space–time of FFFEs are completely stochastized because any FFFE localization region isn’t exactly determined. Riemannian space and space–time with particle–like excitations are not stochastized since geometrical and physical properties is chosen from one element to other.
## III. Functional integral in the space–time of FFFEs.
In the previous section the calculations method with use of invariant structures of the FFFE space was discussed.
The other way of obtaining the physical equations and mathematical operations in the space and the space–time of FFFEs is calculations with use of a continual (functional) integral. In agreement with the central idea of the continual integral theory the calculation of quantum quantities is the integrating over all possible configurations of the space (space–time) of FFFEs (i.e. coverings of the space (space–time)) with into account the corresponding probability measure taken.
Consider the general construction of a functional integral. In the plane space it is:
$$Z=𝒟Ve^{S\left(s_i\right)},$$
(2)
where $`𝒟V`$ is a measure in the set of coverings, $`S\left(s_i\right)`$ is the plane space (space–time) vacuum action, $`s_i`$ is the set of element parameters (sizes, areas, volume). Here integrating is over all coverings of the continual space (space–time) by non-overlapping simply connected regions of any forms and any sizes (see below). Average value of a function on a separate FFFE is defined by
$$<f\left(\left\{a\right\}\right)>=\frac{𝒟Ve^{S\left(s_i\right)}f_{\left\{a\right\}}\left(x^i\right)}{𝒟Ve^{S\left(s_i\right)}}$$
(3)
where $`f_{\left\{a\right\}}\left(x^i\right)`$ is values of the function $`f`$ at regions of coverings set which forms the element $`\left\{a\right\}`$, $`f\left(x^i\right)`$ is a function defined in the continual space (space–time).
In the curved space (space–time) a vacuum functional integral is
$$Z=𝒟V𝒟g_{ik}e^{S(s_i,g_{ik})}$$
(4)
Here $`S`$ is the curved space (space–time) action. Full Riemannian space–time action includes particle terms (see in detail in the section VI). This action is the one of the space–time with excited states, i.e. vacuum action + action of excitations. Average value of an operator in the Riemannian space (space–time) is represented by
$$<A\left(\left\{a\right\}\right)>=\frac{𝒟V𝒟g_{ik}𝒟\phi _mA_{\left\{a\right\}}(g_{ik},\phi _m,x^i)e^{S(g_{ik},\phi _m,s_i)}}{𝒟V𝒟g_{ik}𝒟\phi _me^{S(g_{ik},\phi _m,s_i)}}$$
(5)
where $`\phi _m`$ is any fields defined in the continual space. Thus the functional integrating operation is the one of the averaging over all configurations that form the space of FFFEs with the corresponding action. This construction is similar to the functional integral over surfaces in the Polyakov superstring theory .
Integrating in the functional integral is over all possible coverings of the continual space by non-overlapping simply connected regions with arbitrary sizes and forms. The application of FFFE functional integrals require the information about the action $`S`$. This problem is discussed in the section IV for the vacuum case.
In the space of FFFEs the wave function of each FFFE could be introduced (see the remark about the coordinates introducing in FFFE space in the section II). This function squared determines the probability of finding the concrete FFFE in the state with an average localization point $`\stackrel{}{r}`$, the volume $`V`$, total $`m`$–dimensional areas $`S_m`$ and sizes $`l_i`$. Denote it $`\psi _{\left\{a\right\}}(\stackrel{}{r},V,S_m,l_i)`$. Here $`\left\{a\right\}`$ is the set of fundamental element coordinates in the space of FFFEs. This wave function squared $`\left|\psi \right|^2`$ is a density of probability in the set of coverings, i.e. $`\left|\psi \right|^2d\sigma `$ is the probability of finding the element with FFFE space coordinates $`\left\{a\right\}`$ in the state with continual parameters $`l_i,S_m,V`$. Here $`d\sigma `$ is a measure in the set of coverings. In principle any physical quantities could be found as matrix elements
$$<A\left(\left\{a\right\}\right)>=<\psi _{\left\{a\right\}}(\stackrel{}{r},V,S_m,l_i)\left|\widehat{A}\right|\psi _{\left\{a\right\}}(\stackrel{}{r},V,S_m,l_i)>$$
(6)
where $`A`$ is a function in the continual space. The notation $`<A\left(\left\{a\right\}\right)>`$ means the summing over all regions of coverings set which form the element with the coordinates set $`\left\{a\right\}`$ in the space (or space–time) of FFFEs. Summarized quantities are average values of $`<A>`$ on the each configuration element, multiplied on $`\left|\psi \right|^2`$.
The state of the space of FFFEs is a covering of the continual space. It is also possible to introduce the wave function of a space manifold or all space wave function (see the discussion of the corrections below):
$$\mathrm{\Psi }=\underset{\left\{a\right\}}{}\psi _{\left\{a\right\}}(\stackrel{}{r},V,S_m,l_i)+com\left(\left\{\psi _{\left\{a\right\}}\right\}\right)$$
(7)
This $`\mathrm{\Psi }`$ describes a covering of the continual space by non-overlapping simply connected regions. $`\left|\mathrm{\Psi }\right|^2`$ is the probability density in the set of coverings. This function describes the state of the space–time without particles, i.e. particle–like excited states of FFFE. However, the vacuum itself has excited states, where the elements sizes and localization points differ greatly from the average values. Functions $`\psi _{\left\{a\right\}}`$ are not independent because they describe non-overlapping regions. Therefore the expression (7) contains the second term that is determined by the commutation relations.
## IV. The Minkowski space–time action.
The functional integral construction considered in the previous section requires the information about the action. In this section the case of the pure vacuum is discussed. Regarding physics the case of the vacuum is described by the plane Minkowski space–time. Any nonvacuum excitation, including the virtual particle vacuum polarization and the all space constant fields leads to arising of the connection and the curved structure of the space–time.
Suppose that one term of the space–time action is proportional to the volume of FFFE:
$$S_V=\underset{FFFE}{}\alpha \sqrt{g}𝑑V$$
(8)
Here $`dV`$ and $`\sqrt{g}`$ are continual space values. Integrating in (8) is over one region from some covering of the continual space–time. This term is analogous of ”space–time foam” action proportional to the volume . But this term of an action is unsufficient for describing the equilibrium configuration of the space–time of FFFE. Total actions (8) of all configurations from the set of coverings are equal.
Consider the problem of the space–time action minimum. On the one hand as a rule in the method of functional integrating the minimum of considered action describes the corresponding classical system (a moving pointlike particle in the Feynman integral, the space–time with the classical value of a metric tensor in the integrating over space–time metrics in quantum gravity etc.). But in the case of the space–time the classical system is the continual space–time. It means that the space–time action must have the minimum in the configuration with infinite number of fundamental elements which are points. However, the action minimum in the continual space–time configuration, and as the consequence the finite value of the probability measure of this configuration, leads to divergence of some integrals, for example, the average value of a number of fundamental elements on a continual manifold. It is seems that the probability measure must be small in the configurations with small number of fundamental elements.
On the other hand the equilibrium state of the FFFE space–time is the one with average FFFE sizes at a fundamental length. The more correct approach to the equilibrium action problem is to find the fundamental length, firstly, as the average value of a FFFE size, secondly, as realizing of some quantized action minimum. It means that the average number of FFFEs in a manifold is determined by the solution of the action minimum problem.
The complete expression for the space–time action, meeting this requirement, must contain other terms besides the volume term. The possible term is the one proportional to the total $`n1`$– dimensional area of FFFE. In this supposition the vacuum space–time action is
$$S=\alpha \underset{i}{}V_i+\beta \underset{i}{}\left(S_{n1}\right)_i$$
(9)
where $`_i`$ is summarizing over all FFFEs.
The constants in the action expression can be product of the universal constants only. Thus in the four–dimensional space–time
$$\alpha =A\mathrm{}^1G^2c^6$$
(10)
where $`A`$ is a numerical factor.
Let us direct the way, on which the Nambu–Goto term of the string action might be obtained from the space–time action (8). The expression of the action of a space–time element for this analysis is required. This action is the average value of an action $`S`$ with the functional integrating (3) using. This action is denoted by $`S_{FFFE}`$:
$$S_{FFFE}=<S_{\left\{a\right\}}>$$
(11)
or
$$S_{FFFE}=\sqrt{\eta }𝑑x^1𝑑x^2𝑑x^3𝑑x^4$$
(12)
for the four–dimensional space–time. This construction also could be introduced axiomatically as the invariant structure (see the section II).
A string in the space–time of FFFEs can be considered as an excitation of a number of FFFEs forming one–dimensional space–like curve (in the meaning of FFFE space–time). In the own reference frame the action of this excitation is represented in the form
$$S=A\mathrm{}^1G^2c^6\underset{FFFEs}{}\sqrt{\eta }𝑑x^1𝑑x^2𝑑l𝑑\tau $$
(13)
where $`\tau `$ is the own time, $`l`$ is the own space–like coordinate of an excitation, $`x^1,x^2`$ are the transverse space–like coordinates. Here integrating is over a set of FFFEs, participated in the excitation propagation. After integration over transverse coordinates we might obtain ($`\gamma `$ is the two–dimensional metric tensor determinant):
$$S=AG^1c^3𝑑l𝑑\tau \sqrt{\gamma }$$
(14)
i.e. the Nambu–Goto action for a string. In (14) the equality of the FFFE average size on the one dimension to the fundamental length is taken into account. This result is not completely correct due to the transformation problem of the four–dimensional metric tensor determinant $`\eta `$ in (13) to the two–dimensional one $`\gamma `$ in (14) and absence of the correct definition of one–dimensional integrating. In this concept the $`p`$–branes are considered as the $`p`$–dimensional space–like excitations of FFFEs, and the volume term (13) of FFFE space–time excitation yields the bosonic term of $`p`$–branes action analogically.
## V. The compactification in the space–time of FFFEs.
In the concept of the FFFE space–time the multidimensional space–time with the motion possible on four dimensions only can be described without any special compactification procedure. The multidimensional space–time with average sizes on the higher dimensions by order $`l_f`$ can be constructed using the postulates about the absence of nighbour elements for all ones on all the dimensions without four. These average sizes on the higher dimensions are $`l_f`$ despite the configuration with sizes, which are significantly more $`l_f`$ on these dimensions contributes to the space–time structure.
But this compactification description is not satisfying as it requires the introducing the special postulate. The deeper approach to the compactification problem is to formulate the neighbour element absence requirement, caused space–time action structure analysis or some geometrical requirement. Erenfest’s investigations about stability of systems with Coulomb interaction shows that 4–dimensionality of the real space–time connect with particle action, more precisely, with the interaction part of particles action. This way the case $`n=4`$ of the FFFE space–time dimensions, on which the motion is possible, yields the minimum of an action of $`n`$–dimensional curved space–time.
## VI. Physical fields in the space–time of FFFEs.
Fermionic and bosonic fields are excited states of FFFEs. As any quantum particle, excluding a free particle, has a wave function with different values $`\left|\psi \right|^2`$ in different points of the space–time, the space–time with excitations couldn’t be the Minkowski space–time. Different values of $`\left|\psi \right|^2`$ in different FFFEs violate the Minkowski space–time specific symmetry under rearranging of any number of FFFEs (see section III). Therefore two interacting particles in the space–time result in the curved space–time with changeable curvature. One particle in all space–time or the uniform vacuum polarization leads to the particle–like excited curved space–time with the constant curvature.
Particle–like excitations of the FFFE space–time are finite in each FFFE of considered manifold. The state of FFFE with particles excitations in the approximation of ininteracting particles is described by a wave function
$$\psi _{\left\{a\right\}}^{ex}=|\left\{a\right\},\psi _i>,$$
(15)
where $`\left\{a\right\}`$ is a set of coordinates in the FFFE space–time, $`\psi _i`$ is one-particle wave functions values in this FFFE.
A wave function $`|\left\{a\right\},0>`$ is the sum of all excited vacuum states. In the classical space–time $`\left|\psi \right|^2dV`$ is interpreted as a probability of a particle localization in the volume $`dV`$. In the FFFE space–time the interpretation $`\left|\psi _{\left\{n\right\}}\right|^2dV_{fund}`$ is a probability of finding this FFFE in the excited state with the set of quantum numbers (charges) $`\left\{n\right\}`$. This probabilities equality means the influence of the other FFFEs with excited states on the excited state of this FFFE. In other words, the space–time of FFFEs is the self–organizing system.
The states of a manifold of the space–time of FFFEs is described by an overall wave function, that could be obtained by summarizing of each FFFE wave function with the commutation relations taken into account. Functional integral construction in the curved space–time requires to include particles terms of the space–time action in consideration.
It is to suppose that the number of possible fermionic excitation in one FFFE is finite. In this case the space–time supersymmetry leads to the supposition about the equality of a number of possible fermionic excitations in one FFFE and a number of FFFEs in all the space–time.
## VII. Final remarks.
Ideologically the suggested concept is consistent geometrical approach to the physics of fundamental particles and interactions. This concept may help to solve some problems of lattice space–time geometry due to more consistent quantum approach to the space–time structure problem. In this space–time the particles are geometrical objects - excited states of space–time elements. The superstrings can be constructed as propagating excitations of space–time elements. At this approach the Nambu–Goto action term is considered as a result of this excitation volume space–time action term analysis. Consideration of the superstrings as the excitations of the quantized space–time is the step to the understanding of the superstrings properties at the Plank distances. With this superstrings and $`p`$–branes consideration all these objects are identical at the Plank distances because the excitation of one element doesn’t have dimension in the sense of FFFE space–time.
But some problems of elementary particle and field theories are not obvious in this concept, i.e. the appearance and the role of gauge invariance in the FFFE space–time, appearance of charges of Riemannian space–time excited states, positing of the cosmological evolution problems and some others.
In conclusion it is some words about axiomatical introducing of the quantized space–time (see the section II). Certainly the author can’t be sure that the postulates 1 - 4 are most correct, complete and minimal system of the quantized space–time axioms. It is not improbable that the postulate 1 about the fundamental length existence is the consequence of some other axioms system, and at the same time the quantized space–time properties are defined by the axioms introducing the set of FFFEs operationally. But it is to be noted that on this way the uncertainty relations, quantum field theory and quantization phenomenon itself are the consequences of these quantized space–time axioms system. In particular, it is supposed that the Dirac term of lagrangian density can come as a consequence of coordinates matrix introducing and the Weyl structure of FFFE curved space–time. The Weyl structure of FFFE curved space–time is connected with the finite accuracy of any geometrical operations. Thus the vector length at its parallel transport from one FFFE to other is determined with an accuracy of $`l_f`$. Therefore the FFFE curved space–time is the Weyl space–time automatically, and the Weyl distortion of the Riemannian structure of this space–time is caused by quantized structure of the space–time principally.
The author thanks M.S. Orlov, A.V. Klochkov, E.V. Klochkova, A.V. Sokolov, G.S. Sokolova, D.V. Sokolov, A.B. Vankov for their friendly support during the work time, A.A. Amerikantsev, M.E. Golod for their technical assistance and M.A. Tyntarev for useful discussions. |
warning/0001/physics0001071.html | ar5iv | text | # Untitled Document
Non equilibrium in statistical and fluid mechanics.
Ensembles and their equivalence. Entropy driven intermittency.
Giovanni Gallavotti
Fisica, Università di Roma 1
P.le Moro 2, 00185 Roma, Italia
Abstract:We present a review of the chaotic hypothesis and discuss its applications to intermittency in statistical mechanics and fluid mechanics proposing a quantitative definition. Entropy creation rate is interpreted in terms of certain intermittency phenomena. An attempt to a theory of the experiment of Ciliberto–Laroche on the fluctuation law is presented.
§1. Introduction.
A general theory of non equilibrium stationary phenomena extending classical thermodynamics to stationary non equilibria is, perhaps surprisingly, still a major open problem more than a century past the work of Boltzmann (and Maxwell, Gibbs,…) which made the breakthrough towards an understanding of properties of matter based on microscopic Newton’s equations and the atomic model.
In the last thirty years, or so, some progress appears to have been achieved since the recognition that non equilibrium statistical mechanics and stationary turbulence in fluids are closely related problems and, in a sense, in spite of the apparently very different nature of the equations describing them they are essentially the same.
The unifying principle, originally proposed for turbulent motions by Ruelle, \[Ru78\], in the early 970’s, has been extended to statistical mechanics and eventually called the “chaotic hypothesis”, \[GC95\]:
Chaotic hypothesis: Asymptotic motions of a chaotic system, be it a multi particle system of microscopic particles or a turbulent macroscopic fluid, can be regarded as a transitive Anosov system for the purposes of computing time averages in stationary states.
It may be useful to make a few comments on how this is supposed to be interpreted. The conclusions that we draw here from the chaotic hypothesis are summarized in §13 which might be consulted at this point. For a review on the subject seen from a different perspective see \[Ru99a\]
§2. Meaning of the chaotic hypothesis.
Anosov systems are well understood dynamical systems: they play a paradigmatic role with respect to chaotic systems parallel to the one harmonic oscillators play with respect to orderly motions. They are so simple, and yet very chaotic, that their properties are likely to be the ones everybody develops in thinking about chaos, even without having any familiarity with Anosov systems which certainly are not (yet) part of the background of most contemporary physicists.<sup>1</sup>
<sup>1</sup>Informally a map $`xSx`$ is a Anosov map if at every point $`x`$ of the bounded phase space $`M`$ one can set up a local system of coordinates with origin at $`x`$, continuously dependent on $`x`$ and covariant under the action of $`S`$ and such that in this comoving system of coordinates the point $`x`$ appears as a hyperbolic fixed point for $`S`$. The corresponding continuous time motion, when the evolution is $`xS_tx,tR`$, requires that the local system of coordinates contains the phase space velocity $`\dot{x}`$ as one of the coordinate axes and that the motion transversal to it sees $`x`$ as a hyperbolic fixed point: note that a motion in continuous time cannot possibly be hyperbolic in all directions and it has to be neutral in the direction of $`\dot{x}`$ because the velocity has to be bounded if $`M`$ is bounded, while hyperbolicity would imply exponential growth as either $`t+\mathrm{}`$ or $`t\mathrm{}`$. Furthermore there should be no equilibrium points and the periodic points should be dense in phase space. When the system has one or more (the so called “hysteresis phenomenon’’) attracting sets which do not occupy the whole phase space the chaotic hypothesis can be interpreted as saying that each attracting set is a smooth surface on which the time evolution flow (or map) acts as an Anosov flow (map).
In general an Anosov system has asymptotic motions which approach one out of finitely many invariant closed sets $`C_1,\mathrm{},C_q`$ each of which contains a dense orbit,:one says that the systems $`(C_j,S_t)`$ are “transitive”. One of them, at least, must be an attractive set.
To say that “the asymptotic motions form a transitive Anosov system” means that (1) each of the sets $`C_j`$ which is attractive is a smooth surface in phase space and (2) only one of them is attractive:
The last “transitivity” assumption is meant to exclude the trivial case in which there are more than one attractive sets and the system de facto consists of several independent systems.
The smoothness of $`C_j`$ is a strong assumption that means that one does not regard possible lack of smoothness, i.e. fractality, as a really relevant property in systems with large number of degrees of freedom. In any event one could consider (if necessary) replacing “Anosov systems” with some slightly weaker property like “axiom A” systems which could permit more general asymptotic motions. Here we adhere strictly to the chaotic hypothesis in the stated original form, \[GC95\].
§3. Basic implications of the chaotic hypothesis and relation with the ergodic hypothesis.
The chaotic hypothesis boldly extends to non equilibrium the ergodic hypothesis: applied to equilibrium systems, i.e. to systems described by Hamiltonian equations, it implies the latter, \[Ga98\]. This means that if a Hamiltonian system at a given energy is assumed to verify the chaotic hypothesis, i.e. to be a transitive Anosov system, then for all observables $`F`$ (i.e. for all smooth functions $`F`$ defined on phase space)
$$T^1_0^TF(S_tx)𝑑t\begin{array}{c}T\mathrm{}\hfill \end{array}_MF(y)\mu _L(dy)$$
$`(3.1)3.1`$
where $`\mu _L`$ is the Liouville distribution on the constant energy surface $`M`$, and (3.1) holds for almost all points $`xM`$, i.e. for $`x`$ outside a set $`𝒩`$ of zero Liouville volume on $`M`$.
Being very general one cannot expect that the chaotic hypothesis will solve any special problem typical of non equilibrium physics, like “proving” the Fourier’s law of heat conduction, the Ohm’s law of electric conduction or the K41 theory of homogeneous turbulence.
Nevertheless, like the ergodic hypothesis in equilibrium, the chaotic hypothesis accomplishes the remarkable task of giving us the “statistics” of motions. If $`M`$ is the phase space, which we suppose a smooth bounded surface, and $`tS_tx`$ is the motion starting at $`xM`$, the time average:
$$T^1_0^TF(S_tx)𝑑t\begin{array}{c}T\mathrm{}\hfill \end{array}_MF(y)\mu _{SRB}(dy)$$
$`(3.2)3.2`$
of the observable $`F`$ exists for $`x`$ outside a set $`𝒩`$ of zero phase space volume and it is $`x`$–independent, thus defining the probability distribution $`\mu _{SRB}`$ via (3.2).
Note, in fact, that the probability distribution $`\mu _{SRB}`$ defined by the l.h.s. of (3.2) is uniquely determined (provided it exists): it is usually called the “statistics of the motion” or the “SRB distribution” associated with the dynamics of the system.
To appreciate the above property (existence and uniqueness of the statistics) the following considerations seem appropriate.
An essential feature, and the main novelty, with respect to equilibrium systems is that non conservative forces may act on the system: this is in fact the very definition of “non equilibrium system”.
Since non conservative forces perform work it is necessary that on the system act also other forces that take energy out of it, at least if we wish that the system reaches a stationary state, showing a well defined statistics.
As a consequence any model of the system must contain, besides non conservative forces which keep it out of equilibrium by establishing “flows” on it (like a heat flow, a matter flow, …), also dissipative forces preventing the energy to increase indefinitely and forcing the motion to visit only a finite region of phase space.
The dissipation forces, also called “thermostatting forces”, will in general be such that the volume in phase space is no longer invariant under time evolution. Mathematically this means that the divergence $`\sigma (x)`$ of the equations of motion will be not zero and its time average $`_M\sigma (y)\mu _{SRB}(dy)\stackrel{def}{=}\sigma _+`$ will be positive or zero as it cannot be negative (“because phase space is supposed bounded”: see \[Ru96\]).
One calls a system “dissipative” if $`\sigma _+>0`$ and we expect this to be the case as soon as there are non conservative forces acting on it.
We see that if a system is dissipative then its statistics $`\mu _{SRB}`$ must be concentrated on a set of zero volume in $`M`$: this means that $`\mu _{SRB}`$ cannot be very simple, and in fact it is somewhat hard to imagine it.
If the acting forces depend on a parameter $`E`$, “strength of the non conservative forces”, and for $`E=0`$ the system is Hamiltonian we have a rather unexpected situation. At $`E=0`$ the chaotic hypothesis and the weaker ergodic hypothesis imply that the statistics $`\mu _{SRB}`$ is equal to the Liouville distribution $`\mu _L`$; but if $`E0`$, no matter how small, itwill not be possible to express $`\mu _{SRB}`$ via some density $`\rho _E(y)`$ in the form $`\mu _{SRB}(dy)=\rho _E(y)\mu _L(dy)`$, because $`\mu _{SRB}`$ attributes probability $`1`$ to a set $`𝒩`$ with zero volume in phase space (i.e. $`\mu _L(𝒩)=0`$). Nevertheless natura non facit saltus (no discontinuities appear in natural phenomena) so that sets that have probability $`1`$ with respect to $`\mu _{SRB}`$ may be all still dense in phase space, at least for $`E`$ small. In fact this is a “structural stability” property for systems which verify the chaotic hypothesis (see \[Ga96c\])
The above observations show one of the main difficulties of non equilibrium physics: the unknown $`\mu _{SRB}`$ is intrinsically more complex than a function $`\rho _E(y)`$ and we cannot hope to proceed in the familiar way we might have perhaps expected from previous experiences: namely to just set up some differential equations for the unknown $`\rho _E(y)`$.
Hence it is important that the chaotic hypothesis not only guarantees us the existence of the statistics $`\mu _{SRB}`$ but also that it does so in a “constructive way” giving at the same time formal expressions for the distribution $`\mu _{SRB}`$ which should possibly play the same role as the familiar formal expressions used in equilibrium statistical mechanics in writing expectations of observables with respect to the microcanonical distribution $`\mu _L`$.
For completeness we write a popular expression for $`\mu _{SRB}`$. If $`\gamma `$ is a periodic orbit in phase space, $`x_\gamma `$ a point on $`\gamma `$, $`T(\gamma )`$ the period of $`\gamma `$ then
$$F(y)\mu _{SRB}(dy)=\underset{T\mathrm{}}{lim}\frac{\underset{\gamma :T(\gamma )T}{}e^{_0^{T(\gamma )}\sigma (S_tx_\gamma )𝑑t}_0^{T(\gamma )}F(S_tx_\gamma )𝑑t}{_{\gamma :T(\gamma )T}e^{_0^{T(\gamma )}\sigma (S_tx_\gamma )𝑑t}T(\gamma )}$$
$`(3.3)3.3`$
This is simple in the sense that it does not require, to be formulated, an even slight understanding of any of the properties of Anosov or hyperbolic dynamical systems. But in many respects it is not a natural formula: as one can grasp from the fact that it is far from clear that in the equilibrium cases (3.3) is an alternative definition of the microcanonical ensemble (i.e. of the Liouville distribution $`\mu _L`$), in spite of the fact that in this case $`\sigma \mathrm{\hspace{0.33em}0}`$ and (3.3) becomes slightly simpler.
To prove (3.3) one first derives alternative and much more useful expressions for $`\mu _{SRB}`$ which, however, require a longer discussion to be formulated, see \[Ga99a\], \[Ga86c\]: the original work is due to Sinai and in cases more general than Anosov systems, to Ruelle and Bowen. §4. What can one expect from the chaotic hypothesis?
In equilibrium statistical mechanics we know the statistics of the motions, if the ergodic hypothesis is taken for granted. However this hardly solves the problems of equilibrium physics simply because evaluating the averages is a difficult task which is also model dependent. Nevertheless there are a few general consequences that can be drawn from the ergodic hypothesis: the simplest (and first) is embodied in the “heat theorem” of Boltzmann.
Imagine a system of $`N`$ particles in a box of volume $`V`$ subject to pair interactions and to external forces with potential energy $`W_V`$, due to the walls and providing the confinement of the particles to the box. Define
$$\begin{array}{cc}\hfill T=& \text{average kinetic energy}\hfill \\ \hfill U=& \text{total energy}\hfill \\ \hfill p=& \text{average of }_VW_V\hfill \end{array}$$
$`(4.1)4.1`$
where the averages are taken with respect to the Liouville distribution on the surface of energy $`U`$.
Imagine varying the parameters on which the system depends (e.g. the energy $`U`$ and the volume $`V`$) so that $`dU,dV`$ are the corresponding variations of $`U,V`$, then
$$(dU+pdV)/T=\mathrm{exact}$$
$`(4.2)4.2`$
expresses the heat theorem of Boltzmann.
It is a consequence of the ergodicity assumption, but it is not equivalent to it as it only involves a relation between a few averages ($`U,p,V,T`$), see \[Bo66\], \[Bo84\], \[Ga99a\]. Not only it gives us a relation which is a very familiar property of macroscopic systems, but it also suggests us that even if the ergodic hypothesis is not strictly valid some of its consequences might, still, be regarded as correct.
The proposal is to regard the chaotic hypothesis in the same way: it is possible to imagine that mathematically speaking the hypothesis is not strictly valid and that, nevertheless, it yields results which are physically correct for the few macroscopic observables in which one is really interested in.
The ergodic hypothesis implies the heat theorem as a general (“somewhat trivial”) mechanical identity valid for systems of $`N`$ particles with $`N=1,2,\mathrm{},10^{23},\mathrm{}`$. For small $`N`$ it might perhaps be regarded as a curiosity: such it must have been considered by most readers of the key paper \[Bo84\] who were possibly misled by several examples with $`N=1`$ given by Boltzmann in this and other previous papers. Like the example of the system consisting of one “averaged” Saturn ring, i.e. one homogeneous ring of mass rotating around Saturn with energy $`U`$, kinetic energy $`T`$ and “volume” $`V`$ (improbably identified with the strength of the gravitational attraction!). But for $`N=10^{23}`$ it is no longer a curiosity and it is a fundamental law of thermodynamics in equilibrium: which, therefore, can be regarded on the same footing of a symmetry being a direct consequence of the structure of the equations of motion, \[Ga99a\] appendices to Ch.1 and Ch.9. It reflects in macroscopic terms a simple microscopic assumption (i.e. Newton’s equations for atomic motions, in this case).
No new consequences of even remotely comparable importance are known to follow from the chaotic hypothesis besides the fact that it implies the validity of the ergodic hypothesis itself (hence of all its consequences, first of them classical equilibrium statistical mechanics).
Nevertheless the chaotic hypothesis does have some rather general consequences. We mention here the fluctuation theorem. Let $`\sigma (x)`$ be the phase space contraction rate and $`\sigma _+`$ be its SRB average (i.e. $`\sigma _+=\sigma (x)\mu _{SRB}(dx)`$), let $`\tau >0`$ and define
$$p(x)=\tau ^1_{\tau /2}^{\tau /2}\frac{\sigma (S_tx)}{\sigma _+}𝑑x$$
$`(4.3)4.3`$
and study the fluctuations of the observable $`p(x)`$ in the stationary state $`\mu _{SRB}`$. We write $`\pi _\tau (p)dp`$ the probability that, in the distribution $`\mu _{SRB}`$, the quantity $`p(x)`$ has actually value between $`p`$ and $`p+dp`$ as
$$\pi _\tau (p)dp=\mathrm{const}e^{\zeta _\tau (p)\tau }dp$$
$`(4.4)4.4`$
Them $`lim_\tau \mathrm{}\zeta _\tau (p)=\zeta (p)`$ exists and is convex in $`p`$; and
Theorem: (fluctuation theorem) Assume the chaotic hypothesis and suppose that the dynamics is reversible, i.e. that there is an isometry $`I`$ of phase space such that
$$IS_t=S_tI,I^2=1$$
$`(4.5)4.5`$
and that the attracting set is the full phase space.<sup>2</sup>
<sup>2</sup>It is perhaps important to stress that we distinguish between attracting set and attractor: the first is a closed set such that the motions that start close enough to it approach it ever closer; an attractor is a subset of an attracting set that (1) has probability $`1`$ with respect to the statistics $`\mu `$ of the motions that are attracted by the attracting set (a notion which makes sense when such statistics exists, but for a zero volume set of initial data, and is unique) and that (2) has the smallest Hausdorff dimension among such probability $`1`$ sets. Hence density of an attracting set in phase space does not mean that the corresponding attractor has dimension equal to that of the phase space: it could be substantially lower, see \[GC95\]. Then
$$\zeta (p)=\zeta (p)\sigma _+p,\mathrm{for}\mathrm{all}p$$
$`(4.6)4.6`$
where $`\sigma _+=\mu _{SRB}(\sigma )`$.
It should be pointed out that the above relation was first discovered in an experiment, see \[ECM93\], where also some theoretical ideas were presented, correctly linking the result to the SRB distributions theory and to time reversal symmetry. Although such hints were not followed by what can be considered a proof, \[CG99\], still the discovery has plaid a major role and greatly stimulated further research.
The interest of (4.6) is that, in general, it is a relation without free parameters. The above theorem, proved in \[GC95\] for discrete evolutions (maps) and in \[Ge97\] for continuous time systems (flows), is one among the few general consequences of the chaotic hypothesis, see \[Ga96a\], \[Ga96b\], \[Ga99b\] for others.
§5. Non equilibrium ensembles. Thermodynamic limits. Equivalence.
The chaotic hypothesis gives us, unambiguously, the probability distribution $`\mu _{SRB}`$ which has to be employed to compute averages of observables in stationary states.
For each value of the parameters on which the system depends we have, therefore a well defined probability distribution $`\mu _{SRB}`$. Calling $`\underset{¯}{\alpha }=(\alpha _1,\mathrm{},\alpha _p)`$ the parameters and $`\mu _{\underset{¯}{\alpha }}`$ the corresponding SRB distribution we consider the collection $``$ of probability distributions $`\mu _{\underset{¯}{\alpha }}`$ obtained by letting the parameters $`\underset{¯}{\alpha }`$ vary. We call such a collection an “ensemble”.
For instance $`\underset{¯}{\alpha }`$ could be the average energy $`U`$ of the system, the average kinetic energy $`T`$, the volume $`V`$, the intensity $`E`$ of the acting non conservative forces, etc
Non equilibrium thermodynamics can be defined as the set of relations that the variations of the parameters $`\underset{¯}{\alpha }`$ and of other average quantities are constrained to obey as some of them are varied. In equilibrium the heat theorem is an example of such relations. In reversible non equilibria the fluctuation theorem (4.6) is an example.
In non equilibrium systems the equations of motion play a much more prominent role than in equilibrium: in fact one of the main properties of equilibrium statistical mechanics is that dynamics enters only marginally in the definition of the statistical distributions of the equilibrium states.
The necessity of a reversibility assumption in the fluctuation theorem already hints at the usefulness of considering the equations of motion themselves as “parameters” for the ensembles describing non equilibrium stationary states: we are used to irreversible equations in describing non equilibrium phenomena (like the heat equation, the Navier Stokes equation, etc) and unless we are able to connect our experiments with reversible dynamical models we shall be unable to make use of the fluctuation theorem.
Furthermore it is quite clear that once a system is not in equilibrium and thermostatting forces act on it, the exact nature of such forces might be irrelevant within large equivalence classes: i.e. it might be irrelevant which particular “cooling device” we use to take heat out of the system. Hence one would like to have a frame into which to set up a more precise analysis of such arbitrariness. Therefore we shall set
Definition 1: A stationary ensemble $``$ for a system of particles or for a fluid is the collection of $`SRB`$ distributions, for given equations of motion, obtained by varying the parameters entering into the equations.
It can happen that for the same system one can imagine different models. In this case we would like that the models give the same results, i.e. the same averages to the same observables, at least in some relevant limit. Like in the limit of infinite size in which the number $`N`$, the volume $`V`$ and the energy $`U`$ tend to infinity but $`N/V`$ and $`U/V`$ stay constant. Or in the limit in which the Reynolds number $`R`$ tends to infinity in the case of fluids.
This gives the possibility of giving a precise meaning to the equivalence of different thermostatting mechanisms. We shall declare
Definition 2, (equivalence of ensembles): Two thermostatting mechanisms are equivalent “in the thermodynamic limit” if one can establish a one to one correspondence between the elements of the ensembles $``$ and $`^{}`$ of SRB distributions associated with the two models in such a way that the same observables, in a certain class $``$ of observables, have the same averages in corresponding distributions, at least when some of the parameters of the system are sent to suitable limiting values to which we assign the generic name of “thermodynamic limit”.
In the following sections we illustrate possible applications of this concept.
§6. Drude–Lorentz’ electric conduction models.
Understanding of electric conduction is in a very unsatisfactory state. It is usually based on linear response theory and very seldom a fundamental approach is attempted. Of course this is so for a good reason, because a fundamental approach would require imposing an electric field $`E`$ on the system and, at the same time, a thermostatting force to keep the system from blowing up and to let it approach a steady state with a current $`J_E`$ flowing in it, and then taking the ratio $`J_E/E`$ (with or without taking also the limit as $`E0`$).
However, as repeatedly mentioned, it is an open problem to study steady states out of equilibrium. Hence most theories have recourse to linear response where the problem of studying stationary non equilibria does not even arise.
The reason why this is unsatisfactory is that as long as we are in principle unable to study stationary non equilibria we are also in principle unable to estimate the size of the approximation and errors of linear response.
In spite of many attempts the old theory of Drude, see \[Be59\], \[Se87\], seems to be among the few conduction theories which try to establish a conductivity theory based on the study of electric current at non zero fields.
We imagine a set of obstacles distributed randomly or periodically and among them conduction electrons move, roughly with density of one per obstacle.
The (screened) interactions between the electrons are, at a first approximation, ignored. The collisions between electrons and obstacles (“nuclei”) will take place in the average after the electrons have traveled a distance $`\lambda =(\rho a^2)^1`$ if $`\rho `$ is the nuclei density and $`a`$ is their radius.
Between collisions the electrons, with electric charge $`e`$, accelerate in the direction of an imposed field $`\underset{¯}{E}`$ incrementing, in that direction, velocity by
$$\delta v=\frac{eE\lambda }{mv}=\frac{eE(\rho a^2)^1}{m\sqrt{k_BT/m}}$$
$`(6.1)6.1`$
where $`k_B`$ is Boltzmann’s constant. At collision they are “thermalized”: an event that is modeled by giving them a new velocity of size $`v=\sqrt{k_BT/m}`$ and a random direction.
The latter is the “thermostatting mechanism” which is a, somewhat rough, description of the energy transfer from electrons to lattice which physically corresponds to electrons losing energy in favor of lattice phonons, which in turn are kept at constant temperature by some other thermostatting mechanism which prevents the wire melting. All things considered the total current that flows will be
$$J_E=\frac{e^2}{\rho a^2\sqrt{mk_BT}}E\stackrel{def}{=}\chi E$$
$`(6.2)6.2`$
obtaining Ohm’s law.
To the same conclusion we arrive by a different thermostat model. We imagine that the electrons move exchanging energy with lattice phonons but keeping their total energy constant and equal to $`Nk_BT`$: i.e. $`2^1_{j=1}^Nm\underset{¯}{\overset{\dot{}}{x}}_j^2=3Nk_BT/2`$, where $`k_B`$ is Boltzmann’s constant. There are several forces that can achieve this result
we select the “Gaussian minimal constraint” force.<sup>3</sup>
<sup>3</sup>Not because it plays any fundamental role but because it has been studied by many authors and because it represents a mechanism very close to that proposed by Drude.We recall, for copleteness, that the effort of a constraint reaction on a motion on which the active force is $`\underset{¯}{f}`$ (with $`3N`$ components) and $`\underset{¯}{a}`$ is the acceleration of the particles (with $`3N`$ components) and $`m`$ is the mass is $`\left(\underset{¯}{a}\right)=\left(\underset{¯}{f}m\underset{¯}{a}\right)^2/m`$; then Gauss’ principle is that the effort is minimal if $`\underset{¯}{a}`$ is given the actual value of the acceleration, at fixed space positions and velocities. This is the force that is required to keep $`m\underset{¯}{\overset{\dot{}}{x}}_j^2`$ strictly constant and that is determined by “Gauss least effort” principle, see \[Ga99a\], ch. 9, appendix 4, for instance: as is well known this is, on the $`i`$–th particle, a force
$$\alpha \underset{¯}{\overset{\dot{}}{x}}_i\stackrel{def}{=}\frac{e\underset{¯}{E}\underset{j}{}\underset{¯}{\overset{\dot{}}{x}}_j}{_j\underset{¯}{\overset{\dot{}}{x}}_j^2}\underset{¯}{\overset{\dot{}}{x}}_i\frac{m\underset{¯}{E}N\underset{¯}{J}}{3Nk_BT}\underset{¯}{\overset{\dot{}}{x}}_i$$
$`(6.3)6.3`$
If there are $`N`$ particles and $`N`$ is large it follows that $`\underset{¯}{J}=N^1e_j\underset{¯}{\overset{\dot{}}{x}}_j`$ is essentially constant, see \[Ru99b\], and each particle evolves, almost independently of the others, according to an equation:
$$m\underset{¯}{\overset{¨}{x}}_i=e\underset{¯}{E}\nu \underset{¯}{\overset{\dot{}}{x}}_i$$
$`(6.4)6.4`$
between collisions, with a suitably fixed constant $`\nu `$. If we imagine that the velocity of the particles between collisions changes only by a small quantity compared to the average velocity the “friction term” which in the average will be of order $`E^2`$ will be negligible except for the fact that its “only” effect will be of insuring that the total kinetic energy stays constant and the speeds of the particles are constantly renormalized. In other words this is the same as having continuously collisions between electrons and phonons even when there is no collision between electrons and obstacles. Hence the resulting current is the same (if $`N`$ is large) as in (6.2).
§7. Ensemble equivalence: the example of electric conduction theories.
We have derived three models for the conduction problem, namely
(1) the classical model of Drude, \[Se87\], in which at every collision the electron velocity is reset to the average velocity at the given temperature, with a random direction, c.f.r. (6.1) and (6.2).
(2) the Gaussian model in which the total kinetic energy is kept constant by a thermostat force
$$m\underset{¯}{\overset{¨}{x}}_i=\underset{¯}{E}\frac{m\underset{¯}{E}\underset{¯}{J}}{3k_BT}\underset{¯}{\overset{\dot{}}{x}}_i+\mathrm{`}\mathrm{`}\mathrm{collisional}\mathrm{forces}^{\prime \prime }$$
$`(7.1)7.1`$
where $`3Nk_BT`$ is the total kinetic energy (a constant of motion in this model). The model has been widely studied and it was introduced by Hoover and Evans (see for instance \[HHP87\] and \[EM90\]).
(3) a “friction model” in which particles independently experience a constant friction
$$m\underset{¯}{\overset{¨}{x}}_i=\underset{¯}{E}\nu \underset{¯}{\overset{\dot{}}{x}}_i+\mathrm{`}\mathrm{`}\mathrm{collisional}\mathrm{forces}^{\prime \prime }$$
$`(7.2)7.2`$
where $`\nu `$ is a constant tuned so that the average kinetic energy is $`eNk_BT/2`$. This model was considered in the perspective of the conjectures of ensemble equivalence in \[Ga95\], \[Ga96b\].
The first model is a “stochastic model” while the second and third are deterministic: the third is “irreversible” while the second is reversible because the involution $`I(\underset{¯}{x}_i,\underset{¯}{v}_i)=(\underset{¯}{x}_i,\underset{¯}{v}_i)`$ anticommutes with the time evolution flow $`S_t`$ defined by the equation (7.1): $`IS_t=S_tI`$ (as the “friction term” is odd under $`I`$).
Let $`\mu _{\delta ,T}`$ be the SRB distribution for (7.1) for the stationary state that is reached starting from initial data with energy $`3Nk_BT/2`$. The collection of the distributions $`\mu _{\delta ,T}`$ as the kinetic energy $`T`$ and the density $`\delta =N/V`$ vary, define a “statistical ensemble” $``$ of stationary distributions associated with the equation (7.1).
Likewise we call $`\stackrel{~}{\mu }_{\delta ,\nu }`$ the class of SRB distributions associated with (7.2) which forms an “ensemble” $`\stackrel{~}{}`$.
We establish a correspondence between distributions of the ensembles $``$ and $`\stackrel{~}{}`$: we say that $`\mu _{\delta ,T}`$ and $`\stackrel{~}{\mu }_{\delta ^{},\nu }`$ are “corresponding elements” if
$$\delta =\delta ^{},T=\frac{1}{2}(\underset{j}{}m\underset{¯}{\overset{\dot{}}{x}}_j^2)\stackrel{~}{\mu }_{\delta ,\nu }(d\underset{¯}{x}d\underset{¯}{\overset{\dot{}}{x}})$$
$`(7.3)7.3`$
Then the following conjecture was proposed in \[Ga96b\].
Conjecture 1: (equivalence conjecture) Let $`F`$ be a “local observable”, i.e. an observable depending solely on the microscopic state of the electrons whose positions is inside a fixed box $`V_0`$. Then, if $``$ denotes the local smooth observables
$$\underset{N\mathrm{},N/V=\delta }{lim}\stackrel{~}{\mu }_{\delta ,\nu }(F)=\underset{N\mathrm{},N/V=\delta }{lim}\mu _{\delta ,T}(F)F$$
$`(7.4)7.4`$
if $`T`$ and $`\nu `$ are related by (7.3).
This conjecture has been discussed in \[Ga95\], sec. 5, and \[Ga96a\], see sec. 2 and 5: and in \[Ru99b\] arguments in favor of it have been developed.
Clearly the conjecture is very similar to the equivalence in equilibrium between canonical and microcanonical ensembles: here the friction $`\nu `$ plays the role of the canonical inverse temperature and the kinetic energy that of the microcanonical energy.
It is remarkable that the above equivalence suggests equivalence between a “reversible statistical ensemble”, i.e. the collection $``$ of the SRB distributions associated with (7.1) and a “irreversible statistical ensemble”, i.e. the collection $`\stackrel{~}{}`$ of SRB distributions associated with (7.2).
Furthermore it is natural to consider also the collection $`^{}`$ of stationary distributions for the original stochastic model (1) of Drude, whose elements $`\mu _{\nu ,T}^{}`$ can be parameterized by the quantities $`T`$, temperature (such that $`\frac{1}{2}_jm\underset{¯}{\overset{\dot{}}{x}}_j^2=\frac{3}{2}Nk_BT`$, and $`N/V=\delta `$). This is an ensemble $`^{}`$ whose elements can be put into one to one correspondence with the elements of, say, the ensemble $``$ associated with model (2), i.e. with (7.1): an element $`\mu _{\nu ,T}^{}^{}`$ corresponds to $`\mu _{\delta ,\nu }`$ if $`T`$ verifies (7.3). Then
Conjecture 2: If $`\mu _{\delta ,T}`$ and $`\mu _{\delta ,\nu }^{}^{}`$ are corresponding elements (i.e. (7.3) holds) then
$$\underset{N\mathrm{},N/V=\delta }{lim}\mu _{\delta ,T}(F)=\underset{N\mathrm{},N/V=\delta }{lim}\mu _{\delta ,T}^{}(F)F$$
$`(7.5)7.5`$
for all local observables $`F`$.
Hence we see that there can be statistical equivalence between a viscous irreversible dissipation model and either a stochastic dissipation model or a reversible dissipation model, at least as far as the averages of special observables are concerned.
The argument in \[Ru99b\] in favor of conjecture 1 is that the coefficient $`\alpha `$ in (6.3) is essentially the average $`J`$ of the current over the whole box containing the system of particles, $`J=N^1e_j\underset{¯}{\overset{\dot{}}{x}}_i`$: hence $`J`$ should be constant with probability $`1`$, at least if the stationary SRB distributions can be reasonably supposed to have some property of ergodicity with respect to space translations.
§8. Entropy driven intermittency in reversible dissipation.
A further argument for the equivalence conjectures in the above electric conduction models can be related to the fluctuation theorem: the quantity $`\alpha (x)`$ is also proportional to the phase space contraction rate $`\sigma (x)=(3N1)\alpha (x)`$. Therefore, denoting in general with a subscript $`+`$ the SRB average (or the time average) of an observable, the probability that $`\sigma (x)`$ deviates from its average $`\sigma _+=(3N1)\alpha _+`$ can be studied as follows.
If the number $`N`$ of particles is large the time scale $`\tau _0`$ over which $`\sigma (S_tx)`$ evolves will be large compared to the microscopic evolution rates, because $`\sigma _t(x)`$ is the sum of the $`6N`$ rates of expansion and contraction of the $`6N`$ phase space directions out of $`x`$ (sometimes called the “local Lyapunov exponents”).<sup>4</sup>
<sup>4</sup>The exact number of exponents depends on how many constants of motion the system has: for instance in the case of the conduction model (1) in §6 above the number of exponents is $`6N1`$ because the kinetic energy is conserved and the system has no other (obvious) first integrals. Furthermore one of such exponents is $`0`$ since every dynamical system in continuous time has one zero exponent (corresponding to the direction $`\dot{x}`$ of the flow).
Consider a large number $`m`$ of time intervals $`I_1,I_2,\mathrm{},I_m`$ of size $`\tau _0`$ and let $`\sigma _j`$ be the (average) value of $`\sigma (S_tx)`$ for $`tI_j`$. Then the fraction of the $`j`$’s such that $`\sigma _j\sigma _+\sigma _+p`$ will be proportional to
$$\pi _{\tau _0}(p)e^{\tau _0\zeta (p)}$$
$`(8.1)8.1`$
and $`\zeta (p)<\zeta (1)`$ if $`p1`$. Since we can expect that$`\zeta (p)`$ is proportional to $`N`$ we see that the fraction of time intervals $`I_j`$ in which $`\sigma _j\sigma _+`$ will be exponentially small with $`N`$.
For instance the fraction of time intervals in which $`\sigma _j\sigma _+`$ will be, by the fluctuation theorem
$$e^{(3N1)\alpha _+\tau _0}$$
$`(8.2)8.2`$
In order that the above argument holds it is essential that $`N`$ is large to the point that we can think that the time scale $`\tau _0`$ over which $`\sigma (S_tx)`$ varies is much larger than the microscopic scales: so that we can regard $`\tau _0`$ large enough for the fluctuation theorem to apply. In this respect this is not really different from the previously quoted argument in \[Ru99b\]. However the change of perspective gives further information.
In fact we get the following picture: $`N`$ is large and for most of the time the (stationary) evolution uneventfully proceeds as if $`\sigma (S_tx)\sigma _+`$ (thus justifying conjecture 1). Very rarely, however, it proceeds as if $`\sigma (S_tx)\sigma _+`$, for instance with $`\sigma (S_tx)=\sigma _+`$: such “bursts of anomalous behavior” occur very rarely. But when they occur “everything else goes the wrong way” because, as discussed in detail in \[Ga99c\], while the phase space contraction is opposite to what it “should be” (in the average) then it also happens that all observables evolve following paths that are the time reversal of the expected paths, This is the content, see \[Ga99c\], of the following theorem which is quite close (particularly if one examines its derivation) to the Machlup–Onsager theory of fluctuation patterns (note that, however, it does not require closeness to equilibrium)
Theorem(conditional reversibility theorem): If $`F`$ is an observable with even (or odd), for simplicity, time reversal parity and if $`\tau `$ is large then the evolution or “fluctuation pattern” $`\phi (t)`$ and its time reversal $`I\phi (t)\phi (t)`$, $`t[\tau _0/2,\tau _0/2]`$, will be followed with equal likelihood if the first is conditioned to an entropy creation rate $`p`$ and the second to the opposite $`p`$.
In other words systems with reversible dynamics can be equivalent to systems with irreversible dynamics but they show “intermittent behavior” with intermittency lapses that become extremely rare very quickly as $`N\mathrm{}`$. Sometimes they can be really dramatic, as in the cases in which $`\sigma =\sigma _+`$: alas they are unobservable just for this reason and one can wonder (see §9 below) whether this is really of any interest.
§9. Local fluctuations and observable intermittency.
As a final comment upon the analysis of the equivalence of ensembles attempted above we consider a very large system with volume $`V`$ and a small subsystem of volume $`V_0`$ which is large but not yet really macroscopic so that the number of particles in $`V_0`$ is not too large, a nobler way to express the same notion is to say that we consider a “mesocopic” subsystem of our macroscopic system.
Here it is quite important to specify the system because we want to make use of aspects of the equivalence conjectures that are model dependent. Therefore we consider the conduction models (2) or (3) of §5: these are models in which dissipation occurs “homogeneously” throughout the system. In this case we can imagine to look at the part of the system in the box $`V_0`$: if $`j_1,\mathrm{},j_{N_0}`$ are the particles which at a certain instant are inside $`V_0`$ and $`\underset{¯}{\overset{\dot{}}{x}}_j=\underset{¯}{f}_j(\underset{¯}{x})`$ are the equations of motion we can define
$$\sigma _{V_0}(x)=\underset{i=1}{\overset{N_0}{}}_{x_{j_k}}f_{j_k}(x)$$
$`(9.1)9.1`$
which is (by definition) the part of phase space contraction due to the particles in $`V_0`$.
Since the part of the system inside the microscopically large but macroscopically small $`V_0`$ can be regarded as a new dynamical system whose properties should not be different from the ones of the full system enclosed in the full volume $`V`$ we may expect that the subsystem inside $`V_0`$ is in a stationary state and the quantity $`\sigma _{V_0}`$ has the same fluctuation properties as $`\sigma _V`$, i.e.
$$\begin{array}{cc}\hfill (1)& \sigma _{V_0}_+=V_0\overline{\sigma }_+,\sigma _V_+=V\overline{\sigma }_+\hfill \\ \hfill (2)& \pi _\tau ^{V_0}(p)=e^{\overline{\zeta }(p)\tau V_0},\pi _\tau ^V(p)=e^{\overline{\zeta }(p)\tau V}\hfill \end{array}$$
$`(9.2)9.2`$
where $`\overline{\zeta },\overline{\sigma }_+`$ are the same for $`V,V_0`$ and $`p=\tau ^1_{\tau /2}^{\tau /2}\sigma _{V_0}(S_tx)/\sigma _{V_0}𝑑t`$ or respectively $`p=\tau ^1_{\tau /2}^{\tau /2}\sigma _V(S_tx)/\sigma _V𝑑t`$. Here $`\sigma _{V_0}`$ is naively defined as the contribution to $`\sigma `$ coming from the particles in $`V_0`$.
In other words in large stationary systems with homogeneous reversible dissipation phase space contractions fluctuate in an extensive way, i.e. they are regulated by the same deviation function $`\overline{\zeta }(p)`$ (volume independent).
This is very similar to the well known property of equilibrium density fluctuations in a gas of density $`\rho `$: if $`VV_0`$ are a very large volume $`V`$ in a yet larger container and $`V_0`$ is a small but microscopically large (i.e. mesoscopic) volume $`V_0`$ then the total numbers of particles in $`V`$ and $`V_0`$ will be $`N`$ and $`N_0`$ and the average numbers will be $`\rho V`$ and $`\rho V_0`$ respectively. Then setting
$$p=(N\rho V)/\rho V,or,p=(N_0\rho V_0)/\rho V_0$$
$`(9.3)9.3`$
the probability that the variable $`p`$ has a given value will be proportional to
$$\pi ^V(p)=e^{\overline{\zeta }(p)V},\pi ^{V_0}(p)=e^{\overline{\zeta }(p)V_0}$$
$`(9.4)9.4`$
again with the same function $`\overline{\zeta }(p)`$.
This means that we can observe $`\overline{\zeta }(p)`$ by performing fluctuations experiments in small boxes, ideally carved out of the large container, where the density fluctuations are not too rare. A “local fluctuation law” should hold more generally in cases of models in which dissipation occurs homogeneously across the system, like the above considered conduction models.
The intuitive picture for the above “local fluctuation relation” inspired (and was substantiated) a mathematical model in which a local fluctuation relation can be proved as a theorem: it has een discussed in \[Ga99c\], see also below.
Going back to the conduction model we see that the intermittency phenomena discussed above can be actually observed by looking at the fluctuations of the contribution to phase space contraction due to a small subsystem.
And such “entropy driven” intermittency will be model independent for models which are equivalent in the sense of the previous sections provided the models used are equivalent and one of them is reversible.
An extreme case is provided by models (1)%(3), §7, for electric conduction (conjectured to be equivalent, see §7). In fact at first the model (3), the viscous thermostat, might look uninteresting as, obviously, in this case
$$\sigma ^V(x)=\mathrm{\hspace{0.17em}3}N\nu ,\sigma ^{V_0}(x)=\mathrm{\hspace{0.17em}3}N_0\nu $$
$`(9.5)9.5`$
and $`\sigma ^V/V`$ has no fluctuations.
However the equivalence conjecture makes a statement about expectation values of the same observable: hence we should consider the quantity $`\widehat{\sigma }^{V_0}(x)=\underset{¯}{E}\underset{¯}{J}_{V_0}/_j\underset{¯}{\overset{\dot{}}{x}}_j^2`$ and we should expect that its statistics with respect to an element of the ensemble $`^{}`$ is the same as that of the same quantity with respect to the corresponding elements of the ensembles $`,\stackrel{~}{}`$. Hence in particular the functions $`\overline{\zeta }(p)`$ which control the large fluctuations of $`\sigma ^V(p)`$ will verify
$$\overline{\zeta }(p)=\overline{\zeta }(p)p\widehat{\sigma }^{V_0}_+/V_0=\overline{\zeta }(p)3\rho \nu p=\overline{\zeta }(p)\frac{eEmJ_+}{k_BT}p$$
$`(9.6)9.6`$
where the first equality expresses the validity of a fluctuation theorem type of relation due to the fact that the small system, by the equivalence conjecture, should behave as a closed system; the second equality expresses a consequence of the equivalence conjecture between models (2) and (3) while the third is obtained by expressing the current via Drude’s theory (again assuming the conjectures of equivalence 1,2 of §7).
§10. Fluids.
The chaotic hypothesis was originally formulated to understand developed turbulence, \[Ru78\]: it is therefore interesting to revisit fluid motions theory.
The incompressible Navier Stokes equation for a velocity field $`\underset{¯}{u}`$ in a periodic container $`V`$ of side $`L`$ can be considered as an equation for the evolution in time of its Fourier coefficients $`\underset{¯}{u}_{\underset{¯}{k}}`$ where the “mode$`\underset{¯}{k}`$ has the form $`2\pi L^1\underset{¯}{n}`$ with $`\underset{¯}{n}\underset{¯}{0}`$ and $`\underset{¯}{n}`$ an integer components vector.<sup>5</sup>
<sup>5</sup>The value $`\underset{¯}{n}=\underset{¯}{0}`$ is excluded because, having periodic boundary conditions, it is not restrictive to suppose that the space average of $`\underset{¯}{u}`$ vanishes (galilean invariance). The convention for the Fourier transform that we use is $`\underset{¯}{u}\left(\underset{¯}{x}\right)=_{\underset{¯}{k}}e^{i\underset{¯}{k}\underset{¯}{x}}\underset{¯}{u}_{\underset{¯}{k}}`$. Furthermore $`\underset{¯}{u}_{\underset{¯}{k}}=\overline{\underset{¯}{u}}_{\underset{¯}{k}}`$ and $`\underset{¯}{k}\underset{¯}{u}_{\underset{¯}{k}}\mathrm{\hspace{0.33em}0}`$. If $`p`$ is the pressure field and $`\underset{¯}{f}`$ a simple forcing we shall fix the ideas by considering $`\underset{¯}{f}(\underset{¯}{x})=f\underset{¯}{e}\mathrm{sin}\underset{¯}{k}_f\underset{¯}{x}`$ where $`\underset{¯}{k}_f`$ is some prefixed mode and $`\underset{¯}{e}`$ is a unit vector orthogonal to $`\underset{¯}{k}_f`$.
The Navier Stokes equation is then
$$\underset{¯}{\overset{\dot{}}{u}}+u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=\underset{¯}{}p+\underset{¯}{f}+\nu \mathrm{\Delta }\underset{¯}{u}$$
$`(10.1)10.1`$
and it is convenient to use dimensionless variables $`\underset{¯}{u}_0,p_0,\underset{¯}{\phi }_0,\underset{¯}{\xi },\tau `$: so we define them as
$$\begin{array}{cc}\hfill \underset{¯}{u}(\underset{¯}{x},t)=& fL^2\nu ^1\underset{¯}{u}_0(L^1\underset{¯}{x},L^2\nu t),\underset{¯}{\xi }=L^1\underset{¯}{x},\tau =L^2\nu t\hfill \\ \hfill p(\underset{¯}{x},t)=& fLp_0(L^1\underset{¯}{x},L^2\nu t),R\stackrel{def}{=}fL^3\nu ^2\hfill \\ \hfill \underset{¯}{f}(\underset{¯}{x},t)=& f\underset{¯}{\phi }_0(L^1\underset{¯}{x})\hfill \end{array}$$
$`(10.2)10.2`$
with $`\mathrm{max}|\underset{¯}{\phi }_0|=1`$. The result, dropping the label $`0`$ and calling again $`\underset{¯}{x},t`$ the new variables $`\underset{¯}{\xi },\tau `$, is that the Navier Stokes equations become an equation for a divergenceless field $`\underset{¯}{u}`$ defined on $`V=[0,1]^3`$, with periodic boundary conditions and equations
$$\underset{¯}{\overset{\dot{}}{u}}+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=\underset{¯}{}p+\underset{¯}{\phi }+\mathrm{\Delta }\underset{¯}{u},\underset{¯}{}\underset{¯}{u}=0$$
$`(10.3)10.3`$
with $`\mathrm{max}|\underset{¯}{\phi }|=1`$.
Equation (10.3) is our model of fluid motion, where $`R`$ plays the role of “forcing intensity” and the term $`\mathrm{\Delta }\underset{¯}{u}`$ represents the “thermostatting force”. As $`R`$ varies the stationary distributions $`\mu _R`$ which describe the SRB statistics of the motions (10.3) define a set $``$ of probability distributions which forms an “ensemble”.
The mathematical theory of the Navier Stokes equations is far from being understood: however phenomenology establishes quite clearly a few key points. The main property is that if (10.3) is written as an equation for the Fourier components of $`\underset{¯}{u}`$ then one can assume that $`\underset{¯}{u}_{\underset{¯}{k}}\underset{¯}{0}`$ for $`|\underset{¯}{k}|>K(R)`$, for some finite $`K(R)`$.
Therefore the equation (10.3) should be thought of as a “truncated equation” in momentum space by identifying it with the equation obtained by projecting also $`u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}`$ on the same function space.
Should one develop anxiety about the mathematical aspects of the Navier Stokes equation one should therefore think that an equally good model for a fluid is the mentioned truncation provided $`K(R)`$ is chose large enough.
The idea is that for $`K(R)=R^\kappa `$, with $`\kappa `$ larger than a suitable $`\kappa _0`$ the results of the theory, i.e. the statistical properties of $`\mu _R`$ become $`\kappa `$–independent for $`R`$ large.
The simplest evaluation of $`\kappa _0`$ gives $`\kappa _0=9/4`$ as a consequence of the so called K41 theory of homogeneous turbulence, see \[LL71\].
If (10.3) is a good model for a fluid when $`L`$ is large then it provides us with an “ensemble” $``$ of SRB distributions (on the space of the velocity fields components $`\underset{¯}{u}_{\underset{¯}{k}}`$ of dimension $`8\pi K(R)^3/3`$).<sup>6</sup>
<sup>6</sup>There are about $`4\pi K\left(R\right)^3/3`$ vectors with integer components inside a sphere of radius $`K\left(R\right)`$, thus the number of complex Fourier components with mode label $`\left|\underset{¯}{k}\right|<K\left(R\right)`$ would be $`3`$ times as much, but the divergenceless condition leaves only $`2`$ complex components for $`\underset{¯}{u}_{\underset{¯}{k}}`$ along the two unit vectors orthogonal to $`\underset{¯}{k}`$ and the reality condition further divides by $`2`$ the number of “free” components.
We should expect, following the discussion of the statistical mechanics cases, that there can be other “ensembles” $`\stackrel{~}{}`$ which are equivalent to $``$.
Here $`R`$ plays the role of the volume in non equilibrium statistical mechanics, so that $`R\mathrm{}`$ will play the role of the thermodynamic limit, a limit in which the effective number of degrees of freedom, $`4\pi R^{3\kappa }/3`$, becomes infinite. The role of the local observables will be plaid by the (smooth) functions $`F(\underset{¯}{u})`$ of the velocity fields $`\underset{¯}{u}`$ which depend on $`\underset{¯}{u}`$ only via its Fourier components that have mode $`\underset{¯}{k}`$ with $`|\underset{¯}{k}|<B`$ for some $`B`$: $`F(\underset{¯}{u})=F(\{\underset{¯}{u}_{\underset{¯}{k}}\}_{|\underset{¯}{k}|B})`$.
We shall call $``$ the space of such observables: examples can be obtained by setting $`F(\underset{¯}{u})=|e^{i\underset{¯}{k}\underset{¯}{x}}\underset{¯}{u}(\underset{¯}{x})𝑑\underset{¯}{x}|^2`$ or $`F(\underset{¯}{u})=\underset{¯}{\psi }(\underset{¯}{x})\underset{¯}{u}(\underset{¯}{x})𝑑\underset{¯}{x}`$ where the function has only a finte number of harmonics, $`\underset{¯}{\psi }(\underset{¯}{x})=_{|\underset{¯}{k}|<B}e^{i\underset{¯}{k}\underset{¯}{x}}\underset{¯}{u}(\underset{¯}{x})d\underset{¯}{x}`$, etc.
As in non equilibrium statistical mechanics we can expect that the equations of motion themselves become part of the definition of the ensembles. For instance one can imagine defining the ensemble $`\stackrel{~}{}`$ of the SRB distributions $`\stackrel{~}{\mu }_V`$ for the equations
$$\underset{¯}{\overset{\dot{}}{u}}+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=\underset{¯}{}p+\underset{¯}{\psi }+\nu (\underset{¯}{u})\mathrm{\Delta }\underset{¯}{u}$$
$`(10.4)10.4`$
called GNS equations in \[Ga97a\], or “gaussian Navier Stokes” equations, where $`\nu (\underset{¯}{u})`$ is so defined that
$$\mathrm{\Xi }=_V(_{\stackrel{~}{}\text{ }}\underset{¯}{u})^2𝑑\underset{¯}{x}/(2\pi )^3=\underset{\underset{¯}{k}}{}\underset{¯}{k}^2|\underset{¯}{u}_{\underset{¯}{k}}|^2$$
$`(10.5)10.5`$
is exactly constant and equal to $`\mathrm{\Xi }`$. The equations (10.4) are interpreted as above with the same momentum cut off $`K(R)=R^\kappa `$.
An element $`\stackrel{~}{\mu }_\mathrm{\Xi }`$ of $`\stackrel{~}{}`$ and one $`\mu _R`$ of $``$, SRB distributions for the two different dynamics (10.3) and (10.4), “correspond to each other” if
$$\mathrm{\Xi }=\mu _R(d\underset{¯}{u})\left(_V(_{\stackrel{~}{}\text{ }}\underset{¯}{u})^2𝑑\underset{¯}{x}/(2\pi )^3\right)\stackrel{def}{=}\mathrm{\Xi }_R$$
$`(10.6)10.6`$
where $`\mu _R`$ is the SRB distribution at Reynolds number $`R`$ for the previous viscous Navier Stokes equation, (10.3), and we naturally conjecture
Conjecture 3(equivalence GNS–NS): If $`R\mathrm{}`$ then for all local observables $`F`$ it is $`\mu _R(F)=\stackrel{~}{\mu }_{\mathrm{\Xi }_R}(F)`$ if (10.6) holds.
It is easy to check that the GNS model “viscosity” $`\nu (\underset{¯}{u})`$, having to be such that the quantity $`\mathrm{\Xi }`$ in (10.5) is exactly constant must be
$$\nu (\underset{¯}{u})=\frac{_V\left(\underset{¯}{\phi }\mathrm{\Delta }\underset{¯}{u}R\mathrm{\Delta }\underset{¯}{u}(u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u})\right)𝑑\underset{¯}{x}}{_V(\mathrm{\Delta }\underset{¯}{u})^2𝑑\underset{¯}{x}}$$
$`(10.7)10.7`$
and we see that while (10.3) is an irreversible equation the (10.4) is reversible, with time reversal symmetry given by
$$I\underset{¯}{u}(\underset{¯}{x},t)=\underset{¯}{u}(\underset{¯}{x},t)$$
$`(10.8)10.8`$
as one can check.
More generally one may wish to leave the “Kolmogorov parameter” $`\kappa `$ as a free parameter: in this case the SRB distributions will form an ensemble whose elements can be parameterized by $`R,\kappa `$ and the equivalence conjecture can be extended to this case yielding equivalence between $`\mu _{R,\kappa }`$ and $`\stackrel{~}{\mu }_{\mathrm{\Xi },\kappa }`$. This is of interest, particularly if one has numerical experiments in mind.
If $`\kappa >\kappa _0`$ then the value of $`\kappa `$ should be irrelevant: but if $`\kappa <\kappa _0`$ the phenomenology will be different from the one of the Navier Stokes equation and equivalence might still hold but one cannot expect either equation to have the properties that we expect for the usual Navier Stokes equations (i.e. in this situation one would have to be careful in making statements based on common experience).
If we take $`\kappa `$ to be exactly equal to the value $`\kappa _0=9/4`$ (i.e. if we take the ultraviolet cut–off to be such that, according to the K41 theory, for larger values it is needlessly large and for lower values it is incorrectly low and shows a phenomenology which will depend on its actual value) then we may speculate that the “attracting set” is the full phase space (available compatibly with the constraint $`\mathrm{\Xi }=\mathrm{\Xi }_R`$). Therefore the divergence of the equations of motion, which is given by a rather involved expression in which only the first term seems to dominate at large $`R`$, namely
$$\begin{array}{cc}& \sigma (\underset{¯}{u})=(\underset{|\underset{¯}{k}|<K(R)}{}\underset{¯}{k}^2)\nu (\underset{¯}{u})(_V\mathrm{\Delta }\underset{¯}{\phi }\mathrm{\Delta }\underset{¯}{u}d\underset{¯}{x})(_V[(\mathrm{\Delta }\underset{¯}{u})^2\hfill \\ & R\mathrm{\Delta }\underset{¯}{u}(\mathrm{\Delta }(u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}))R(\mathrm{\Delta }u_{\stackrel{~}{}\text{ }})(\mathrm{\Delta }\underset{¯}{u})(_{\stackrel{~}{}\text{ }}\underset{¯}{u})R\mathrm{\Delta }\underset{¯}{u}(\mathrm{\Delta }_{\stackrel{~}{}\text{ }}\underset{¯}{u})u_{\stackrel{~}{}\text{ }}+\hfill \\ & +\nu (\underset{¯}{u})\mathrm{\Delta }\underset{¯}{u}\mathrm{\Delta }^2\underset{¯}{u}]d\underset{¯}{x})/_V(\mathrm{\Delta }\underset{¯}{u})^2d\underset{¯}{x}\hfill \end{array}$$
$`(10.9)10.9`$
will verify the fluctuation theorem, i.e. the rate function $`\zeta (p)`$ for the average phase space contraction $`p=\tau ^1_{\tau /2}^{\tau /2}\sigma (S_t\underset{¯}{u})𝑑t/\tau \sigma _+`$ will be such that $`\zeta (p)=\zeta (p)p\sigma _+`$.
If the chaotic hypothesis is valid together with the equivalence conjecture the validity of the fluctuation relation can be taken as a criterion for determining $`\kappa `$: it would be the last $`\kappa `$ before which the fluctuation relation between $`\zeta (p)`$ and $`\zeta (p)`$ holds. However this conclusion can only be drawn if the attracting set in phase space is the full ellipsoid $`\mathrm{\Xi }=\mathrm{\Xi }_R`$ at least for $`K(R)=R^{\kappa _0}`$.
The latter property might not be realized: and in such case the fluctuation theorem does not apply directly, although the equivalence conjectures still hold. In fact one can try to extend the fluctuation theorem to cover reversible cases in which the attracting set is smaller than the full phase space left available by the constraints. In such cases under suitable geometric assumptions, \[BG97\] and the earlier work \[BGG97\], one can derive a relation like
$$\zeta (p)=\zeta (p)p\sigma _+\vartheta ,0\vartheta 1$$
$`(10.10)10.10`$
where $`\vartheta `$ is a coefficient that can be related to the Lyapunov spectrum of the system, c.f.r. \[BG97\], \[Ga97a\]. In fact numerical work to check the theory proposed in \[Ga97a\] is currently being performed (private communication by Rondoni and Segre) with not too promising results which, optimistically, can be attributed to the fact that the ultraviolet cut off is too small due to numerical litmitations: clearly there is more work to do here. The preliminary numerical results give, so far, the somewhat surprising linearity in $`p`$ but with a slope that, although of the correct order of magnitude, seems to have a value that does not match the theory within the error bounds.
Coming back to the Navier Stokes equation we mention that we may imagine to write it as (10.3) but with the different constraint
$$U=_V\underset{¯}{u}^2𝑑\underset{¯}{x}=\mathrm{const}$$
$`(10.11)10.11`$
rather than (10.5).
This case has been considered in \[RS99\] and the multiplier $`\nu (\underset{¯}{u})`$ is in this case
$$\nu (\underset{¯}{u})=\frac{_V\underset{¯}{\phi }\underset{¯}{u}𝑑\underset{¯}{x}}{_V\underset{¯}{u}^2𝑑\underset{¯}{x}},\sigma (\underset{¯}{u})=(3\underset{|\underset{¯}{k}|<K(R)}{}|\underset{¯}{k}|^21)\nu (\underset{¯}{u})$$
$`(10.12)10.12`$
and we can (almost) repeat the above considerations and equivalence conjectures. This constraint is a gaussian constraint that $`U`$ is constant obtained by imposing its constancy on the Euler evolution via Gauss’ principle with a suitable definition of the notion of “constraint effort” (this notion is not unique, see \[Ga97a\] for another definition) and we do not discuss it here to avoid overlapping with §12 below.
Thd intuitive motivation for the equivalence conjectures is that for large $`R`$ the phase space contraction $`\sigma (\underset{¯}{u})`$ and the coefficient $`\nu (\underset{¯}{u})`$<sup>7</sup>
<sup>7</sup>Which in the case (10.9) are simply proportional and in the case of (10.4) they are related in a more involved way, see (10.8),(10.9), but which are still probably proportional to leading order as $`R\mathrm{}`$. are “global quantities” and depend on the global properties of the system (e.g. $`\sigma (\underset{¯}{u})`$ is the sum of all the local Lyapunov exponents of the system whose number is $`O(K(R)^3)`$): they will “therefore” vary over time more slowly than any time scale of the system and can be considered constant.
The argument is not very convincing in the case of the equations with the constraint (10.11) because the $`\sigma (\underset{¯}{u})`$ in (10.12) is proportional to $`_V\underset{¯}{\phi }\underset{¯}{u}𝑑\underset{¯}{x}`$ which clearly depends only on harmonics of $`\underset{¯}{u}`$ with $`\underset{¯}{k}`$ small, i.e. it is a “local observable”. Note that this does not apply to the GNS equations with the constrained vorticity $`\mathrm{\Xi }`$, (10.6) where the “main” contribution to $`\sigma (\underset{¯}{u})`$, see (10.7), comes from the term proportional to $`R`$ which contains all harmonics. Therefore the result in \[RS99\] about the equivalence between the GNS equations, (10.4) with the constraint (10.5), and the equations with constraint (10.11) is interesting and puzzling: it might be an artifact of the smallness of the cut off that one has to impose in order to have numerically feasible simulations.
Finally $`\sigma _+(\underset{¯}{u})/\sigma _+`$, i.e. essentially $`\nu (\underset{¯}{u})/\nu _+`$ will fluctuate taking values sensibly different from their average value $`1`$, at very rare intervals of time: but when such fluctuations will occurr one shall see “bursts” of anomalous behavior: i.e. the motion will be “intermittent” as in the case discussed in non equilibrium statistical mechanics.
11. Entropy creation rate and entropy driven intermittency.
Of course if $`R`$ is large the number of degrees of freedom is large and intermittency on the scale of the fluid container will not be observable due to its extreme unlikelyhood (expected and quantitatively predicted by the fluctuation theorem).
Therefore we look also here, in fluid motions, for a local fluctuation relation. Fluids seem particularly suitable for verifying such local fluctuations relations because dissipation occurs homogeneously, i.e. friction strength is translation invariant.
This implies that we can regard a very small volume $`V_0`$ of the fluid as a system in itself (as always done in the derivation of the basic fluid equations, e.g. see \[Ga97b\]) and we can expect that the phase space contraction due to such volume elements is simply $`\sigma (\underset{¯}{u})`$, given by (10.9) or (10.12) (“equivalently” because of our equivalence conjectures) with the integrals in the numerator and denominator being extended to the volume $`V_0`$ rather than to the whole box, and expressing (essentially by definition) the “local phase space contraction” $`\sigma _{V_0}(\underset{¯}{u})`$.
Then $`p=\tau ^1_{\tau /2}^{\tau /2}\sigma _{V_0}(S_t\underset{¯}{u})/\sigma _{V_0}_+`$ will have a rate function $`\zeta (p)`$ which will verify, under the same assumptions as in (10.10), a large deviation relation as
$$\zeta (p)=\zeta (p)p\sigma _{V_0}_+\vartheta $$
$`(11.1)11.1`$
for some $`\vartheta `$: as mentioned the theoretical value of this slope $`\vartheta `$ seems currently inaccessible to theory (as the theory proposed in \[BG67\], \[Ga97a\] may need substantial modifications, c.f.r. comment following (10.10)). The $`\sigma _{V_0}`$ and $`\zeta (p)`$ will be proportional to $`V_0`$: $`\zeta (p)=V_0\overline{\zeta }(p)`$ with a $`V_0`$–independent $`\overline{\zeta }(p)`$. Note that $`\zeta (p)`$ depends also on $`R`$.
The small volume element of the fluid will therefore be subject to rather frequent variations: in spite of $`\zeta (p)`$ being proportional to $`V_0`$, because now $`V_0`$ is not large. The consequent intermittency phenomena can therefore be observed. And as in §9 once the phase space contraction is intermittent all properties of the system show the same behavior.
And in fact intermittency in observations averaged ove a time span $`\tau `$ will appear with a time frequency of the form $`e^{V_0(\overline{\zeta }(p)\overline{\zeta }(1))\tau }`$: the quantity $`p`$ can be interpreted as a measure of the “strength of intermittency” observable in easurements averaged over a time $`\tau `$ because as noted in §9 and in \[Ga99b\] the size of $`p`$ controls the statistical properties of “most” other observables. Therefore the function $`\overline{\zeta }(p)`$ (hence $`\zeta (p)`$) might be directly measurable and it should be rather directly related to the quantities that one actually observes in intermittency experiments. And the difference $`\zeta (p)\zeta (p)`$ can be tested for linearity in $`p`$ as predicted by the analysis above.
Note that in an extended system the volume $`V`$ is much larger than $`V_0`$ and we shall see “for sure” intermittency (for observables averaged over a time $`\tau `$) of strength $`p`$ in a region of volume $`V_0`$ somewhere within a volume $`W`$ such that
$$\frac{W}{V_0}e^{V_0(\overline{\zeta }(p)\overline{\zeta }(1))\tau }1$$
$`(11.2)11.2`$
At this point it seems relevant to recall that it is rather heatedly being debated whether the name of “entropy creation rate” that some authors (including the present one) give to the phase space contraction rate is justified or not, see \[An82\]. The above properties not only propose the physical meaning of the quantity $`p`$ and bring up the possibility of measuring its rate function $`\zeta (p)`$ in actual experiments but also provide a further justification of the name given to $`\sigma `$ as “entropy creation rate” and fuel for the debate that inevitably the word entropy generates at each and every occurrence.
§12. Benard convection, intermittency and the Ciliberto–Laroche experiment.
A very interesting attempt at checking some of the above ideas has been made recently by Cilberto and Laroche in an experiment on real fluids which has been performed with the aim of testing the relation (11.1) locally in a small volume element, \[CL98\]. By “real” we mean here non numerical: a distinction that, however, has faded away together with the XX–th century but that some still cherish: the system is physically macroscopic (water in a container of a size of the order of a liter).
This being a real experiment one has to stretch quite a bit the very primitive theory developed so far in order to interpret it and one has to add to the chaotic hypothesis other assumptions that have been discussed in \[BG97\], \[Ga97a\] in order to obtain the fluctuation relation (10.12) and its local couterpart (11.1).
The experiment attempts at measuring a quantity that is eventually interpreted as the difference $`\zeta (p)\zeta (p)`$, by observing the fluctuations of the product $`\vartheta u^z`$ where $`\vartheta `$ is the deviation of the temperature from the average temperature in a small volume element $`\mathrm{\Delta }`$ of water at a fixed position in a Couette flow and $`u^z`$ is the velocity in the $`z`$ direction of the water in the same volume element.
The result of the experiment is in a way quite unexpected: it is found that the function $`\zeta (p)`$ is rather irregular and lacking symmetry around $`p=1`$: nevertheless the function $`\zeta (p)\zeta (p)`$ seems to be strikingly linear. As discussed in \[Ga97a\], predicting the slope of the entropy creation rate would be difficult but if the equivalence conjecture considered above and discussed more in detail in \[Ga97a\] is correct then we should expect linearity of $`\zeta (p)\zeta (p)`$.
In the experiment of \[CL98\] the quantity $`\vartheta u_z`$ did not appear to be the divergence of the phase space volume simply because there was no model proposed for a theory of the experiment. Nevertheless Ciliberto–Laroche select the quantity $`_\mathrm{\Delta }\vartheta u^z𝑑\underset{¯}{x}`$ on the basis of considerations on entropy and dissipation so that there is hope that in a model of the flow this quantity can be related to the entropy creation rate discussed in §10,§11.
Here we propose that a model for the fluid, that can be reasonably used, is Rayleigh’s model of convection, \[Lo63\], \[LL71\] and \[Ga97b\] sec. 5. An attempt for a theory of the experiment could be the following.
One supposes that the equations of motion of the system in the whole container (of linear size of the order of $`30cm`$) are written for the quantities $`t,x,z,\vartheta ,\underset{¯}{u}`$ in terms of the height $`H`$ of the container (assumed to be a horizontal infinite layer), of the temperature difference between top and bottom $`\delta T`$ and in terms of the phenomenological “friction constants$`\nu ,\chi `$ of viscosity, dynamical thermal conductivity and of the thermodynamic dilatation coefficient $`\alpha `$. We suppose that the fluid is $`3`$–dimensional but stratified, so that velocity and temperature fields do not depend on the coordinate $`y`$, and gravity is directed along the $`z`$–axis: $`\underset{¯}{g}=g\underset{¯}{e},\underset{¯}{e}=(0,0,1)`$. The temperature deviation $`\vartheta `$ is defined as the difference betwen the temperature $`T(x,y,z)`$ and the temperature that the fluid would have at height $`z`$ in absence of convection, i.e. $`T_0z\delta T/H`$ if $`T_0`$ is the bottom temperature.
In such conditions the equations, including the boundary conditions (of fixed temperature at top and bottom and zero normal velocity at top and bottom), the convection equations in the Rayleigh model, see \[Lo63\] eq. (17), (18) where they are called the Saltzman equations, and \[Ga97b\] §1.5, become
$$\begin{array}{ccc}& \underset{¯}{}\underset{¯}{u}=0,u_x𝑑\underset{¯}{x}=u_y𝑑\underset{¯}{x}=0\hfill & \\ & \underset{¯}{\overset{\dot{}}{u}}+u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=\nu \mathrm{\Delta }\underset{¯}{u}\alpha \vartheta \underset{¯}{g}\underset{¯}{}p^{},\hfill & (12.1)12.1\hfill \\ & \dot{\vartheta }+u_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\vartheta =\chi \mathrm{\Delta }\vartheta +\frac{\delta T}{H}u_z\vartheta (0)=0=\vartheta (H),u_z(0)=0=u_z(H),\hfill & \end{array}$$
The function $`p^{}`$ is related to the pressure $`p`$: within the approximations it is $`p=p_0\rho _0gz+p^{}`$. We shall impose for simplicity horizontal periodic boundary conditions in $`x,y`$ so that the fluid can be considered in a finite container $`V`$ of side $`a`$ for some $`a>0`$ prefixed (which in theoriginal variable would correspond to a container of horizntal size $`aH`$).
It is useful to define the following adimensional quantities
$$\begin{array}{cc}\hfill \tau & =t\nu H^2,\xi =xH^1,\eta =yH^1,\zeta =zH^1,\hfill \\ \hfill \vartheta ^0& =\frac{\alpha \vartheta }{\alpha \delta T},\underset{¯}{u}^0=(\sqrt{gH\alpha \delta T})^1\underset{¯}{u}\hfill \\ \hfill R^2& =\frac{gH^3\alpha \delta T}{\nu ^2},R_{Pr}=\frac{\nu }{\kappa }\hfill \end{array}$$
$`(12.2)12.2`$
and one checks that the Rayleigh equations take the form
$$\begin{array}{cc}& \underset{¯}{\overset{\dot{}}{u}}+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=\mathrm{\Delta }\underset{¯}{u}R\vartheta \underset{¯}{e}\underset{¯}{}p,\hfill \\ & \dot{\vartheta }+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\vartheta =R_{Pr}^1\mathrm{\Delta }\vartheta +Ru_z,\hfill \\ & \underset{¯}{}\underset{¯}{u}=0\hfill \\ & u_z(0)=u_z(1)=0,\vartheta (0)=\vartheta (1)=0,\hfill \\ & _Vu_x𝑑\underset{¯}{x}=_Vu_y𝑑\underset{¯}{x}=0\hfill \end{array}$$
$`(12.3)12.3`$
where we again call $`t,x,y,z,\underset{¯}{u},\vartheta `$ the adimensional coordinates $`\tau ,\xi ,\eta ,\zeta ,\underset{¯}{u}^0,`$ $`\vartheta ^0`$ in (12.2). The numbers $`R,R_{Pr}`$ are respectively called the Reynolds and Prandtl numbers of the problem: $`R_{Pr}=6.7`$ for water while $`R`$ is a parameter that we can adjust, to some extent, from $`0`$ up to a rather large value.
According to the principle of equivalence stated in \[Ga97a\] here one could impose the constraints
$$_V\left(\underset{¯}{u}^2+\frac{1}{R_{Pr}}\vartheta ^2\right)𝑑\underset{¯}{x}=C$$
$`(12.4)12.4`$
on the “frictionless equations” (i.e. the ones without the terms with the laplacians) and determine the necessary forces via Gauss’ principle of minimal effort, see footnote <sup>3</sup> and \[Ga96a\], \[Ga97a\]. We use as effort functional of an acceleration field $`\underset{¯}{a}`$ and of a temperature variation field $`s`$ the quantity
$$\begin{array}{ccc}& (\underset{¯}{a},s)\stackrel{def}{=}((\underset{¯}{a}+\underset{¯}{}p\underset{¯}{f}),(\mathrm{\Delta })^1(\underset{¯}{a}+\underset{¯}{}p\underset{¯}{f}))+\hfill & (12.5)12.5\hfill \\ & +((s\phi ),(\mathrm{\Delta })^1(s\phi ))\mathrm{with}\hfill & \\ & \underset{¯}{f}\stackrel{def}{=}R\vartheta \underset{¯}{e},\phi \stackrel{def}{=}Ru_z\hfill & \end{array}$$
and require it to be minimal over the variations $`\underset{¯}{\delta }(\underset{¯}{x})`$ of $`\underset{¯}{a}=\frac{d\underset{¯}{u}}{dt}`$ and $`\tau (\underset{¯}{x})`$ of $`s=\frac{d\vartheta }{dt}`$ with the constraints that for all $`\underset{¯}{x}`$ it is $`\underset{¯}{}\underset{¯}{\delta }=0`$, besides those due to the boundary conditions. The result is
$$\begin{array}{cc}& \underset{¯}{}\underset{¯}{u}=0\hfill \\ & \underset{¯}{\overset{\dot{}}{u}}+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=R\vartheta \underset{¯}{e}\underset{¯}{}p^{}+\underset{¯}{\tau }_{th}\hfill \\ & \dot{\vartheta }+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\vartheta =Ru^z+\lambda _{th}\hfill \\ & \vartheta (0)=0=\vartheta (H),_Vu_x𝑑\underset{¯}{x}=_Vu_y𝑑\underset{¯}{x}=0\hfill \end{array}$$
$`(12.6)12.6`$
where the frictionless equations are modified by the thermostats forces $`\underset{¯}{\tau }_{th},\lambda _{th}`$: the latter impose the nonholonomic constraint in (12.4) with the effort functional defined by (12.5). Looking only at the bulk terms we see that the equations obtained by imposing the con1straints via Gauss’ principle become the (12.3) with coefficients in front of the Laplace operators equal to $`\nu _G,\nu _GR_{Pr}^1`$, respectively, with the “gaussian multiplier” $`\nu _G`$ being an odd functions of $`\underset{¯}{u}`$, see \[Ga97a\]: setting $`\stackrel{~}{C}_V(\underset{¯}{u},\vartheta )=_V\left((_{\stackrel{~}{}\text{ }}\underset{¯}{u})^2+R_{Pr}^1(\underset{¯}{}\vartheta )^2\right)𝑑\underset{¯}{x}`$ one finds
$$\nu _G=\stackrel{~}{C}_V(\underset{¯}{u},\vartheta )^1R(1+R_{Pr}^1)_Vu^z\vartheta 𝑑\underset{¯}{x}$$
$`(12.7)12.7`$
And the equations become, finally
$$\begin{array}{cc}& \underset{¯}{}\underset{¯}{u}=0\hfill \\ & \underset{¯}{\overset{\dot{}}{u}}+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\underset{¯}{u}=R\vartheta \underset{¯}{e}\underset{¯}{}p^{}+\nu _G\mathrm{\Delta }\underset{¯}{u}\hfill \\ & \dot{\vartheta }+Ru_{\stackrel{~}{}\text{ }}_{\stackrel{~}{}\text{ }}\vartheta =Ru^z+\nu _G\frac{1}{R_{Pr}}\mathrm{\Delta }\vartheta \hfill \\ & \vartheta (0)=0=\vartheta (H),u_x𝑑\underset{¯}{x}=u_y𝑑\underset{¯}{x}=0\hfill \end{array}$$
$`(12.8)12.8`$
If one wants the equivalence between the ensembles of SRB distributions for the equation (12.8) and for (12.3) one has to tune, \[Ga97a\], the value of the constant $`C`$ in (12.4) so that the time average value $`\nu _G_+`$ of $`\nu _G`$ is precisely the physical one: namely $`\nu _G=1`$ by (12.3). This is (again) the same, in spirit, as fixing the temperature in the canonical ensemble so that it agrees with the microcanonical temperature thus implying that the two ensembles give the same averages to the local observables.
The equations (12.8) are time reversible (unlike the (12.3)) under the time reversal map:
$$(\underset{¯}{u},\vartheta )=(\underset{¯}{u},\vartheta )$$
$`(12.9)12.9`$
and they should be supposed, by the arguments in \[Ga97a\] and §10,11: “equivalent” to the irreversible ones (12.3),
The (12.8) should therefore have a “divergence” $`\sigma (\underset{¯}{u},\vartheta )`$ whose fluctuation function $`\zeta (p)`$ verifies a linear fluctuation relation, i.e. $`\zeta (p)\zeta (p)`$ should be linear in $`p`$. Note that the divergence of the above equations is proportional to $`\nu _G`$ if one supposes that the high momenta modes with $`|\underset{¯}{k}|>K(R)=R^\kappa `$ with $`\kappa `$ suitable can be set equal to $`0`$ so that the equation (12.8) becomes a system of finite differential equations for the Fourier components of $`\underset{¯}{u},\vartheta `$.
For instance the Lorenz’ equations, \[Lo63\] see also §17 of \[Ga97b\], reduced the number of Fourier components necessary to describe (12.3) to just three components, thus turning it into a system of three differential equations.
Proceeding in this way the divergence of the equations of motion can be computed as a sum of two integrals one of which proportional to $`\nu _G`$ in (12.7). If instead of integrating over the whole sample we integrate over a small region $`\mathrm{\Delta }`$, like in the experiment of \[CL98\], we can expect to see a fluctuation relation for the entropy creation rate if the fluctuation theorem holds locally, i.e. for the entropy creation in a small region.
As for the cases in §11 this is certainly not implied by the proof in \[GC95\]: however when the dissipation is homogeneous through the system, as it is the case in the Rayleigh model there is hope that the fluctuation relation holds locally because “a small subsystem should be equivalent to a large one”. As noted in §9 the actual possibility of a local fluctuation theorem in systems with homogeneous dissipation has been shown in \[Ga99c\], after having been found through numerical simulations in \[GP99\], and this example was relevant because it gave us some justification to imagine that it might apply to the present situation as well.
The entropy creation is due to the term $`R_\mathrm{\Delta }u_z\vartheta 𝑑\underset{¯}{x}/\stackrel{~}{C}_\mathrm{\Delta }(\underset{¯}{u},\vartheta )`$, where $`\mathrm{\Delta }`$ is the region where the measurements of \[CL98\] are performed, hence we have a proposal for the explanation of the remarkable experimental result. Unfortunately in the experiment \[CL98\] the contributions not explicitly proportional to $`R`$ to the entropy creation rates have not been measured nor has been the $`\widehat{C}`$ in (12.7) which also fluctuates (or might fluctuate). In any event they might be measurable by improving the same apparatus, so that one can check whether the above attempt to an explanation of the experiment is correct, or try to find out more about the theory in case it is not right. If correct the above “theory” the experiment in \[CL98\] would be quite important for the status of the chaotic hypothesis.
§13. Conclusions.
The chaotic hypothesis promises a point of view on non equilibrium that has proved so far of some interest. Here we have exposed the basic ideas and attempted at drawing some consequences: admittedly the most interesting rely on rather phenomenological and heuristic grounds. They are summarized below. (1) The definition of nonequilibrium ensembles with the proposal that out of equilibrium also the equation of motion should be considered as part of the definition of ensemble. This is take into account that while in equilibrium the system is uniquely defined by its microscopic forces and constituents in non equilibrium it is not so. Systems must be put in contact with thermostats if we want them to become stationary after a transient time. And (for large systems) there may be several equivalent ways of taking heat out of a system, i.e. several thermostats, without affecting the properties of stationary state that is eventually reached by the system itself.
(2) Equivalence of ensembles has the most striking aspect that systems which evolve with equations that are very different may exhibit the same statistical properties. In particular reversible evolutions might be equivalent to non reversible ones, thus making it possible to apply results that require reversibility, in particular the fluctuation relations, to cases in which it is not valid.
(3) An interpretation of the quantity $`p`$ that intervenes in the fluctuation theorems in terms of an intermittency phenomenon and as a further quantitative measure of it.
(4) The possibility of applying the theory to strongly turbulent motions was the origin of the Ruelle’s principle that evolved into the chaotic hypothesis: therefore not surprisingly the ideas can be applied to fluid dynamics. We have discussed a possible approach. The approach leads again to a proposal for the theory of certain intermittency phenomena which appear quantitatively related to entropy creation fluctuations.
(5) The possibility of measurement of the rate $`\zeta (p)`$ leads to a a possible prediction of the spatial frequency of internittent events of strength $`p`$ or, as I prefer, with entropy creation rate $`p`$ (see (4) above, (11.2) and §12). This seems testable in concrete experiments (both real and numerical).
(6) We have used the results in (2)%(8) to hint at an interpretation of the experiment by Ciliberto and Laroche on Benard convection in water.
Although the theory is still at its beginning and it migh turn out to be not really of interest it seems that at this moment it is worth trying to test it both in its safest, c.f.r. §2%§8, and in its most daring, c.f.r. §9%§12, predictions.
Acknowledgements: I am greatly indebted to Professor R. Newton for giving me the opportunity to collect the above thoughts which, collected together, go far beyond what I originally planned after he proposed to me to write this review. Work partially supported by Rutgers University and by MPI through grant # ??????????.
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Internet: Authors’ preprints downloadable (latest version) at:
http://ipparco.roma1.infn.it
(link) http://www.math.rutgers.edu/$``$giovanni
e-mail: giovanni.gallavotti@roma1.infn.it |
warning/0001/astro-ph0001317.html | ar5iv | text | # Statistics of clustering of ultra-high energy cosmic rays and the number of their sources
## I Introduction
Recent analysis of arrival directions of ultra-high energy cosmic rays (UHECR) reveals groups of events (clusters) with arrival directions lying within $`3^{}`$, the typical angular resolution of the experiment. The set of 92 observed events with energy $`E>4\times 10^{19}`$ eV contains 7 doublets and 2 triplets . The small probability of chance coincidence, of the order of $`10^3`$ , suggests that clustering is a result of the existence of compact sources. At higher energies, $`E>10^{20}`$ eV, one doublet out of 14 events is observed.
Compact sources of UHECR are naturally explained in astrophysical models where they are associated with possible UHECR production sites, such as AGN , hot spots of powerful radio-galaxies , dead quasars and gamma-ray bursts (GRB) . These models have much in common. They assume that primary particles are protons; the sources of the observed UHECR have, therefore, to lie within the GZK cutoff distance. For energies $`E10^{20}`$ eV the GZK radius is $`R_{\mathrm{GZK}}50`$ Mpc, while at $`E>2\times 10^{20}`$ eV it drops to $`20`$ Mpc . In all these models the distribution of sources in space within the GZK sphere is close to uniform, while the distribution in luminosity does not depend on space and peaks around a certain value.
An important common feature of these models is a small local density of sources. The number density of dead quasars is estimated as $`10^4`$ Mpc<sup>-3</sup> ; the number of AGN is $`10\%`$ of the number of galaxies , which gives $`5\times 10^4`$ Mpc<sup>-3</sup>. Most likely, only a small fraction of them is capable of producing UHECR with energies $`E>10^{20}`$ eV. In the case of GRB the effective density of sources is determined by the rate $`\gamma `$ of GRB and the typical time delay $`\tau `$ of UHECR particles. Taking $`\tau <10^5`$ yr and the rate $`\gamma 2\times 10^{10}h^3`$ Mpc<sup>-3</sup> yr<sup>-1</sup> gives the density of sources $`10^5`$ Mpc<sup>-3</sup>.
The purpose of this letter is to show that the observed clustering favors larger density of sources, provided the propagation of UHECR with energy $`E>10^{20}`$ eV is not strongly affected by extra-galactic magnetic fields. The latter assumption is justified if the existing bound on extra-galactic magnetic field $`B<10^9`$ G is valid.
## II Statistics of clustering
The observable quantities which characterize clustering are $`\overline{N}_m`$, the expected numbers of clusters of different multiplicities $`m`$ (e.g., $`\overline{N}_1`$ and $`\overline{N}_2`$ are the expected numbers of single and double events, respectively). They depend on the total exposure of the experiment $`B`$ and the distribution of sources in the flux they produce<sup>*</sup><sup>*</sup>*Here and below we mean the integral flux of cosmic rays with energies above some energy threshold. It measures the average number of events per unit time per unit area of the detector., $`n(F)`$, which is defined in such a way that the number of sources with the flux from $`F`$ to $`F+dF`$ is $`dS=n(F)dF`$. The events which come from the same source at different times are statistically independent and therefore have the Poisson distribution. Thus, the expected number of clusters is
$$\overline{N}_m=_0^{\mathrm{}}\frac{(FB)^m}{m!}\text{e}^{FB}n(F)𝑑F.$$
(1)
This equation implies that the expected total number of events $`N_{\mathrm{tot}}`$ is
$$\overline{N}_{\mathrm{tot}}=\underset{m}{}m\overline{N}_m=B_0^{\mathrm{}}Fn(F)𝑑F=BF_{\mathrm{tot}},$$
(2)
as it should be. The probability to observe $`k`$ clusters of multiplicity $`m`$ is also given by the Poisson distribution,
$$P_m(k)=\frac{(\overline{N}_m)^k}{k!}\text{e}^{\overline{N}_m}.$$
(3)
Any model of UHECR can be characterized by the distribution of sources in distance and luminosity $`f(r,L)`$ (in the case of anisotropic distribution it should be understood as average over the sphere, $`f(r,L)f(𝐫,L)d\mathrm{\Omega }/4\pi )`$). In order to express $`n(F)`$ and $`\overline{N}_m`$ in terms of the distribution function $`f(r,L)`$, consider the sources at distances from $`r`$ to $`r+dr`$. The number of such sources with luminosities from $`L`$ to $`L+dL`$ is
$$dS=f(r,L)\mathrm{\hspace{0.17em}4}\pi r^2drdL.$$
(4)
Making use of the relation $`F=L/4\pi r^2`$ and integrating over $`r`$ one finds $`dS=n(F)dF`$, where
$$n(F)=(4\pi )^2_0^{\mathrm{}}𝑑rr^4f(r,4\pi r^2F).$$
(5)
Here we have neglected the curvature effects since they are small at distances of order 50 Mpc.
In the case of the astrophysical models, the distribution function $`f(r,L)`$ is uniform in space and depends only on the luminosity. The GZK effect, however, makes distant sources fainter by a factor $`\mathrm{exp}(r/R)`$. This is equivalent to substituting
$$f(r,L)\text{e}^{r/R}h(L\text{e}^{r/R})$$
(6)
in Eq.(5), where $`h(L)`$ is the distribution of sources in the luminosity. The exact value of $`R`$ can be determined by numerical simulations of UHECR propagation with full account of the energy dependence. For protons with $`E>10^{20}`$ eV the simulations give $`R25`$ Mpc .
## III Number of sources
A key parameter which enters the distribution $`f(r,L)`$ is the normalization, or the spatial density of sources, which can be characterized by the number of sources $`S`$ within the sphere of a radius $`R`$. An important information about $`S`$ can be obtained from statistical analysis of clustering even if the functional form of the distribution $`f(r,L)`$ is not known. The idea is to find the distribution $`f(r,L)`$ which corresponds to the minimum number of sources $`S_{}`$ with total number of events, $`\overline{N}_{\mathrm{tot}}`$, and the number of events in clusters, $`\overline{N}_{\mathrm{cl}}\overline{N}_{\mathrm{tot}}\overline{N}_1`$, being fixed. We will show in a moment that in the case $`\overline{N}_{\mathrm{cl}}\overline{N}_{\mathrm{tot}}`$ the number $`S_{}`$ is surprisingly large, much larger than the number of the sources already observed.
It is intuitively clear why in the case $`\overline{N}_{\mathrm{cl}}\overline{N}_{\mathrm{tot}}`$ the number of sources is much larger than $`\overline{N}_{\mathrm{tot}}`$ . In order to produce $`N_{\mathrm{tot}}`$ single events by $`N_{\mathrm{tot}}`$ sources each of them has to be bright enough. But then a large number of doublets would be produced as well. Since this is not the case, i.e. most of the resolved sources are dim and produce at most one event, one concludes that there is a large number of sources which have not yet revealed themselves. Assuming that all sources have the same flux $`F`$ one finds from Eq. (1) $`\overline{N}_1S\overline{n}`$ and $`\overline{N}_2S\overline{n}^2/2`$, where $`\overline{n}=FB`$ is the average number of events produced by one source. Therefore, $`S\overline{N}_1^2/(2\overline{N}_2)\overline{N}_{\mathrm{tot}}^2/\overline{N}_{\mathrm{cl}}`$, i.e. much larger than $`\overline{N}_{\mathrm{tot}}`$. Using methods described in the Appendix it is possible to show that the case of equal fluxes corresponds to the absolute minimum of S. However, this distribution is unphysical. Many realistic situations correspond to a homogeneous distribution of sources in space when more distant sources are fainter; consequently, their number has to be even larger than predicted by the above estimate.
In astrophysical models the distribution $`f(r,L)`$ is given by Eq. (6) containing one unknown function $`h(L)`$. The minimum density of sources is determined by minimizing over $`h(L)`$. As is shown in the Appendix, the minimum is reached at the delta-function distribution
$$h(L)=h_{}\delta (LL_{}),$$
(7)
where $`L_{}`$ is the luminosity of the sources and $`h_{}`$ is their spatial density. The unknown parameters $`h_{}`$ and $`L_{}`$ can be related to $`\overline{N}_{\mathrm{tot}}`$ and $`\overline{N}_{\mathrm{cl}}`$ by making use of Eqs.(1) and (2). Introducing the notations
$`S_{}`$ $`=`$ $`(4\pi /3)R^3h_{},`$ (8)
$`\nu _{}`$ $`=`$ $`BL_{}/(4\pi R^2),`$ (9)
where $`\nu _{}`$ is the expected number of events from one source at the distance $`R`$ in the absence of the GZK cutoff, one has the following equations,
$`\overline{N}_{\mathrm{tot}}`$ $`=`$ $`3S_{}\nu _{},`$ (10)
$`\overline{N}_1`$ $`=`$ $`3S_{}\nu _{}{\displaystyle _0^{\mathrm{}}}𝑑x\mathrm{exp}\left(x\nu _{}x^2\text{e}^x\right).`$ (11)
These equations can be solved perturbatively at small $`\overline{N}_{\mathrm{cl}}\overline{N}_{\mathrm{tot}}`$. One finds
$$\nu _{}\frac{1}{\pi }\frac{\overline{N}_{\mathrm{cl}}^2}{\overline{N}_{\mathrm{tot}}^2},$$
(12)
$$S_{}\frac{\pi }{3}\frac{\overline{N}_{\mathrm{tot}}^3}{\overline{N}_{\mathrm{cl}}^2}.$$
(13)
If $`\overline{N}_{\mathrm{cl}}\overline{N}_{\mathrm{tot}}`$, the minimum number of sources $`S_{}`$ is indeed much larger than $`\overline{N}_{\mathrm{tot}}`$ and, therefore, is much larger than the number of sources already observed. From Eq. (12), each source produces much less than 1 event in average.
## IV Discussion
Let us apply these arguments to the observed events with energies $`E>10^{20}`$ eV. In this case, $`N_{\mathrm{tot}}=14`$ and $`N_{\mathrm{cl}}=2`$. The solution to Eqs.(10) and (11) is $`S_{}400`$, which at $`R=25`$ Mpc corresponds to the density
$$h_{}6\times 10^3\mathrm{Mpc}^3.$$
(14)
This number is large as compared to the density of sources in most of astrophysical models. However, it should be interpreted with care. One may expect large statistical fluctuations because both $`N_{\mathrm{tot}}`$ and $`N_{\mathrm{cl}}`$ are small and may not coincide with their expected values.
In order to address this issue quantitatively, let us find the model which has the largest probability $`p(h_{})`$ to reproduce the observed clustering at fixed density of sources $`h_{}`$. To this end, consider the set of models which are described by Eqs.(6) and (7) and are characterized by two parameters $`h_{}`$ and $`L_{}`$. At fixed density of sources, there remains a freedom of changing $`L_{}`$. The probability to reproduce the observed data is maximum for some $`L_{}`$; this probability is $`p(h_{})`$. By construction, there are no models with the density of sources smaller than $`h_{}`$, in which the probability to reproduce the observed data is larger than $`p(h_{})`$.
| proba- | 14 events | | 30 events | | 60 events | |
| --- | --- | --- | --- | --- | --- | --- |
| bility | 1 doublet | | 1 doublet | | 1 doublet | |
| $`p`$ | $`h_{}`$ | $`\nu `$ | $`h_{}`$ | $`\nu `$ | $`h_{}`$ | $`\nu `$ |
| 0.1 | $`23`$ | 0.51 | $`320`$ | 0.065 | $`3100`$ | 0.012 |
| 0.01 | $`3.2`$ | 5.7 | $`63`$ | 0.38 | $`690`$ | 0.058 |
| 0.001 | | | $`21`$ | 1.3 | $`260`$ | 0.15 |
TABLE I. Minimum density of sources $`h_{}`$ in the units of $`10^5`$ Mpc<sup>-3</sup>, and corresponding source luminosity in the units of $`\nu =BL_{}/(4\pi R^2)`$, which are required to reproduce the observed clustering with given probability $`p`$ for the real experimental data (1 doublet out of 14 events) and for two hypothetical data sets with larger number of events (one doublet out of 30 events and one doublet out of 60 events).
There is some ambiguity in defining what is “to reproduce the observed data”. In the case at hand we request that the number of singlets is 12 or larger, the number of doublets is 1 or smaller, and the number of clusters with the multiplicity 3 and larger is zero. Eq. (3) determines the probability $`p(h_{},L_{})`$ of such clustering as a function of two parameters $`h_{}`$ and $`L_{}`$. The probability $`p(h_{})`$ is found by maximizing $`p(h_{},L_{})`$ at fixed $`h_{}`$. We have performed this calculation numerically. The results are summarized in Table 1 in the form of lower bounds on the density of sources. We also present the source luminosity in the units of $`\nu =BL_{}/(4\pi R^2)`$, i.e. the number of events from a single source at the distance $`R`$.
The models where the observed clustering occurs with probability 1% have minimum $`2`$ sources inside the 25 Mpc sphere. In the latter case most of the observed 14 events are produced by the sources which are further than 25 Mpc and thus have to be bright enough. This is reflected in Table 1 from which we see that these sources would produce in average $`6`$ events each (in the absence of the GZK cutoff) if placed at 25 Mpc.
It is worth noting that the numbers of Table 1 correspond to the extreme situation when the distribution of sources is given by Eq. (7) with a particular value of $`L_{}`$. In realistic models, the distribution of sources in luminosity is usually spread over an order of magnitude at least. There may also be constraints on the luminosity of the sources. In these cases, the lower bounds on the number of sources are higher than in Table 1.
When the new large-area detectors like the Pierre Auger array will start operating, the number of observed events will increase and the statistical errors in determination of the density of sources will go down. Correspondingly, the lower bounds on the density will become higher. To show that the bounds may become very high when the total number of events is still small, we have performed calculations for two hypothetical situations, 1 doublet out of 30 events, and 1 doublet out of 60 events. The results are also listed in Table 1. The bounds grow roughly like cube of the number of events, in agreement with Eq.(13).
To summarize, the statistical analysis of clustering may provide tight constraints on astrophysical models of UHECR when the number of clusters is small. In this situation, a key quantity is the density of sources which can be bound from below in a model-independent way. The bound grows very fast with the number of single events above $`E=10^{20}`$ eV and is potentially dangerous for astrophysical models which associate production of UHECR with GRB or exceptional galaxies such as AGN, powerful radio-galaxies and dead quasars.
Our method equally applies to models in which UHECR are produced in the Galactic halo, or in which primary particles are immune to the background radiation. The relation (13) remains valid with a different numerical coefficient of order one and different meaning of $`S_{}`$. In the first case $`S_{}`$ is the number of sources in the halo and detailed analysis shows that statistical properties of clustering of UHECR are compatible with clumpiness of super-heavy dark matter in decays of which UHECR may be produced. In the second case our method counts the number of UHECR sources within the cosmological horizon, which is inaccessible by other means.
## Acknowledgments
The authors are grateful to K.A. Postnov for the references concerning the number of galaxies and to V.A. Rubakov and M.E. Shaposhnikov for valuable comments and discussions. This work is supported in part by the Swiss Science Foundation, grant 21-58947.99. S.D. acknowledges the hospitality of the University of Lausanne and was supported in part by RFBR grant 99-02-18410, by the RAS JRP grant 37, and by ISSEP.
## APPENDIX: minimum number of sources
Consider the problem in general terms. First note that by changing the integration variable in Eq. (5) one can show that any distribution is equivalent to a factorizable one. So, let us take the distribution of sources in the form
$$f(r,L)=g(r)h(L).$$
(15)
Let us fix $`g(r)`$ and minimize the number of sources
$`S=4\pi {\displaystyle _0^{\mathrm{}}}r^2g(r)𝑑r{\displaystyle _0^{\mathrm{}}}h(L)𝑑L`$
with respect to the distribution $`h(L)`$ under the constraints fixing $`\overline{N}_{\mathrm{tot}}`$ and $`\overline{N}_1`$,
$`B{\displaystyle _0^{\mathrm{}}}g(r)𝑑r{\displaystyle _0^{\mathrm{}}}Lh(L)𝑑L`$ $`=`$ $`\overline{N}_{\mathrm{tot}},`$ (16)
$`B{\displaystyle _0^{\mathrm{}}}𝑑r𝑑LLh(L)g(r)\mathrm{exp}\left({\displaystyle \frac{BL}{4\pi r^2}}\right)`$ $`=`$ $`\overline{N}_1.`$ (17)
This is equivalent to minimizing the functional
$`W=4\pi {\displaystyle _0^{\mathrm{}}}dLh(L){\displaystyle _0^{\mathrm{}}}g(r)dr\{r^2+\lambda {\displaystyle \frac{BL}{4\pi }}`$
$$\mu \frac{BL}{4\pi }\mathrm{exp}(\frac{BL}{4\pi r^2})\}\lambda \overline{N}_{\mathrm{tot}}+\mu \overline{N}_1$$
(18)
with respect to $`h(L)`$. Here $`\lambda `$ and $`\mu `$ are the Lagrange multipliers.
The functional (18) is linear in $`h(L)`$; denote the coefficient by $`G(L)`$,
$`G(L)={\displaystyle _0^{\mathrm{}}}g(r)𝑑r\left\{r^2+\lambda {\displaystyle \frac{BL}{4\pi }}\mu {\displaystyle \frac{BL}{4\pi }}\mathrm{exp}\left({\displaystyle \frac{BL}{4\pi r^2}}\right)\right\}.`$
At those values of $`L`$ where $`G(L)`$ is negative, the minimum of $`W`$ is at $`h(L)\mathrm{}`$. The latter, however, is not compatible with Eqs.(16) and (17). Therefore, at the minimum the values of $`\lambda `$ and $`\mu `$ have to be such that $`G(L)`$ is non-negative.
At those values of $`L`$ where $`G(L)`$ is positive, the minimum of $`W`$ is reached at $`h(L)=0`$. If $`G(L)`$ is positive at all $`L`$, then $`h(L)`$ is identically zero and Eqs.(16) and (17) are again violated. Therefore, $`\lambda `$ and $`\mu `$ must be such that $`G(L)`$ touches zero at some $`L_{}`$. The function $`h(L)`$ is non-zero only at this point. Thus, the minimum number of sources corresponds to the situation when all of them have the same luminosity $`L_{}`$, and we arrive at the delta-function distribution, Eq.(7).
It remains to show that for a given positive function $`g(r)`$ satisfying $`g(r)𝑑r<\mathrm{}`$ the Lagrange multipliers $`\lambda `$ and $`\mu `$ can always be chosen in such a way that $`G(L)`$ is positive everywhere except an isolated point. To this end rewrite $`G(L)`$ in the following form,
$$G(L)=C+\lambda F(L),$$
(19)
where $`C=r^2g(r)𝑑r`$ is a positive constant and the function $`F(L)`$ depends only on the ratio $`\mu /\lambda `$,
$`F(L)={\displaystyle \frac{BL}{4\pi }}{\displaystyle _0^{\mathrm{}}}g(r)𝑑r\left\{1{\displaystyle \frac{\mu }{\lambda }}\mathrm{exp}\left({\displaystyle \frac{BL}{4\pi r^2}}\right)\right\}.`$
The behavior of the function $`F(L)`$ is the following. At $`L0`$ it goes to zero. At small $`L`$ it is negative if $`\mu /\lambda >1`$ and positive otherwise. At $`L\mathrm{}`$ it grows linearly with $`L`$, the coefficient being $`B/4\pi g(r)𝑑r>0`$. Therefore, at $`\mu /\lambda >1`$ the function $`F(L)`$ must have an absolute minimum at some $`L_{}>0`$ (which is a function of $`\mu /\lambda `$). Then it is clear from Eq.(19) that by choosing $`\lambda =C/F(L_{})>0`$ one can set $`G(L)`$ to zero in that particular point. The argument can be easily generalized to the case of infinite number of sources, $`g(r)r^2𝑑r=\mathrm{}`$.
In order to apply this argument to the case of astrophysical models, one should find the factorizable distribution $`\stackrel{~}{f}(r,L)`$ which produces the same $`n(F)`$ as Eq.(6). This can be done by substituting Eq.(6) into Eq.(5) and changing the integration variable according to
$$r^2\mathrm{exp}(r/R)=x^2.$$
(20)
The result reads
$`\stackrel{~}{f}(x,L)=g(x)h(L),`$
where
$`g(x)=(1+r(x)/2R)^1\text{e}^{3r(x)/2R}`$
and $`r(x)`$ is defined by Eq.(20). |
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