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warning/0001/math-ph0001010.html | ar5iv | text | # References
Osterwalder–Schrader axioms—Wightman Axioms—The mathematical axiom systems for quantum field theory (QFT) grew out of Hilbert’s sixth problem , that of stating the problems of quantum theory in precise mathematical terms. There have been several competing mathematical systems of axioms, and here we shall deal with those of A.S. Wightman and of K. Osterwalder and R. Schrader , stated in historical order. They are centered around group symmetry, relative to unitary representations of Lie groups in Hilbert space. We also mention how the Osterwalder–Schrader axioms have influenced the theory of unitary representations of groups, making connection with . Wightman’s axioms involve: (1) a unitary representation $`U`$ of $`G:=\mathrm{SL}(2,)^4`$ as a cover of the Poincaré group of relativity, and a vacuum state vector $`\psi _0`$ fixed by the representation, (2) quantum fields $`\phi _1\left(f\right),\mathrm{},\phi _n\left(f\right)`$, say, as operator-valued distributions, $`f`$ running over a specified space of test functions, and the operators $`\phi _i\left(f\right)`$ defined on a dense and invariant domain $`D`$ in $`𝐇`$ (the Hilbert space of quantum states), and $`\psi _0D`$, (3) a transformation law which states that $`U\left(g\right)\phi _j\left(f\right)U\left(g^1\right)`$ is a finite-dimensional representation $`R`$ of the group $`G`$ acting on the fields $`\phi _i\left(f\right)`$, i.e., $`_iR_{ji}\left(g^1\right)\phi _i\left(g\left[f\right]\right)`$, $`g`$ acting on space-time and $`g\left[f\right]\left(x\right)=f\left(g^1x\right)`$, $`x^4`$. (4) The fields $`\phi _j\left(f\right)`$ are assumed to satisfy locality and one of the two canonical commutation relations of $`[A,B]_\pm =AB\pm BA`$, for fermions, resp., bosons; and (5) it is assumed that there is scattering with asymptotic completeness, in the sense $`𝐇=𝐇^{\text{in}}=𝐇^{\text{out}}`$.
The Wightman axioms were the basis for many of the spectacular developments in QFT in the seventies, see, e.g., , and the Osterwalder–Schrader axioms came in response to the dictates of path space measures. The constructive approach involved some variant of the Feynman measure. But the latter has mathematical divergences that can be resolved with an analytic continuation so that the mathematically well-defined Wiener measure becomes instead the basis for the analysis. Two analytical continuations were suggested in this connection: in the mass-parameter, and in the time-parameter, i.e., $`t\sqrt{1}t`$. With the latter, the Newtonian quadratic form on space-time turns into the form of relativity, $`x_1^2+x_2^2+x_3^2t^2`$. We get a stochastic process $`𝐗_t`$: symmetric, i.e., $`𝐗_t𝐗_t`$; stationary, i.e., $`𝐗_{t+s}𝐗_s`$; and Osterwalder–Schrader positive, i.e., $`_\mathrm{\Omega }f_1𝐗_{t_1}f_2𝐗_{t_2}\mathrm{}f_n𝐗_{t_n}𝑑P0`$, $`f_1,\mathrm{},f_n`$ test functions, $`\mathrm{}<t_1t_2\mathrm{}t_n<\mathrm{}`$, and $`P`$ denoting a path space measure.
Specifically: If $`t/2<t_1t_2\mathrm{}t_n<t/2`$, then
(1)
$$\begin{array}{c}\mathrm{\Omega }A_1e^{\left(t_2t_1\right)\widehat{H}}A_2e^{\left(t_3t_2\right)\widehat{H}}A_3\mathrm{}A_n\mathrm{\Omega }\hfill \\ \hfill =\underset{t\mathrm{}}{lim}\underset{k=1}{\overset{n}{}}A_k\left(q\left(t_k\right)\right)d\mu _t\left(q()\right).\end{array}$$
By Minlos’ theorem, there is a measure $`\mu `$ on $`𝒟^{}`$ such that
(2)
$$\underset{t\mathrm{}}{lim}e^{iq\left(f\right)}d\mu _t\left(q\right)=e^{iq\left(f\right)}d\mu \left(q\right)=:S\left(f\right)$$
for all $`f𝒟`$. Since $`\mu `$ is a positive measure, we have
$$\underset{k}{}\underset{l}{}\overline{c}_kc_lS\left(f_k\overline{f}_l\right)0$$
for all $`c_1,\mathrm{},c_n`$, and all $`f_1,\mathrm{},f_n𝒟`$. When combining (1) and (2), we note that this limit-measure $`\mu `$ then accounts for the time-ordered $`n`$-point functions which occur on the left-hand side in formula (1). This observation is further used in the analysis of the stochastic process $`𝐗_t`$, $`𝐗_t\left(q\right)=q\left(t\right)`$. But, more importantly, it can be checked from the construction that we also have the following reflection positivity: Let $`\left(\theta f\right)\left(s\right):=f\left(s\right)`$, $`f𝒟`$, $`s`$, and set
$$𝒟_+=\{f𝒟f\text{ real valued, }f\left(s\right)=0\text{ for }s<0\}.$$
Then
$$\underset{k}{}\underset{l}{}\overline{c}_kc_lS\left(\theta \left(f_k\right)f_l\right)0$$
for all $`c_1,\mathrm{},c_n`$, and all $`f_1,\mathrm{},f_n𝒟_+`$, which is one version of Osterwalder–Schrader positivity.
Since the Killing form of Lie theory may serve as a finite-dimensional metric, the Osterwalder–Schrader idea turned out also to have implications for the theory of unitary representations of Lie groups. In , the authors associate to Riemannian symmetric spaces $`G/K`$ of tube domain type, a duality between complementary series representations of $`G`$ on one side, and highest weight representations of a $`c`$-dual $`G^c`$ on the other side. The duality $`GG^c`$ involves analytic continuation, in a sense which generalizes $`t\sqrt{1}t`$, and the reflection positivity of the Osterwalder–Schrader axiom system. What results is a new Hilbert space where the new representation of $`G^c`$ is “physical” in the sense that there is positive energy and causality, the latter concept being defined from certain cones in the Lie algebra of $`G`$.
A unitary representation $`\pi `$ acting on a Hilbert space $`𝐇(\pi )`$ is said to be reflection symmetric if there is a unitary operator $`J:𝐇(\pi )𝐇(\pi )`$ such that
* $`J^2=\text{id}`$.
* $`J\pi (g)=\pi (\tau (g))J,gG`$,
where $`\tau \mathrm{Aut}\left(G\right)`$, $`\tau ^2=id`$, and $`H:=\{gG\tau \left(g\right)=g\}`$.
A closed convex cone $`C𝔮`$ is hyperbolic if $`C^o\mathrm{}`$ and if $`\mathrm{ad}X`$ is semisimple with real eigenvalues for every $`XC^o`$.
Assume the following for $`(G,\pi ,\tau ,J)`$:
* $`\pi `$ is reflection symmetric with reflection $`J`$.
* There is an $`H`$-invariant hyperbolic cone $`C𝔮`$ such that $`S(C)=H\mathrm{exp}C`$ is a closed semigroup and $`S(C)^o=H\mathrm{exp}C^o`$ is diffeomorphic to $`H\times C^o`$.
* There is a subspace $`0𝐊_0𝐇(\pi )`$ invariant under $`S(C)`$ satisfying the positivity condition
$$v\text{ }v_J:=v\text{ }J(v)0,v𝐊_0.$$
Assume that $`(\pi ,C,𝐇,J)`$ satisfies (PR1)–(PR3). Then the following hold:
* $`S(C)`$ acts via $`s\stackrel{~}{\pi }(s)`$ by contractions on $`𝐊`$ ($`=`$ the Hilbert space obtained by completion of $`𝐊_0`$ in the norm from (PR3)).
* Let $`G^c`$ be the simply connected Lie group with Lie algebra $`𝔤^c`$. Then there exists a unitary representation $`\stackrel{~}{\pi }^c`$ of $`G^c`$ such that $`d\stackrel{~}{\pi }^c(X)=d\stackrel{~}{\pi }(X)`$ for $`X𝔥`$ and $`id\stackrel{~}{\pi }^c(Y)=d\stackrel{~}{\pi }(iY)`$ for $`YC`$, where $`𝔥:=\{X𝔤\tau \left(X\right)=X\}`$.
* The representation $`\stackrel{~}{\pi }^c`$ is irreducible if and only if $`\stackrel{~}{\pi }`$ is irreducible.
Palle E.T. Jorgensen: jorgen@math.uiowa.edu
Gestur Ólafsson: olafsson@math.lsu.edu |
warning/0001/nucl-th0001031.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The success of the shell model in predicting the nuclear magic numbers is related to the presence of a strong spin-orbit term in the nuclear average potential. Few years after the formulation of the nuclear shell–model , evidences of spin-orbit terms in the nuclear interaction were identified by analyzing polarized proton scattering data off complex nuclei .
The strong spin–orbit term of the nuclear average potential should be generated by an analogous term present in the nucleon–nucleon interaction. However, the connection between realistic nucleon–nucleon interactions, i.e. those built to reproduce the two nucleon scattering data and the deuteron properties, and the effective interactions, those used in nuclear structure effective theories, is still unclear.
In the present article we present a study on the relationship between the spin–orbit terms used in realistic nucleon–nucleon interaction and the analogous terms used in Hartree–Fock (HF) calculations. Our work has been done within the non relativistic framework where the spin–orbit terms can be easily isolated. In effect, in our HF calculations, we have used the explicit spin-orbit terms of the Argonne–Urbana nucleon-nucleon realistic potentials describing, within the non relativistic framework, nucleon–nucleon elastic scattering data up to energies of about 300 MeV.
The spin–orbit interactions commonly used in HF calculations are of zero–range type and, in general, they are parametrized following the expressions proposed by Skyrme . Also the Gogny interaction , which has a finite range for all the other channels, uses a Skyrme–like expression for the spin–orbit term.
We explicitly develop the expressions for a finite range spin–orbit interaction to be used in HF calculations. As expected, in addition to the direct term these expressions produce a contribution also in the Fock–Dirac exchange term of the HF equations. This term is not present when zero–range interactions are used. Also in the direct term there are some new contributions with respect to the expression obtained with the Skyrme interaction.
In the next paragraphs we present the detailed expressions of the HF equations when finite range spin–orbit terms are considered, then we discuss the importance of the finite–range and the effect of using spin–orbit terms taken from realistic nucleon–nucleon interactions. Finally we draw our conclusions.
## 2 The formalism
In ref. we made an explicit presentation of the HF formalism with finite–range interactions. In the present article we extend the formalism in order to treat also the spin–orbit terms. For this reason we recall here only those parts of the formalism involving the spin-orbit terms.
The effective interaction used in our calculations has the form:
$$V(𝐫_1,𝐫_2)=\underset{p=1}{\overset{8}{}}V_p(𝐫_1,𝐫_2)O^p(1,2)$$
(1)
where to the first 6 components used in ref. we have added the spin–orbit terms defined as $`O^7(1,2)=𝐋𝐒`$ and $`O^8(1,2)=𝐋𝐒𝝉_1𝝉_2`$ with $`𝐋=(𝐫_1𝐫_2)\times (𝐩_1𝐩_2)`$ and $`𝐒=\frac{1}{2}(𝝈_1+𝝈_2)`$.
To take advantage of the spherical symmetry of the problem, we describe the single particle wave functions by separating the angular and the radial parts $`\varphi _k(𝐫)\stackrel{~}{k}(\mathrm{\Omega })u_k(r)/r`$ where the subindex $`k`$ indicates all the quantum numbers necessary to identify the state and $`\mathrm{\Omega }`$ the two angular coordinates $`\theta `$ and $`\varphi `$.
The specific expression of the single particle wave functions allows us to reduce the HF equations into a set of differential equations of the type:
$$\frac{\mathrm{}^2}{2m_k}\left(\frac{\mathrm{d}^2}{\mathrm{d}r^2}\frac{l_k(l_k+1)}{r^2}\right)u_k(r)+U_k(r)u_k(r)W_k(r)=ϵ_ku_k(r),$$
(2)
where we have defined:
$$U_k(r)=\underset{i}{}dr^{}u_i^{}(r^{})d\mathrm{\Omega }d\mathrm{\Omega }^{}\stackrel{~}{k}^{}(\mathrm{\Omega })\stackrel{~}{i}^{}(\mathrm{\Omega }^{})V(|𝐫𝐫^{}|)\stackrel{~}{k}(\mathrm{\Omega })\stackrel{~}{i}(\mathrm{\Omega }^{})u_i(r^{}),$$
(3)
and
$$W_k(r)=\underset{i}{}dr^{}u_i^{}(r^{})d\mathrm{\Omega }d\mathrm{\Omega }^{}\stackrel{~}{k}^{}(\mathrm{\Omega })\stackrel{~}{i}^{}(\mathrm{\Omega }^{})V(|𝐫𝐫^{}|)\stackrel{~}{i}(\mathrm{\Omega })\stackrel{~}{k}(\mathrm{\Omega }^{})u_i(r)u_k(r^{}).$$
(4)
While the interaction depends from the relative distance between two nucleons, the HF equations (2) depend upon the distance of the particles from the origin of the reference system. The implementation of finite range interactions in the HF equations requires the separation of the coordinate variables in the interaction. For the central and tensor channels this separation is done by considering the interaction in coordinate space as Fourier transform of the interaction expressed in momentum space (see ref. for details). For the spin–orbit channels ($`p=7,8`$) we use a different strategy consisting in expanding in multipoles the interaction:
$$V_p(r_{12})=4\pi \underset{LM}{}\frac{1}{\widehat{L}^2}𝒱_L^p(r_1,r_2)Y_{LM}^{}(\widehat{r}_1)Y_{LM}(\widehat{r}_2).$$
(5)
with $`\widehat{L}\sqrt{2L+1}`$. ¿From the previous equation, making use of the orthogonality of the spherical harmonics, we obtain a close expression for the coefficients of the expansion:
$$𝒱_L^p(r_1,r_2)=\frac{\widehat{L}^2}{2}_1^1\mathrm{d}\mathrm{cos}\theta _{12}V_p(r_{12})P_L(\mathrm{cos}\theta _{12}),$$
(6)
In the previous equation we have indicated with $`P_L`$ the Legendre polynomials and with $`\theta _{12}`$ the angle between $`𝐫_1`$ and $`𝐫_2`$.
The details of the calculations of the spin–orbit matrix elements are given in the Appendix. We obtain for the direct term in the HF equations the following result:
$$\left[U_k(r)\right]_{p=7,8}=\mathrm{\hspace{0.17em}2}\pi I_k^pdr^{}r^2\left\{\left[j_k\left(j_k+1\right)l_k\left(l_k+1\right)\frac{3}{4}\right]𝒰_C^p(r,r^{})+𝒰_{LS}^p(r,r^{})\right\},$$
(7)
where
$$I_k^p=\{\begin{array}{cc}1,\hfill & p\text{=7}\hfill \\ 2t_k,\hfill & p\text{=8},\hfill \end{array}$$
(8)
and the potentials $`𝒰_C^p`$ and $`𝒰_{LS}^p`$ are given by:
$$𝒰_C^p(r,r^{})=\left[𝒱_0^p(r,r^{})\frac{1}{3}\frac{r^{}}{r}𝒱_1^p(r,r^{})\right]\mathrm{\Omega }_C^p(r^{}),p=7,8$$
(9)
and
$$𝒰_{LS}^p(r,r^{})=\left[𝒱_0^p(r,r^{})\frac{1}{3}\frac{r}{r^{}}𝒱_1^p(r,r^{})\right]\mathrm{\Omega }_{LS}^p(r^{}),p=7,8.$$
(10)
The function $`\mathrm{\Omega }_C^p(r)`$ used in the previous equations has been defined as:
$$\mathrm{\Omega }_C^p(r)=\{\begin{array}{cc}\rho (r),\hfill & p\text{=7}\hfill \\ & \\ \rho ^\pi (r)\rho ^\nu (r),\hfill & p\text{=8},\hfill \end{array}$$
(11)
where $`\rho ^\pi (r)`$ and $`\rho ^\nu (r)`$ are the proton and neutron densities such as $`\rho ^\pi (r)+\rho ^\nu (r)=\rho (r)`$. The other function used in eq. (10) has been defined as:
$$\mathrm{\Omega }_{LS}^p(r)=\{\begin{array}{cc}\rho _{LS}(r),\hfill & p\text{=7}\hfill \\ & \\ \rho _{LS}^\pi (r)\rho _{LS}^\nu (r),\hfill & p\text{=8},\hfill \end{array}$$
(12)
where $`\rho _{LS}`$ is the nucleon spin-density,
$$\rho _{LS}(r)=\frac{1}{4\pi }\underset{i}{}\left[j_i(j_i+1)l_i(l_i+1)\frac{3}{4}\right]\widehat{j}_i^2\left(\frac{u_i(r)}{r}\right)^2,$$
(13)
and $`\rho _{LS}^\pi (r)`$ and $`\rho _{LS}^\nu (r)`$ are the analogous functions for protons and neutrons respectively.
For the exchange terms of eq. (2) we obtain:
$$\left[W_k(r)\right]_{p=7,8}=\underset{iL}{}I_{ki}^p\underset{\alpha =1,5}{}\epsilon _{kiL}^{(\alpha )}dr^{}𝒲_{kiL}^{p(\alpha )}(r,r^{}),$$
(14)
where
$$I_{ki}^p=\{\begin{array}{cc}\delta _{t_k,t_i},\hfill & p=7\hfill \\ & \\ 2\delta _{t_k,t_i}+\delta _{t_k,t_i},\hfill & p=8.\hfill \end{array}$$
(15)
The new five functions $`\epsilon `$ have been defined as:
$$\epsilon _{kiL}^{(\alpha )}=\sqrt{3}(1)^{j_i+l_i+\frac{1}{2}}\widehat{l}_k\widehat{l}_i\widehat{j}_i^2\underset{K}{}(1)^K\widehat{K}^2\left(\begin{array}{ccc}1& K& L\\ 1& 1& 0\end{array}\right)\zeta _{ki}^{(\alpha )}(L,K),\alpha =1,\mathrm{},5,$$
(16)
with $`\zeta _{ki}^{(\alpha )}(L,K)`$ given by:
$`\zeta _{ki}^{(1)}(L,K)`$ $`=`$ $`\xi (l_k+l_i+L)𝒯_{ki}(L,K)`$
$`\left[\sqrt{l_i(l_i+1)}\left(\begin{array}{ccc}l_i& l_k& K\\ 1& 0& 1\end{array}\right)\sqrt{l_k(l_k+1)}\left(\begin{array}{ccc}l_i& l_k& K\\ 0& 1& 1\end{array}\right)\right]`$
$`\zeta _{ki}^{(\alpha )}(L,K)`$ $`=`$ $`𝒢_{ki}(L,K)\{\begin{array}{cc}(2)\sqrt{l_i(l_i+1)}\left(\begin{array}{ccc}l_i& l_k& K\\ 1& 0& 1\end{array}\right)\left(\begin{array}{ccc}1& K& L\\ 1& 1& 0\end{array}\right),\hfill & \alpha =2\hfill \\ & \\ 2\sqrt{l_k(l_k+1)}\left(\begin{array}{ccc}l_i& l_k& K\\ 0& 1& 1\end{array}\right)\left(\begin{array}{ccc}1& K& L\\ 1& 1& 0\end{array}\right),\hfill & \alpha =3\hfill \\ & \\ (\sqrt{2})\left(\begin{array}{ccc}l_i& l_k& K\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}1& K& L\\ 0& 0& 0\end{array}\right),\hfill & \alpha =4,5.\hfill \end{array}`$
In the previous expressions we have used the following definitions:
$`𝒯_{ki}(L,K)`$ $`=`$ $`\left\{\begin{array}{ccc}l_k& \frac{1}{2}& j_k\\ l_i& \frac{1}{2}& j_i\\ K& 1& L\end{array}\right\}\left(\begin{array}{ccc}j_k& j_i& L\\ \frac{1}{2}& \frac{1}{2}& 0\end{array}\right)`$
$`\widehat{l}_k\widehat{l}_i\left\{\begin{array}{ccc}l_k& \frac{1}{2}& j_k\\ l_i& \frac{1}{2}& j_i\\ L& 1& K\end{array}\right\}\left\{\begin{array}{ccc}l_k& l_i& K\\ j_i& j_k& \frac{1}{2}\end{array}\right\}\left(\begin{array}{ccc}l_k& l_i& L\\ 0& 0& 0\end{array}\right),`$
$`𝒢_{ki}(L,K)`$ $`=`$ $`\xi (l_k+l_i+L+1){\displaystyle \underset{L^{}}{}}\xi (L+L^{}+1)\widehat{L^{}}^2\left(\begin{array}{ccc}1& K& L^{}\\ 1& 1& 0\end{array}\right)𝒯_{ki}(L^{},K).`$
In the eqs. (14) we have used five new potentials:
$`𝒲_{kiL}^{p(\alpha )}(r,r^{})`$ $`=`$ $`u_i^{}(r^{})𝒱_L^p(r,r^{})u_k(r^{})u_i(r)\{\begin{array}{cc}1,\hfill & \alpha =1,\hfill \\ & \\ \frac{r^{}}{r},\hfill & \alpha =2,\hfill \\ & \\ \frac{r}{r^{}},\hfill & \alpha =3,\hfill \end{array}`$
$`𝒲_{kiL}^{p(\alpha )}(r,r^{})`$ $`=`$ $`\{\begin{array}{cc}r^{}u_i^{}(r^{})u_k(r^{})𝒱_L^p(r,r^{})\frac{\text{d}}{\text{d}r}u_i(r),\hfill & \alpha =4,\hfill \\ & \\ ru_k(r^{})\frac{\text{d}}{\text{d}r^{}}[u_i^{}(r^{})𝒱_L^p(r,r^{})]u_i(r),\hfill & \alpha =5.\hfill \end{array}`$
Like in ref. the numerical solution of eq. (2) has been obtained iteratively using the plane wave expansion method of refs. . The center of mass motion has been considered in its simplest approximation, consisting in inserting the nucleon reduced mass in the hamiltonian. The single particle wave functions used to start the iterative procedure have been generated by a Saxon–Woods potential without spin–orbit and Coulomb terms. Therefore the starting wave functions for spin–orbit partners are the same.
## 3 Results
In the same spirit of the work of ref. we are more interested in investigating the validity of the commonly used approximations rather than proposing a new effective interaction to be used in HF calculations. This study has been conducted by adding different kinds of spin–orbit terms to a basic interaction composed by the four central terms of the force. These terms are described as a sum of two gaussians:
$$V_p(r)=\underset{i=1}{\overset{2}{}}A_{pi}exp(b_ir^2),$$
(27)
with $`p=1,2,3,4`$. The parameters of this part of the interaction, which we call $`B1a`$, are compared in tab. 1 with the parameterization $`B1`$ of Brink and Boeker . The small differences are due to the fact that we have considered the Coulomb interaction and therefore we had to readjust the parameters of the force in order to reproduce the binding energy of <sup>4</sup>He. We added to the $`B1a`$ interaction a finite range spin–orbit term of gaussian form:
$$V_7(r_{12})=A_7exp(b_7r_{12}^2).$$
(28)
The finite range effects have been investigated by comparing the results obtained with the above interaction with those produced by adding to the $`B1a`$ force a zero range spin-orbit term of the form:
$$V_7(r_{12})=A_7\delta ^3(𝐫_\mathrm{𝟏}𝐫_\mathrm{𝟐}).$$
(29)
The straightforward insertion of this expression in our formalism gives a contribution exactly equal to zero. The reason of this result can be traced back to the fact that we have developed our expressions using $`𝐋=(𝐫_1𝐫_2)\times (𝐩_1𝐩_2)`$. To get results different from zero for a zero-range spin–orbit interaction we set to zero the quantity $`𝒱_1^p(r,r^{})`$ in eqs. (9) and (10), and after inserting eq. (29) we obtained:
$$𝒰_{C,LS}^7(r,r^{})=\frac{A_7}{2r}\mathrm{\Omega }_{C,LS}^7(r^{})\delta (rr^{}).$$
(30)
The calculations done with this approach are labelled as $`z`$.
In addition to these effective interactions we have also used spin–orbit terms taken from microscopic forces: the Urbana V14 , Argonne V14 , and Argonne V18 potentials. Our study has been restricted to the investigation of the doubly magic nuclei <sup>12</sup>C, <sup>16</sup>O, <sup>40</sup>Ca, <sup>48</sup>Ca and <sup>208</sup>Pb.
The finite range interaction (28) has been used to study the role played by the various terms of the spin–orbit potential. In a first set of calculations, labelled as $`c`$, only the $`𝒰_C`$ term of eq. (7) has been used. This is the only spin–orbit term present in shell–model calculations. In another set of calculations, denoted as $`d`$, we considered the full direct term, and, finally, the results identified with $`so`$ have been obtained with all the spin–orbit terms. These calculations have been done by changing every time the parameters of the force (28) to reproduce the 6.3 MeV splitting between the protons 1p levels in <sup>16</sup>O. The parameter $`b_7`$ was fixed to the arbitrary value of 1.2 fm<sup>-2</sup> and the fit of the splitting was obtained by changing $`A_7`$. The values of $`A_7`$ obtained in this way are, -108.75, -107.50, -97.86 MeV for the $`c`$, $`d`$, and $`so`$ calculations respectively. The three interactions do not differ very much as it is shown in the panel I of the figure. This result indicates that the largest contribution from the spin–orbit force is coming from the $`𝒰_C`$ factor of the direct term. Since the other terms are small we have explored the possibility of avoiding their explicit calculation by simulating their effects with a readjustment of the force parameters. This is the reason why each type of calculation has been done with a different parametrization of the force, each of them reproducing the same empirical quantity.
In tab. 2 we compare the binding energies obtained with our calculations with the experimental ones . In spite of the fact that we handle with a non–linear problem, the spin–orbit terms acts on the binding energies as expected. The main contribution to the binding energy is obtained by the sum of the single particle energies. In nuclei where all the spin–orbit partners are occupied the spin–orbit term lowers the energy of the $`l+1/2`$ level and increases that of the $`l1/2`$ level, in such a way that the contribution to the nuclear binding energy is almost zero. In table 2 this is observed by looking at the values of the energies of <sup>16</sup>O and <sup>40</sup>Ca which are practically the same, independently from the spin-orbit force used. Clearly those nuclei where not all the spin–orbit partners are occupied are sensitive to the spin–orbit force, since the single particle energy of the last occupied level is lowered. The effect is seen in <sup>12</sup>C, <sup>48</sup>Ca and <sup>208</sup>Pb where the binding energy increases, in absolute value, the stronger the spin-orbit force is.
The quantity most sensitive to the spin–orbit interaction is the energy splitting between spin–orbit partners levels. The energy splittings calculated for the various nuclei under investigation with the interactions proposed are compared in tabs. 3 and 4 with the Skyrme III results and with the empirical values . The experimental spectrum is more compressed than the theoretical one. This fact is well known , and it is related to the intrinsic limitations of the HF theory in the description of an interacting many–body system.
As expected, the splittings increase with increasing value of $`l`$. The splittings obtained with the zero range interaction $`z`$ become larger than those obtained with finite range interaction as the mass number of the nucleus increases. The results obtained with zero–range Skyrme interaction do not present this effect. In the Skyrme interaction there are velocity dependent terms generating spin–orbit like contributions which add to those produced by the genuine spin–orbit term. These velocity dependent terms simulate the effects of the finite range. We observe that the value of the splittings obtained with the Skyrme III interaction are comparable with those obtained with our finite range interactions.
¿From the comparison of the results of the $`c`$ and $`d`$ columns of tabs. 3 and 4 we infer information on the role of the terms $`𝒰_{LS}^p`$ in eq. (7). The inclusion of $`𝒰_{LS}^p`$ increases the splitting for all the nuclei considered but the magnitude of this increase is rather different for the various nuclei. We should not consider in our analysis the nucleus <sup>16</sup>O since it has been used to fit the interaction. We notice that the addition of $`𝒰_{LS}^p`$ produces quite small differences in the splitting of <sup>40</sup>Ca and <sup>208</sup>Pb nuclei while they are remarkable in <sup>12</sup>C and <sup>48</sup>Ca.
These results can be understood by considering that $`𝒰_{LS}^p`$ is related to the nuclear spin density, eq. (13). If we assume that the radial wave functions $`u(r)`$ are the same for spin–orbit partners levels, the contribution of these two levels to the spin density is exactly zero. In real calculations these wave functions are slightly different, but the contribution to the spin density remains small. This explains the small increase of the splitting in <sup>40</sup>Ca and the relatively large modifications produced in <sup>12</sup>C and <sup>48</sup>Ca. One should remark that only the unoccupied levels contribute to the spin density. For this reason the effect of $`𝒰_{LS}^p`$ is relatively large with respect to that of $`𝒰_C^p`$ in <sup>12</sup>C and <sup>48</sup>Ca where the number of single particle levels is relatively small. In a heavy nucleus like <sup>208</sup>Pb there are many levels contributing in $`𝒰_C^p`$ and the effects of $`𝒰_{LS}^p`$ produced by a single level is relatively small.
The contribution of the exchange term can be seen by comparing the results of the $`d`$ and $`so`$ columns. The variations with respect to the calculations done with only the direct terms can be as big as 10-15%, but not all of them have the same sign. It seems that for all the $`p`$ states the splitting is reduced when the exchange term is considered, but it is increased in the $`f`$, $`g`$ and $`h`$ states. The situation for the $`d`$ states is even more complicated, since the splitting is reduced for the $`1d`$ states in <sup>208</sup>Pb but it has increased for all the other $`d`$ states. The contribution of the exchange term cannot be taken into account in calculations with the direct term only by modifying the force parameters.
The values of the splittings produced by the Urbana (U), and Argonne V<sub>18</sub> interactions are comparable with those of our interactions, while the Argonne V<sub>14</sub> (A14) generates smaller values. The radial dependence of the spin–orbit terms of these interactions are shown in the panel II of the figure. It is remarkable that the results of U and A18 are similar in spite of the large difference in the depth. The depth value of A14 is intermediate between those of the previous two forces, but its splittings are smaller. Th U and A18 forces have similar range, while that of A14 is smaller. These facts indicate that our calculations are more sensitive to the range of the interaction than to its minimum value. In effect we recall that, in our calculations, a zero–range interaction does not produce any splitting.
The calculations done with the microscopic interactions include both spin–orbit and spin–orbit isospin terms. In order to study the importance of the isospin part of the spin–orbit interaction we have repeated each calculation leaving out this terms. The differences of the results obtained with the full interaction and those without the isospin part are very small. In order to avoid a long list of numbers we give in tab. 5, for each nucleus under investigation, the minimum, the maximum and the average difference, in absolute value, between the calculated splittings. It appears clear the relatively small importance of this term of the interaction. This fact can be understood considering that the major contribution to the spin–orbit interaction is coming from the direct $`𝒰_C^p`$ term. For the spin-orbit isospin term of the interaction, the case $`p=8`$, the $`𝒰_C^p`$ term contains a function which is given by the difference between the proton and neutron density distributions, eq. (11). In all the nuclei we have considered this difference is small and particularly small in those nuclei having the same number of protons and neutron (<sup>16</sup>O and <sup>40</sup>Ca). In effect the maximum differences are larger in <sup>48</sup>Ca and <sup>208</sup>Pb than in <sup>16</sup>O and <sup>40</sup>Ca.
## 4 Summary and Conclusions
In this article we have presented a formalism to treat finite range spin–orbit interactions in HF calculations. The finite range of the interaction generates additional terms with respect to the usual shell model expression. One of these is the contribution to the exchange Fock–Dirac term in the HF equation (2). Also in the direct (Hartree) term of this equation there is a new part: the $`𝒰_{LS}`$ piece of eq. (7). The major goal of our work was the investigation of the effects produced by these new components. This has been done by adding different type of spin–orbit terms to a fixed interaction active only in the four central channels. We have used a spin–orbit interaction of a gaussian form whose parameters have been fixed to reproduce the energy splitting of the proton $`1p`$ levels in <sup>16</sup>O.
We have shown that the largest part of the spin–orbit effects in HF calculations is produced by the traditional shell model term, $`𝒰_C`$ in eq. (7). The contribution of the other term, $`𝒰_{LS}`$, is very small and it can be simulated by a redefinition of the parameters of the force. The role of the exchange term is more complicated: its inclusion in the calculations modifies by a maximum of 15% the values of the spin–orbit splittings. The complication arises because these modifications do not have the same sign for all the nuclei studied. In calculations done with only the direct terms, it is not possible to simulate the exchange effects by simply readjusting the parameters of the interaction.
Forcing our formalism to handle zero–range spin–orbit interactions we have studied, by comparison, the importance of the finite range. We found that calculations done with zero–range interaction produce energy splittings which, in heavy nuclei, are much larger than the empirical ones. Traditional HF calculations use spin–orbit zero–range terms of Skyrme type . These expressions produce contributions to the hamiltonian which are related to the derivative of the density distribution, while our expressions depend directly from the density distribution. The dependence from the derivative of the density distribution simulate effects produced by the finite range of the force.
We have done calculations with spin-orbit terms taken from microscopic interactions and we have obtained splittings close to those produced by our effective spin–orbit interactions. This would indicate that the medium does not affect the spin-orbit term of the realistic interaction, in agreement with the findings of G-matrix calculations . We would like to point out, however, that the observables we have investigated are more sensitive to the global properties of the spin-orbit potential than to its details. Modifications of the local properties of the interaction would not produce effects on our results.
We have also investigated the effects of the spin–orbit isospin term of the interaction, and we found them to be very small. These terms are related to the differences between protons and neutrons density and spin–orbit density distributions. In our calculations these quantities are rather small even for a nucleus with a large neutron excess like <sup>208</sup>Pb. There are however indications for observables which seems to be sensitive to this part of the potential .
The comparison of the results of our calculations with the empirical values of the splittings on the various nuclei investigated is not satisfactory. The empirical splittings are smaller than those we have obtained, except for <sup>12</sup>C. This is a well known problem of the HF theory, and it could be solved only by using theories going beyond the mean field description of the nucleus.
## 5 Appendix
Since we have developed the HF equations in spherical coordinates it is necessary to express the operator $`𝐋𝐒`$ in terms of these coordinates. For this purpose we define an operator $`O(ijk)`$ as:
$$O(ijk)(1)^{i+j}𝐫_i\times 𝐩_j𝐬_k=\sqrt{\frac{2\pi }{3}}(1)^{i+j}r_i\underset{\mu }{}(1)^{1\mu }Y_{1\mu }(\widehat{r}_i)\left[_j\sigma (k)\right]_\mu ^1,$$
(31)
where we have set $`\mathrm{}=1`$ and the indexes $`i,j,k`$ can assume only two values, $`1`$ and $`2`$ for example. The previous formula has been obtained by expressing $`r`$ in terms of spherical harmonics and by making explicit use of the tensor product properties. Using the above operator we can write the spin–orbit channels of the force as:
$`V_p(r_{12})O_p(1,2)`$ $``$ $`V_p(r_{12})^p{\displaystyle \underset{i,j,k=1,2}{}}O(ijk)`$ (32)
$`=`$ $`4\pi \sqrt{{\displaystyle \frac{2\pi }{3}}}^p{\displaystyle \underset{L}{}}{\displaystyle \frac{1}{\widehat{L}^2}}𝒱_L^p(r_1,r_2){\displaystyle \underset{jk}{}}(1)^j{\displaystyle \underset{\mu }{}}(1)^{1\mu }\left[_j\sigma (k)\right]_\mu ^1`$
$`{\displaystyle \underset{iM}{}}(1)^{M+i}r_iY_{LM}(\widehat{r}_1)Y_{LM}(\widehat{r}_2)Y_{1\mu }(\widehat{r}_i),p=7,8.`$
The operator $`^p`$ has been defined as:
$$^p=\{\begin{array}{cc}1,\hfill & p=7\hfill \\ & \\ 𝝉(1)𝝉(2),\hfill & p=8.\hfill \end{array}$$
(33)
In eq. (32) we make the sum on $`M`$ e $`\mu `$, and we obtain a more synthetic expression:
$`V_p(r_{12})O_p(1,2)`$ $`=`$ $`2\sqrt{2\pi }^p{\displaystyle \underset{LL^{}}{}}{\displaystyle \frac{\widehat{L^{}}}{\widehat{L}}}\left(\begin{array}{ccc}L& L^{}& 1\\ 0& 0& 0\end{array}\right)𝒱_L^p(r_1,r_2)`$ (37)
$`{\displaystyle \underset{ijk}{}}O_{00}^{LL^{}}(ijk),p=7,8,`$
where we have defined new set of operators as:
$$O_{00}^{LL^{}}(ijk)=(1)^{i+j}r_i\left[[Y_L^{}(\widehat{r}_i)Y_L(\widehat{r}_{i/})]^1\left[_j\sigma (k)\right]^1\right]_0^0$$
(38)
and where we have defined:
$$𝐫_{i/}=\{\begin{array}{cc}𝐫_1,\hfill & \text{for }i\text{=2}\hfill \\ 𝐫_2,\hfill & \text{for }i\text{=1}.\hfill \end{array}$$
(39)
It is convenient to express these operators such as the coordinates of each particle are separated:
$$O_{00}^{LL^{}}(ijk)=\sqrt{3}\underset{KM}{}(1)^{L+M}\left\{\begin{array}{ccc}L& L^{}& 1\\ 1& 1& K\end{array}\right\}\stackrel{~}{O}_{LL^{}K}^M(ijk),$$
(40)
where we have used the Racah 6-$`j`$ symbol and we have defined the operators:
$$\stackrel{~}{O}_{LL^{}K}^M(ijk)=e(ijk)\overline{A}_{JM}^{\left(ijk\right)}(1)\overline{B}_{JM}^{\left(ijk\right)}(2).$$
(41)
The expressions of the three terms of the above equation, for each value of $`(ijk)`$ are given in tab. 6 as functions of the following operators:
$`𝒞_{\lambda \mu }^L(i)`$ $`=`$ $`\left[Y_L(\widehat{r}_i)_i\right]_\mu ^\lambda `$
$`𝒮_{\lambda \mu }^{LK}(i)`$ $`=`$ $`\left[\left[Y_L(\widehat{r}_i)_i\right]^K\sigma (i)\right]_\mu ^\lambda \left[𝒞_K^L(i)\sigma (i)\right]_\mu ^\lambda `$ (42)
$`_{\lambda \mu }^L(i)`$ $`=`$ $`\left[Y_L(\widehat{r}_i)\sigma (i)\right]_\mu ^\lambda .`$
At the end we express the spin–orbit terms of the interaction as:
$`V_p(r_{12})O_p(1,2)`$ $`=`$ $`4\pi ^p{\displaystyle \underset{LL^{}K}{}}f(L,L^{},K)𝒱_L^p(r_1,r_2)`$
$`{\displaystyle \underset{ijk}{}}{\displaystyle \underset{M}{}}(1)^M\stackrel{~}{O}_{LL^{}K}^M(ijk),p=7,8,`$
where $`f`$ is given by:
$$f(L,L^{},K)=()^{L+K}\xi (L+L^{}+1)\frac{\widehat{L^{}}}{\widehat{L}}\left(\begin{array}{ccc}1& K& L^{}\\ 1& 1& 0\end{array}\right)\left(\begin{array}{ccc}1& K& L\\ 1& 1& 0\end{array}\right).$$
(43)
In the calculation of the HF equations for the spin–orbit channels we have used the results corresponding to the matrix elements $`\stackrel{~}{O}_{LL^{}K}^M(ijk)`$ which in the tab. 6 are shown to be function of $`𝒞_{LM}^K`$, $`_{LM}^K`$, $`𝒮_{LM}^{L^{}K}`$ defined in (5) and of the spherical harmonics $`Y_{LM}(\widehat{r})`$.
Using the function $`\xi (l)`$ =1 if $`l`$ is even and =0 if $`l`$ is odd, we express the reduced matrix elements for the spherical harmonics as:
$$l\frac{1}{2}jY_Ll^{}\frac{1}{2}j^{}=\frac{1}{\sqrt{4\pi }}(1)^{j^{}+L+\frac{3}{2}}\widehat{j}\widehat{j^{}}\widehat{L}\xi (l+l^{}+L)\left(\begin{array}{ccc}j& j^{}& L\\ \frac{1}{2}& \frac{1}{2}& 0\end{array}\right).$$
(44)
For the other three operators we obtain the following operators:
$`l{\displaystyle \frac{1}{2}}j𝒮_L^{L^{}K}l^{}{\displaystyle \frac{1}{2}}j^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2\pi }}}(1)^l^{}\widehat{j}\widehat{j^{}}\widehat{l}\widehat{l^{}}\widehat{L}\widehat{L^{}}\widehat{K}\left\{\begin{array}{ccc}l& \frac{1}{2}& j\\ l^{}& \frac{1}{2}& j^{}\\ K& 1& L\end{array}\right\}`$ (58)
$`\{\sqrt{2}\xi (l+l^{}+L^{}+1)\left(\begin{array}{ccc}l^{}& K& l\\ 1& 1& 0\end{array}\right)\left(\begin{array}{ccc}K& 1& L^{}\\ 1& 1& 0\end{array}\right){\displaystyle \frac{\sqrt{l^{}(l^{}+1)}}{r}}`$
$`+\left(\begin{array}{ccc}l^{}& K& l\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}K& 1& L^{}\\ 0& 0& 0\end{array}\right){\displaystyle \frac{\mathrm{d}}{\mathrm{d}r}}\},`$
$`l{\displaystyle \frac{1}{2}}j𝒞_L^Kl^{}{\displaystyle \frac{1}{2}}j^{}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{4\pi }}}(1)^{l+l^{}+\frac{1}{2}+j^{}+L}\widehat{j}\widehat{j^{}}\widehat{l}\widehat{l^{}}\widehat{L}\widehat{K}\left\{\begin{array}{ccc}l& j& \frac{1}{2}\\ j^{}& l^{}& L\end{array}\right\}`$ (71)
$`\{\sqrt{2}\xi (l+l^{}+K+1)\left(\begin{array}{ccc}l^{}& L& l\\ 1& 1& 0\end{array}\right)\left(\begin{array}{ccc}L& 1& K\\ 1& 1& 0\end{array}\right){\displaystyle \frac{\sqrt{l^{}(l^{}+1)}}{r}}`$
$`+\left(\begin{array}{ccc}l^{}& L& l\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}L& 1& K\\ 0& 0& 0\end{array}\right){\displaystyle \frac{\mathrm{d}}{\mathrm{d}r}}\},`$
$`l{\displaystyle \frac{1}{2}}j_L^Kl^{}{\displaystyle \frac{1}{2}}j^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2\pi }}}(1)^l\widehat{j}\widehat{j^{}}\widehat{l}\widehat{l^{}}\widehat{L}\widehat{K}\left\{\begin{array}{ccc}l& \frac{1}{2}& j\\ l^{}& \frac{1}{2}& j^{}\\ K& 1& L\end{array}\right\}\left(\begin{array}{ccc}l& K& l^{}\\ 0& 0& 0\end{array}\right).`$ (77) |
warning/0001/quant-ph0001018.html | ar5iv | text | # Quantum computation with mesoscopic superposition states
## I Introduction
Quantum Mechanics is now fundamental to the modern world we live and interact with, not being just the abstract realm of theoretical physics. Many new areas of emerging technology depend on the principles contained within it . One of the most striking features of quantum systems are superposition states. They have given rise to a large amount of discussion in the literature and now play a central role for the recent developments made in the area of quantum information. This is due to their possibility to encode information in a way impossible to be attained by any classical system. Quantum computation has become a significant subject within quantum information theory, due to the powerful property of superposition states to execute large parallel processing. Quantum information research has also improved significantly the understanding about the quantum systems involved on the factual realisation of a quantum computer and has raised many interesting problems such as in the encoding of information , entanglement of states and quantum cryptography .
A number of core technologies are currently under investigation for constructing a quantum computer which is necessary to fully implement quantum algorithms. These include ion-traps, cavity QED, solid state and liquid state NMR to name but a few. The proposals to engineer a quantum computer or as a first step a single logic gate in the realm of quantum optics are generally based on discrete atomic states and cavity field number states of zero and one photons. A central proposal which has gained much attention in recent years is the Cirac and Zoller trapped ions scheme to encode a n-conditional gate. We also cite the proposals of Sleator and Weinfurter and of Domokos et al. based on cavity-QED (quantum electrodynamics) technology and dealing with two-bit universal gates. Experimentally, there are few initiatives for logical operations in ion traps and in NMR , which allow for a scalable implementation. These proposals require a technological domain, which to date has not been attained . In cavity-QED technology, for optical frequencies, a conditional interaction between two-modes, the idler and pump, have been proposed to encode a phase gate (P-gate) due to the high non-linearity that can be presented by single atoms. At microwave frequencies, logical elements have been demonstrated experimentally as a means of encoding a quantum memory with a single photon .
In this article we are going to focus on cavity QED and the technology associated with it. Cavity QED has had a very rich past and has been instrumental in a huge amount of fundamental quantum and atom optical research . Such a system has been used for photon number quantum non demolition measurements , generation of single Fock state and generation and measurement of the time of decoherence of Schrödinger cats states . With such a rich history recent attention in cavity QED has been focused on quantum information. With the non demolition measurement of a single photon number in the cavity , the technology became available to encode qubits and realise a quantum gate . The quantum information proposals based on cavity QED technology makes use of only zero and one field number states. More recently there is the significant evidence of generation of trapped states of more than one photon which could be used in an encoding scheme.
With a CNot gate based on an encoding scheme using zero and one Fock states, spontaneous errors have a disastrous effect. Quantum information is irreversibly lost. It is possible to protect the system against such errors. In fact to protect the qubit against general one qubit errors it is necessary to encode the original state by distributing its quantum information over at least five qubits. Basically the 5-qubit quantum circuit takes the initial state with four extra qubits in the state $`|0`$ to an encoded state. This state is then protected versus all single qubit errors. Decoding this state and then applying a simple unitary transformation yields the original state. Implementing a five qubit error correcting code is quite expensive in terms of quantum resources. Other encoding schemes may allow simpler error correction circuits.
There is no fundamental reason to restrict oneself to physical systems with two dimensional Hilbert spaces for the encoding. It may be more natural in some contexts to encode logical states as a superposition over a large number of basis states. Significant advances can be achieved. For instance in the protection against errors incoming due to the coupling of the qubit system to a dissipative environment. Recent work by Cochrane et. al. have proposed how macroscopically distinct quantum superposition states (Schrödinger cat states) may be used as logical qubit encoding. Spontaneous emission causes a bit-flip error in these superposition state qubit encoding, which is easily corrected by a standard 3-qubit error correction circuit (compared to five qubits for Fock states). This is particularly relevant, as the bit-flip error is much easier to fix than spontaneous emission errors in Fock state systems. Another good reason for using superposition of coherent states to encode qubits, is that they are naturally generated in any cavity system, while number states of more than one photon require a large amount of control .
In this paper we propose how even and odd mesoscopic coherent superposition of states can be used to implement and encode a CNot quantum gate in a realistic superconducting cavity-QED system, where those states were already generated . We define the even cat state as the $`0`$ qubit and the odd cat state as the $`1`$ qubit. This encoding can be represented as
$`|0_L`$ $``$ $`{\displaystyle \frac{1}{N_+}}\left(|\alpha +|\alpha \right),`$ (1)
$`|1_L`$ $``$ $`{\displaystyle \frac{1}{N_{}}}\left(|\alpha |\alpha \right).`$ (2)
where $`N_\pm =\sqrt{2(1\pm e^{2\left|\alpha \right|^2})}`$. This normalisation is important and will be retained throughout the paper.
Given the generation of the two logic qubits how does one implement a quantum gate in cavity QED. Essentially any two-bit quantum gate is universal . One of these universal quantum gates is the control not gate (CNot) and consist of a conditional gate - here if the control bit is 0 the target bit will be maintained, but if the control bit is 1 the target bit will suffer a flip transform to 0. The CNot gate can be engineered by two Hadamard transforms plus a phase (P) transform . The Hadamard transform is a single qubit operation that leads to a rotation in the state while the P-transform is a conditional two-bits transform necessary to identify the state of the control bit. The question posed here is how to identify these Hadamard and P transforms in a realizable physical cavity QED system when the encoding for the qubits is in terms of odd and even cat states.
To begin this paper we show how the apparatus similar to the one used to generate Schrödinger cat field states can be generalised to perform a CNot gate conditional transform involving two levels of a Rydberg atom and the field mesoscopic superposition state. Here the two levels of a Rydberg atom are considered to encode the controlled (or target) bit and the field cat state will be the control bit. Since the generation of Schrödinger even and odd cat field states in cavity QED experiments is dependent of a conditional measurement , giving a random outcome, we propose in Section III a strategy based on resonant atomic feedback which allow us to definitely prepare the state of the control bit. The essence of this proposal involves using a feedback scheme based on the injection of appropriately prepared atoms. Basically the state of the cavity is monitored indirectly via the detection of atoms that have interacted dispersively with it. If the cavity field state is not in the required state, a photon is injected into the cavity. Finally in the last section of this paper we present a reasonable detailed discussion of dissipation and their effect on the CNot gate. We explicitly discuss the advantages of encoding with superposition states over zero one photon number states used in previous proposals . Attention is focused on the decoherence phenomenon, as this is one of the main difficulties for quantum computation.
## II Superposition State Encoding
In the last few years a great amount of experimental progress in cavity QED has enabled work at the level of single atoms and single photons, where only two electronic energy states of Rydberg atoms participate in the exchange of a photon with the cavity . This has enabled cavity QED technology to be responsible for a large number of interesting experiments showing, the generation of mesoscopic coherent superposition field states, called Schrödinger cat states , the decoherence phenomenon and non-local entanglement of quantum systems . These systems have gained much attention due to the quantum non demolition (QND) property of measurement on the field photon number by atomic interferometry .
Our experimental proposal is based on the cavity-QED scheme to generate the field superposition states and is depicted schematically in Fig.(1). It consists of a Rydberg atom beam crossing three cavities, R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$, C and R$`{}_{}{}^{\theta }{}_{2}{}^{}`$. Here R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$ and R$`{}_{}{}^{\theta }{}_{2}{}^{}`$ are Ramsey zones and C is a superconducting Fabry-Perot cavity of high quality factor. To achieve our desired encoding the atoms are initially prepared at B in circular states of quantum principal number of the order of 50. Such atoms are well suited for this scheme since their lifetime is over $`3\times 10^2`$s .
The R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$ and R$`{}_{}{}^{\theta }{}_{2}{}^{}`$ cavities, where classical fields resonant with an atomic $`|g`$ $``$ $`|e`$ transition (51.099 GHz) are injected during the time of interaction with the atoms, constitute the usual setup for Ramsey interferometry . There, for a selected atomic velocity, the state of the atom will suffer a rotation in the vector space spanned by $`\{|e,|g\}`$.
The experiment is started when one selects the initial state of an atom prepared in the $`|g`$ or $`|e`$ by the laser field L. This atom has a resonant interaction with the field in R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$ given by
$$H_I=\mathrm{}\mathrm{\Omega }\left(a_r\sigma ^++a_r^{}\sigma ^{}\right)$$
(3)
where $`\sigma ^+|eg|`$ and $`\sigma ^{}|ge|`$ are the atomic pseudo-spin Pauli operators, $`a_r^{}`$ ($`a_r`$) are the creation (annihilation) operator for the mode of the field in R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$ and $`\mathrm{\Omega }`$ is the one photon Rabi frequency. With a proper choice of the field phase $`\varphi `$ in R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$ the atomic states $`|g`$ and $`|e`$ are rotated to
$`|g`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|g+e^{i\varphi }|e\right)`$ (4)
$`|e`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left(|ee^{i\varphi }|g\right)`$ (5)
The cavity C is tuned near the resonance of the transitions between the atomic states $`|e`$ and $`|i`$, a reference state corresponding to the higher level from $`|e`$. The frequency of the transition $`|e`$ $``$$`|i`$ is 48.18 GHz and is distinct of any transition involving the level $`|g`$. The mode geometry inside the cavity is configured in such a way that the intensity of the field rises and decreases smoothly through with the atomic trajectory inside C. For sufficiently slow atoms and for sufficiently large cavity mode detuning from the $`|e`$ $``$$`|i`$ frequency transition, the atom-field evolution is adiabatic and no photonic absorption or emission occurs . On the other hand, dispersive effects emerge \- an atom in the state $`|e`$ crossing C induces a phase shift in the cavity field which can be adjusted by a proper selection of the atomic velocity ($`100`$ m/s) . For a $`\pi `$ phase shift the coherent field $`|\alpha `$ in C transforms to $`|\alpha `$. On the other hand, the phase shift caused by an atom in the $`|g`$ state is null. The atom field interaction can be written effectively as
$$H_{off}=\mathrm{}\mathrm{\Omega }_2a^{}a\sigma ^+\sigma ^{}$$
(6)
where $`\mathrm{\Omega }_2`$ is the effective Rabi frequency for the interaction of the atom with the field and $`a^{}`$ ($`a`$) is the creation (annihilation) operator for the field in C. After the atomic interaction with the field in C, the atom crosses the second Ramsey zone $`R_2^\theta `$ which introduces a new rotation in the atomic vector space, analogously to Eq. (4) and (5), but for the phase $`\theta `$. The atomic state is detected in D by an ionization zone detector, instantaneously giving the atomic state and the field state in C. This is due to the entanglement of their states. The important point we emphasise here is that the resonant interaction of the Ramsey zones can be used as Hadamard transform since they induce rotations in the vector space of the target bit (atomic state) and the off-resonant interaction between atom and field in C can be used for the P-transform.
We begin the description of the implementation of the CNot gate by specifying that the coherent field state will be responsible for the encoding of the control bit and the atomic states $`|g`$ and $`|e`$ will be the target bits $`|0_T`$ and $`|1_T`$, respectively. The procedures to implement the CNot gate is described as follows. The laser field L prepares the target bit in $`|g`$ or $`|e`$; a one bit Hadamard transform is applied to the target qubit by the first Ramsey zone R$`{}_{}{}^{\varphi }{}_{1}{}^{}`$; then the two-bit P-gate is realized by the off-resonant atom-field interaction in C and the second Hadamard transform is realized by R$`{}_{}{}^{\theta }{}_{2}{}^{}`$. Finally the atom is detected simultaneously specifying the atomic and field states. The effective unitary operator related to the evolution of the atom-field in cavity C entangled state, due to the sequential interaction of the atom with the field in $`R_1^\varphi `$, C and $`R_2^\theta `$ is given by
$$U(\varphi ,\theta )=U_2^\theta \mathrm{exp}\left[i\mu a^{}a\sigma ^+\sigma ^{}\right]U_1^\varphi $$
(7)
where $`U_1^\varphi `$ and $`U_2^\theta `$ are the unitary operators related to the evolution of the joint state in $`R_1^\varphi `$ and $`R_2^\theta `$, respectively. In eqn. (7), $`\mu =\mathrm{\Omega }_2t`$, where $`t`$ is the time interval for the off-resonant interaction. Proceeding through the immediate states generated by the atomic passing through each of the cavities it is easy to show, for $`\varphi =\pi `$ and $`\theta =0`$, the following table
| Input | $`R_1^\varphi `$ | | C | | $`R_2^\theta `$ | Output |
| --- | --- | --- | --- | --- | --- | --- |
| $`|g|0_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|g|e\right)|0_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|g|e\right)|0_L`$ | $``$ | $`|g|0_L`$ |
| $`|e|0_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|e+|g\right)|0_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|e+|g\right)|0_L`$ | $``$ | $`|e|0_L`$ |
| $`|g|1_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|g|e\right)|1_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|e+|g\right)|1_L`$ | $``$ | $`|e|1_L`$ |
| $`|e|1_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|e+|g\right)|1_L`$ | $``$ | $`\frac{1}{\sqrt{2}}\left(|g|e\right)|1_L`$ | $``$ | $`|g|1_L`$ |
which verifies the standard CNot truth-table.
Above we have discussed a setup where the atoms encode the target qubit and the cavity field mode encodes the control qubit. Nevertheless, it is also possible to proceed with atoms responsible by both the control and target qubit. In this second case, the state of the control atom must be transferred to the cavity C and with a proper selection of the cavity state (to what we address to the next section) the procedure for implementing the CNot gate follows as above. After the second atom, which encodes the target qubit interaction in the process described above, a third atom is sent across the system to read the cavity state in a process similar to the scheme already proposed by Sleator and Weinfurter. To envisage a quantum network, i.e., the interconnection of quantum gates, the carriers of qubits between gates can be achieved by atoms transferring the state of one cavity to another , or even by the coupling of these cavities by superconducting wave-guides which can be responsible by an exchange of states between two gates.
It is important for this proposal to include a brief discussion of the realistic parameters. We first note that an atom crosses the cavity in a time of order of 10<sup>-4</sup> s, which is well below the relaxation time of the field inside C (typically of the order of 10<sup>-3</sup>-10<sup>-2</sup>s for Niobium superconducting cavities ) and below the atomic spontaneous emission time of (3$`\times 10^2`$s) . Therefore, the limits considered in that proposal must be far away from the problematic limits found in those experiments.
Our entire proposal for encoding a CNot gate discussed here is reliant on being able to generate the zero ($`|0_L`$) and one ($`|1_L`$) logical states. For this reason we address in Section (III) a strategy for guaranteeing the exact choice of the initial cavity field state. Without such a strategy, the logical states can only theoretically be generated with a $`50\%`$ probability. More explicitly there is a $`50\%`$ probability that the $`|0_L`$ state actually contains only even photon number states and a $`50\%`$ probability that it contains only odd photon number states.
## III Initial Conditions for the Control Bit
Our generation of the CNot gate outlined in the previous section relies on our ability to be able to generate the coherent logical state encoding with a high degree of certainty. The initial state of the control bit (the field state of the cavity) has to be prepared with a probability greater than 50% as usually occurs in the preparation of superposition field states by Rydberg atoms. The state of the field in the cavity is $`|0_L`$ or $`|1_L`$ conditioned by the measurement of the atomic $`|g`$ or $`|e`$ state in the process of generation of superposition states. Such a scheme is analogous to the depicted in fig.(1), however here we have $`\theta =\pi `$ in the second Ramsey zone and for the initial cavity state a coherent one, considering that the atom was prepared in the $`|e`$ state. Let us suppose we are interested in preparing the state $`|0_L`$ for the control bit. If the atomic state $`|e`$ was detected, then our scheme would have failed. For it to succeed we have to apply a process conditioned to the measurement of the atomic $`|e`$ state to guarantee the flip of the cavity field state from $`|1_L`$ to $`|0_L`$. Analogously we have to apply a process conditioned to the measurement of the atomic $`|g`$ state to guarantee the flip of the cavity field state from $`|0_L`$ to $`|1_L`$ if we are interested in prepare the control bit in the $`|1_L`$ state.
First noting the fact that an atom interacting resonantly with the field in C, with a controlled velocity, can exchange a single photon and regarding that a single photon emission by the cavity field causes
$`a|0_L`$ $`=`$ $`\alpha {\displaystyle \frac{N_{}}{N_+}}|1_L\alpha |1_L(\alpha \mathrm{large})`$ (8)
$`a|1_L`$ $``$ $`\alpha |0_L.`$ (9)
We can now formulate an atomic feedback scheme that operates whenever the atomic detector clicks, if we are interested in the control qubit $`|0_L`$ or $`|1_L`$. In fact this process is very similar to the stroboscopic feedback proposed by Vitali et al. for the suppression of decoherence of superposition field states. Of course we do not need a stroboscopic action, but just one event conditioned to the atomic state measurement.
The scheme proposed is depicted schematically in fig.(2), where B<sub>2</sub> is a source of atoms which are tuned in resonance with the field in C by the Stark shift conditioned to the atomic state measurement made in the ionization zones D<sub>e</sub> or D<sub>g</sub>. The resonant atom-field interaction is given by the Hamiltonian
$$H_I=\mathrm{}\mathrm{\Gamma }\left(a\sigma _f^++a^{}\sigma _f^{}\right)$$
(10)
where $`\mathrm{\Gamma }`$ is the coupling constant between the field and atomic variables. Here $`\sigma _f^+`$ and $`\sigma _f^{}`$ are rising and lowering operator for the feedback atom. If the feedback atom is prepared in the state $`|e`$ then the field state is given by
$`\rho _f^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}`$ $`=`$ $`\mathrm{cos}(\mathrm{\Gamma }\tau \sqrt{a^{}a+1})\rho _C^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}\mathrm{cos}(\mathrm{\Gamma }\tau \sqrt{a^{}a+1})+a^{}{\displaystyle \frac{\mathrm{sin}(\mathrm{\Gamma }\tau \sqrt{a^{}a+1})}{\sqrt{a^{}a+1}}}\rho _C^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}{\displaystyle \frac{\mathrm{sin}(\mathrm{\Gamma }\tau \sqrt{a^{}a+1})}{\sqrt{a^{}a+1}}}a.`$ (11)
where $`\rho _C^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}`$ is the density operator associated with the field state in C before the feedback action. Here $`\rho _f^g`$ ($`\rho _f^e`$) explicitly is the ground (excited) state density operator. $`\tau `$ is the time of interaction of the feedback atom with the field.
As a measure of the field state in the cavity a second atom is sent through the setup and again measured in D<sub>g</sub> or D<sub>e</sub> . The conditional probability $`P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T)`$ that the second atom will be detected in the $`|g`$ or $`|e`$ state, at the time $`T`$ after detection of the first atom follows
$`P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{1\pm {\displaystyle \frac{1}{1+\mathrm{cos}\phi \text{e}^{2\left|\alpha \right|^2}}}\left[\text{e}^{2\left|\alpha \right|^2\text{e}^{\gamma T}}+\mathrm{cos}\phi \text{e}^{2\left|\alpha \right|^2\left(1\text{e}^{\gamma T}\right)}\right]\right\},`$ (12)
conditioned to $`\phi =0`$ \[$`\pi `$\] if the first atom is detected in the $`|g_1`$ \[$`|e_1`$\] state and to the signal + \[-\] for the second atom be detected in the $`|g_2`$ \[$`|e_2`$\] state. For the computation of Eq. (12) at time T we have included the relaxation of the field state due to dissipation. Considering a reservoir at zero temperature, this state is now given by
$`\rho _C^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T)`$ $`=`$ $`{\displaystyle \frac{1}{N_\pm ^2}}\{|\alpha \text{e}^{\gamma T/2}><\alpha \text{e}^{\gamma T/2}|+|\alpha \text{e}^{\gamma T/2}><\alpha \text{e}^{\gamma T/2}|`$ (13)
$`\pm `$ $`\text{e}^{2|\alpha |^2(1\text{e}^{\gamma T})}[|\alpha \text{e}^{\gamma T/2}><\alpha \text{e}^{\gamma T/2}|+|\alpha \text{e}^{\gamma T/2}><\alpha \text{e}^{\gamma T/2}|]\}.`$ (14)
where $`\gamma `$ is the relaxation constant of the field. By analyzing Eq.(12) we observe that if the second atom is detected instantaneously after the first one ($`\gamma T1`$) then
$$P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T)=\frac{1}{2}\left[1\pm \mathrm{cos}\phi \right],$$
(15)
again with $`\phi =0`$ $`(\pi )`$. This gives the conditional probability of detection of the first atom in $`|e_1`$ and the second atom in $`|e_2`$ as $`P(e,e)P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}|^{\pi ,}=1`$ and the probability of detection of the first in $`|g_1`$and the second in $`|e_2`$ as $`P(g,e)P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}|^{0,}=0`$. Analogously $`P(g,g)P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}|^{0,+}=1`$ and $`P(e,g)P^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}|^{\pi ,+}=0`$. This is a signature of the measurement for the state in which the cavity field was prepared. It presents our undesired results $`P(g,e)`$ and $`P(e,g)`$ equal to zero, that is no probability of them occurring. However if the feedback loop is taken in to account in the calculation of the probabilities $`P(g,e)`$ and $`P(e,g)`$ then, instead of using $`\rho _C^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}`$ in Eq. (12) we must use $`\rho _f^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T+\tau )`$ from Eq. (11), where $`T^{}=T+\tau `$. Substituting Eq. (14) for the field relaxation into Eq. (11) it follows
$`P_f^{\left(\genfrac{}{}{0pt}{}{g}{e}\right)}(T+\tau )`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left\{1\pm {\displaystyle \frac{2}{N^2}}\text{e}^{\left|\alpha \right|^2}{\displaystyle \underset{m}{}}{\displaystyle \frac{\left(\left|\alpha \right|^2\text{e}^{\gamma T}\right)^m}{m!}}\left[1+(1)^m\mathrm{cos}\phi \text{e}^{2\left|\alpha \right|^2\left(1\text{e}^{\gamma T}\right)}\right]\mathrm{cos}\left(2\mathrm{\Gamma }\tau \sqrt{m+1}\right)\right\},`$ (16)
which accounts for the conditional probability of detection of the second atom in the state $`|g_2`$ \[$`|e_2`$\] if the first atom was detected in the state $`|e_1`$ \[$`|g_1`$\]. In Fig. (3) we show the respective four conditional probability of atomic detection, $`P(e,e)`$ and $`P(g,g)`$ without feedback and $`P_f(e,g)`$ and $`P_f(g,e)`$ considering the feedback loop for some values of $`\mathrm{\Gamma }\tau `$. This shows the feasibility of the feedback to control the initial state of the cavity. The figure is plotted until $`\gamma T=1`$ since there is no reason to consider times longer than this once the decoherence of the state has already taken place. In fact the scale of time to be taken into account in figures 3(c) and 3(d) is $`T^{}=T+\tau `$, the time interval after the detection of the first atom plus the time interval of the feedback atom. In these figures the continuous solid line represents the absence of feedback. As can be seen there is an optimum value for the feedback process at $`\mathrm{\Gamma }\tau =\pi /6`$ which gives a 93% of chance for the cavity field qubit to be prepared in the right state. It also must be noted that an optimal value is possible only when the feedback atom is sent instantaneously after the click of the respective detector. The performance of the setup decreases considerably when a time delay exists, as can be observed in Figs. 3(c) and 3(d) for $`\gamma T^{}0.1`$. The limit of those curves around 0.5 means that the field state already decohered, and so there is 50% of chance again for generation of the $`|0_L`$ or $`|1_L`$ states. For $`\gamma T^{}>1.0`$ (not shown in the figures) the effect of dissipation implies amplitude damping. The field asymptotically tends to be in a vacuum state, and when this occurs is easily shown through Eq. (16) that the second atom will always be detected in the $`|g`$ state. With the feedback it tends to always be detected in $`|e`$ state.
## IV Efficiency and Sources of Error
This section discusses in detail the advantages and disadvantages of encoding qubits in superposition states instead of number states of only one and zero photon and the effect of dissipation on these. As is already well known for cavity QED experiments the dominant source of error that will affect the implementation of quantum logic elements is cavity damping. Since the cavities are not isolated, when the states $`|0_L`$ or $`|1_L`$ are constructed, the presence of dissipative effects will alter the free evolution of the cavity field state introducing amplitude damping as well as coherence loss. The zero temperature master equation describing the bosonic damping is simply
$$\frac{d\rho }{dt}=\frac{\gamma }{2}\left(2a\rho a^{}a^{}a\rho \rho a^{}a\right),$$
(17)
and its solution for any initial state can be written as
$$\rho (t)=\underset{k=0}{\overset{\mathrm{}}{}}\mathrm{{\rm Y}}_k(t)\rho (0)\mathrm{{\rm Y}}_k^{}(t),$$
(18)
where
$$\mathrm{{\rm Y}}_k(t)=\underset{n=k}{\overset{\mathrm{}}{}}\sqrt{\left(\genfrac{}{}{0pt}{}{n}{k}\right)}\left(e^{\gamma t}\right)^{(nk)/2}\left(1e^{\gamma t}\right)^{k/2}|nkn|.$$
(19)
We are interested in the effect of dissipation on the information encoded in the qubits. For that we will consider first a superposition of Schrödinger cats qubits and thereafter a superposition of one and zero photon number states qubit encoding.
The action of a single decay event $`\mathrm{{\rm Y}}_1`$ on the state
$$|\psi _1=E_1|0_L+E_2|1_L,$$
(20)
leads to
$$\mathrm{{\rm Y}}_1|\psi _1=\alpha \left(1e^{\gamma t}\right)^{1/2}e^{|\alpha |^2\left(1e^{\gamma t}\right)/2}\left(E_1\frac{N_{}}{N_+}|\stackrel{~}{1}_L+E_2\frac{N_+}{N_{}}|\stackrel{~}{0}_L\right),$$
(21)
that is, a simple bit-flip occurs. Here $`|\stackrel{~}{0}_L\frac{1}{N_+}\left(|e^{\gamma t/2}\alpha +|e^{\gamma t/2}\alpha \right)`$ and $`|\stackrel{~}{1}_L\frac{1}{N_{}}\left(|e^{\gamma t/2}\alpha |e^{\gamma t/2}\alpha \right)`$, account for the amplitude damping. A simple unitary process will transform Eq. (21) back to Eq. (20), meaning the reversibility of the process. Under a double decay event $`\mathrm{{\rm Y}}_2`$,
$$\mathrm{{\rm Y}}_2|\psi _1=\alpha ^2\left(1e^{\gamma t}\right)e^{|\alpha |^2\left(1e^{\gamma t}\right)/2}\left(E_1|\stackrel{~}{0}_L+E_2|\stackrel{~}{1}_L\right),$$
(22)
which is exactly our initial state but with amplitude damping. This special superposition is invariant under even number of decay events. This fact brings one important information about these states. However a single decay event, $`\mathrm{{\rm Y}}_1`$ on the Fock superposition state
$$|\psi _2=F_1|0+F_2|1,$$
(23)
leads to
$$\mathrm{{\rm Y}}_1|\psi _2=F_2\left(1e^{\gamma t}\right)^{1/2}|0.$$
(24)
No unitary operation can recover (23) indicating the irreversibility of the process. This means that in one photon state information processing schemes, one single photon decay is fatal, since there is no way in which the resulting error can be corrected once it occurs. However for qubits consisting of superpositions of odd and even number states, one decay event cause a bit-flip, which could be, in principle be corrected. So here, we classified two different kind of error arising from dissipation, one impossible to be corrected (called irreversible error) and the other a bit-flip which can be corrected (reversible error) by unitary processes.
There is a number of error correction schemes that protect the quantum information against single errors. As we have mentioned previously a spontaneous emission error for the Schrödinger cat encoding results in a bit-flip. It is well known that such errors can easily be prevented by a 3-qubit error correction circuit (schematically depicted in Fig (4a)). This circuit is reasonably simple and the superposition state it produces is relatively simple. In fact for an arbitrary qubit $`|\psi =E_1|0_L+E_2|1_L`$ the correction circuit generates the encoded superposition state
$`|\psi `$ $`=`$ $`E_1|000+E_2|111.`$ (25)
To protect against arbitrary error requires normally a 5-qubit error correction circuit (schematically depicted in Fig (4b)). For an arbitrary qubit $`|\psi =F_1|0+F_2|1`$, the correction circuit generates the superposition state
$`|\psi `$ $`=`$ $`F_1\left[|00000+|00110+|01001|01111+|10011+|10101+|11010+|11100\right]`$ (26)
$`+`$ $`F_2\left[|00011|00101|01010|01100|10000+|10110+|11001+|11111\right].`$ (27)
This is quite a complicated superposition state to create (as can be seen from the quantum circuit). The 3-qubit correction circuit is much simpler and hence we see the advantage of the Schrödinger cat encoding. Also, while here we are only discussing a single gate, a reasonable quantum computer has to be constituted of many gates. Then, if the above 5-qubit protection circuit has to be implemented, this will become much more expensive in terms of qubits in comparison to the 3-qubit circuit for bit-flip protection. The bit-flip protection scheme saves 2 qubits at each needed qubit in comparison to the 5-qubit protection circuit described above. It does however only protect against a specific type of error. An unavoidable error incoming from dissipation over superposition states is decoherence. Let us consider the general effect of dissipation on the quantum coherent superposition state. At zero temperature the state of the cavity field is described by the density operator
$`\rho _C^\pm (t)`$ $`=`$ $`{\displaystyle \frac{1}{N_\pm ^2}}\{|\alpha \text{e}^{\gamma t/2}><\alpha \text{e}^{\gamma t/2}|+|\alpha \text{e}^{\gamma t/2}><\alpha \text{e}^{\gamma t/2}|\pm \text{e}^{2|\alpha |^2(1\text{e}^{\gamma t})}`$ (29)
$`\times [|\alpha \text{e}^{\gamma t/2}><\alpha \text{e}^{\gamma t/2}|+|\alpha \text{e}^{\gamma t/2}><\alpha \text{e}^{\gamma t/2}|]\}.`$
We see that two characteristic times are involved in this evolution. The first one, the *decoherence time* is the time in which the pure state given by Eq.(7) is turned into a statistic mixture
$$\rho _C(t)\frac{1}{2}\left\{|\alpha ><\alpha |+|\alpha ><\alpha |\right\},$$
(30)
the second is the *damping time* or *relaxation time* of the field $`t_c`$ =$`\gamma ^1`$, the time that the dissipative effect reduces the energy of the field leading it in to a vacuum state.
The decoherence of the field state is characterized by the $`\mathrm{exp}\left[2\left|\alpha \right|^2\left(1\text{e}^{\gamma t}\right)\right]`$ factor, that for short times, $`\gamma t1`$, turns to be $`\mathrm{exp}\left[2\left|\alpha \right|^2\gamma t\right]`$ and the coherence decays with the time $`t_d=\left(2\gamma \left|\alpha \right|^2\right)^1`$. Unfortunately the coherence time depends on inversely on $`\left|\alpha \right|^2`$ and hence the larger the $`\left|\alpha \right|^2`$ the smaller the coherence time. Decoherence constitutes the main obstacle to quantum computation , since the encoding is completely based in the purity of the field state.
The relaxation time of microwave fields in superconducting cavities is of the order of $`10^2`$ s , what means $`t_d10^2\left|\alpha \right|^2`$ s. So all the interactions involved in this proposal must consider this time and more specifically the number of photons as critical quantities. Moreover, the initial information encoded in the superpositions given by Eqs. (20) and (23) also suffer the effect of decoherence, which for $`|\alpha |1`$ occur at the same time for both encoding schemes. Again, decoherence prevention schemes play a crucial role for any quantum information encoding.
One favorable point for the superposition state encoding is that proposals for sustaining the coherence of these field states have already been considered which could be well adapted for our case. For example the stroboscopic feedback proposal of Vitali et al. . This proposal is particularly appropriate here since it guarantee that at each single decay event a feedback atom is sent through the setup restituting the coherence and the state parity. In fact in the authors claim that the coherence is restored but for a slightly different state. For the proposal presented here what is important is not the original superposition of states, but the original parity of the state, if it was originally a superposition of even or odd photon number states. It is important to emphasize the experimentally critical values for the physical elements involved. The time of flight of the atom across the setup (10<sup>-4</sup> s), the relaxation time of the field (10<sup>-3</sup>-10<sup>-2</sup>s for Niobium superconducting cavities) and the atomic spontaneous emission time (3$`\times 10^2`$s) .
## V Conclusion
In conclusion we have presented a feasible scheme to encode the CNot quantum gate, based on a field superposition of states. These states have been already generated in superconducting microwave cavities which constitute a system almost dominated by the current technology . The proposal here to encode the CNot gate based on a superposition of states is less susceptible to irreversible errors due to dissipative effect imposed by the environment than number states . The generation of these kind of states is dependent of a conditional measurement giving a random assignment of initial control bits, which would be useless if no further process is considered. Hence we propose a conditional feedback scheme, which guarantees that the initial control bit is prepared in the required state. Once the amplitude damping of a coherent state (at zero temperature) still constitutes a coherent state the method proposed works until the inevitable effect of decoherence takes place. For that a reset of the qubits must be done after a time of the order of the time of the decoherence or a coherence control scheme must be applied. The reset process is done repeating the process here described.
The state of a logic unit can be transferred to another logic unit (if the time of decoherence is respected), constituting a sort of quantum memory circuit . That can be attained by the proper choice of atomic interactions between atoms and the field in the microwave cavity or even by the direct photonic process of coupling two cavities by a superconducting wave-guide, which permits an exchange of information (exchange of states) between the coupled units . This problem is of fundamental importance for the engineering of quantum networks. The major sources of error here are the loss of coherence of the field state and the control bit-flip due to the dissipative effect. Analysis of these kind of errors on quantum networks constituted by the basic element here described is left for further investigation.
###### Acknowledgements.
MCO thanks the Fundação de Amparo à Pesquisa do Estado de São Paulo (Brazil) for financial support while WJM acknowledges the support of the Australian Research Council. |
warning/0001/nucl-th0001025.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Some time ago a bright phenomenon of pion condensation predicted in attracted the common attention and was widely investigated -. The reason for this prediction was the fact that one of the solutions of the pion dispersion equation in medium (let it be called $`\omega _c`$) turns to zero $`\omega _c^2(k)=0`$ (at some momentum $`k0`$) while the density of medium increases . This meant an appearance of excitations with zero energy in medium. In this case the ground state should be reconstructed in correspondence with phase transition.
One variant of reconstruction was to include the pion condensate (taken in one or another form) into the ground state. This resulted in a prediction for pion condensation. A number of efforts was devoted to the search for this phenomenon. The result of discussions given in the book of A.B.Migdal et. al. was that the pion condensation manifested itself weakly and probably was absent (at least not observed) at normal density. Nevertheless, this phenomenon could influence the equation of state at high densities achieved in heavy ion collisions or neutron stars.
Below, we consider in detail the solutions of pion dispersion equation
$$\omega ^2k^2m_\pi ^2\mathrm{\Pi }(\omega ,k,p_F)=0$$
(1)
in the complex plane of the pion frequency $`\omega `$. In (1) $`\mathrm{\Pi }`$ is the pion polarization operator (pion self-energy) in the matter. The consideration in the complex $`\omega `$-plane allows us to obtain additional information about well-known solutions.
The main content is as follows. We consider excitations with quantum numbers $`0^{}`$ in symmetrical nuclear matter. At the equilibrium density there are three branches of solutions of (1). When the density is larger than critical one ($`\rho >\rho _c`$) the fourth branch appears on the physical sheet.
The equation (1) has logarithmic cuts on the physical sheet of the complex $`\omega `$-plane, determined by the structure of the polarization operator $`\mathrm{\Pi }`$. Different branches of solutions are considered on physical and unphysical sheets of the complex $`\omega `$-plane. With the change of $`k`$ they move from physical to unphysical sheet (and backward) through the cuts. Below, to explain the obtained results, we deal with the case when isobar width in medium is equal to that in vacuum, i.e. 115 MeV. In this case all solutions of (1) and certain cuts move from the real (or imaginary) axis to the complex plane. This helps us to trace the $`k`$-dependence of the solutions. For certain important cases the influence of $`\mathrm{\Gamma }_\mathrm{\Delta }`$ on the behaviour of $`\omega _i(k)`$ is studied.
The well-known branches of solutions (1) are: 1) spin-isospin sound branch $`\omega _s(k)`$; 2) pion branch $`\omega _\pi (k)`$; 3) isobar branch $`\omega _\mathrm{\Delta }(k)`$. At $`0kk_f`$ they are located on the physical sheet (the momentum $`k_f`$ is different for each branch), then with further increase of $`k`$ they move to unphysical sheet across the cuts. For these branches Re($`\omega _i^2)>0`$ everywhere on the physical sheet.
Our investigations show that there is the whole set of solutions at unphysical sheets. With increasing density certain solutions come to the physical sheet. In such a way the fourth branch, $`\omega _c(k)`$, appears on the physical sheet at $`p_F283`$ MeV<sup>1</sup><sup>1</sup>1For our values of the parameters, which are presented below.. Below $`\omega _c(k)`$ is referred as the condensate branch. At $`k=0`$ $`\omega _c(k)`$ is located at the same point as $`\omega _\pi (k)`$; with increasing of $`k`$ the branch goes onto the unphysical sheet. At some $`k=k_1`$ it appears on the physical sheet and at $`k=k_2`$ leaves it for the same unphysical sheet. The values of $`k_1`$ and $`k_2`$ depend on the density. At critical density, $`\rho _c`$, corresponding to $`p_F=283`$ MeV ($`\rho _c1.2\rho _0`$) there is the equality $`k_1=k_2=1.8m_\pi `$. It is $`\omega _c`$ which obeys the inequality Re$`\omega _c^20`$. All branches depend on the isobar width: at $`\mathrm{\Gamma }_\mathrm{\Delta }=0`$ the branch $`\omega _c(k)`$ is pure imaginary and $`\omega _c^20`$. Recall that $`\rho _0`$ is the equilibrium density of the matter, $`\rho _0=2p_{F0}^3/3\pi ^2`$, $`p_{F0}=268`$ MeV.
The paper is organized as follows. In section 2 the particle-hole polarization operator $`\mathrm{\Pi }`$ is considered in the complex $`\omega `$-plane. Then the branches of solutions of the pion dispersion equation (1) are presented on the physical and unphysical sheets. It is shown how the branches of solutions go across the logarithmic cuts in the complex plane. The condensate branch $`\omega _c`$ is shown in details at different densities $`\rho `$ and isobar width $`\mathrm{\Gamma }_\mathrm{\Delta }`$.
## 2 Polarization operator and its singularities
Here we write down the formulae for the polarization operator used. Our expressions, as is shown below, differ in some points from the well-known expressions . Only $`S`$\- and $`P`$-waves of the $`\pi NN`$ ($`\pi N\mathrm{\Delta }`$) interactions are taken into account. In this case $`\mathrm{\Pi }`$ is the sum of scalar, $`\mathrm{\Pi }_S`$, and vector, $`\mathrm{\Pi }_P`$, terms:
$$\mathrm{\Pi }(\omega ,k)=\mathrm{\Pi }_S(\omega ,k)+\mathrm{\Pi }_P(\omega ,k).$$
(2)
The scalar polarization operator $`\mathrm{\Pi }_S`$ is constructed using linear PCAC equation obtained by Gell-Mann–Oakes–Renner (GMOR) for the pion mass squared in the matter:
$$m_\pi ^2=\frac{NM|\overline{q}q|NM(m_u+m_d)}{2f_\pi ^2}.$$
(3)
The value $`\kappa =NM|\overline{q}q|NM`$ is the scalar quark condensate calculated in nuclear matter ; $`m_u,m_d`$ are masses of the current $`u`$\- and $`d`$-quarks; $`f_\pi ^{}`$ is the pion decay constant in medium. Here we put $`f_\pi ^{}=f_\pi =92`$ MeV (more details can be found in ).
The scalar quark condensate $`\kappa `$ can be expanded in a power series with respect to $`\rho `$
$$\kappa =\kappa _0+\rho N|\overline{q}q|N+termswithhigherdegreesof\rho ,$$
(4)
where $`\kappa _0`$ is the value of scalar quark condensate in vacuum, $`\kappa _0=0.03`$ GeV<sup>3</sup>; $`N|\overline{q}q|N`$ is the matrix element for scalar quark condensate in nucleon, $`N|\overline{q}q|N8`$. In this place it is enough to keep the first two terms in (4), this correspond to the gas approximation for $`\kappa `$ . Then we get from (3)
$$m_\pi ^2=m_\pi ^2\rho \frac{N|\overline{q}q|N(m_u+m_d)}{2f_\pi ^2}.$$
(5)
On the other hand, in the dispersion equation (1) $`m_\pi ^2`$ is defined as
$$m_\pi ^2=m_\pi ^2+\mathrm{\Pi }(\omega ,k=0).$$
(6)
Whilst $`\mathrm{\Pi }_P(k=0)=0`$, we get from (5) and (6) the following form for $`\mathrm{\Pi }_S`$:
$$\mathrm{\Pi }_S=\rho \frac{N|\overline{q}q|N(m_u+m_d)}{2f_\pi ^2}.$$
(7)
The P-wave polarization operator $`\mathrm{\Pi }_P`$ can be written following the papers as a sum of nucleon and isobar polarization operators
$$\mathrm{\Pi }_P=\mathrm{\Pi }_N+\mathrm{\Pi }_\mathrm{\Delta }.$$
(8)
Here $`\mathrm{\Pi }_N(\mathrm{\Pi }_\mathrm{\Delta }`$) is equal to the sum of the nucleon-hole and isobar-hole loops without pion in the intermediate states :
$$\mathrm{\Pi }_N=\mathrm{\Pi }_N^0\frac{1+(\gamma _\mathrm{\Delta }\gamma _{\mathrm{\Delta }\mathrm{\Delta }})\mathrm{\Pi }_\mathrm{\Delta }^0/k^2}{E},\mathrm{\Pi }_\mathrm{\Delta }=\mathrm{\Pi }_\mathrm{\Delta }^0\frac{1+(\gamma _\mathrm{\Delta }\gamma _{NN})\mathrm{\Pi }_N^0/k^2}{E},$$
(9)
$$E=1\gamma _{NN}\frac{\mathrm{\Pi }_N^0}{k^2}\gamma _{\mathrm{\Delta }\mathrm{\Delta }}\frac{\mathrm{\Pi }_\mathrm{\Delta }^0}{k^2}+\left(\gamma _{NN}\gamma _{\mathrm{\Delta }\mathrm{\Delta }}\gamma _\mathrm{\Delta }^2\right)\frac{\mathrm{\Pi }_N^0\mathrm{\Pi }_\mathrm{\Delta }^0}{k^4}.$$
The expressions for the nucleon-hole loop , $`\mathrm{\Pi }_N^0`$, is as follows:
$$\mathrm{\Pi }_N^0(\omega ,k)=\mathrm{Sp}\frac{d^3p}{(2\pi )^3}\mathrm{\Gamma }_{\pi NN}^2\left[\frac{\theta (pp_F)\theta (p_F|\stackrel{}{p}+\stackrel{}{k}|)}{E_{\stackrel{}{p}+\stackrel{}{k}}E_p\omega }+\frac{\theta (p_Fp)\theta (|\stackrel{}{p}+\stackrel{}{k}|p_F)}{E_pE_{\stackrel{}{p}+\stackrel{}{k}}+\omega }\right].$$
(10)
Here $`E_p=p^2/2m^{}`$. The expression for $`\mathrm{\Pi }_\mathrm{\Delta }^0`$ is analogous.
The $`\pi NB`$ vertex $`\mathrm{\Gamma }_{\pi NB}`$ with $`B`$ labelling nucleon or $`\mathrm{\Delta }`$-isobar is
$$\mathrm{\Gamma }_{\pi NB}=\mathrm{\Gamma }_{\pi NB}^0d_B(k),$$
(11)
$$\mathrm{\Gamma }_{\pi NN}^0=i\frac{g_A}{\sqrt{2}f_\pi }\chi ^{}(\stackrel{}{\sigma }\stackrel{}{k})\chi ,\mathrm{\Gamma }_{\pi N\mathrm{\Delta }}^0=f_{\mathrm{\Delta }/N}i\frac{g_A}{\sqrt{2}f_\pi }\chi ^\alpha (\stackrel{}{S}_\alpha ^+\stackrel{}{k})\chi ,$$
(12)
where $`\chi `$ is nucleon 2-spinors, $`\stackrel{}{\sigma }`$ is nucleon spin, $`\stackrel{}{S}^+`$ turns the spin $`3/2`$ into $`1/2`$. In order to take into account the non-zero baryon size, the vertex $`\mathrm{\Gamma }_{\pi NB}^0`$ is multiplied by the form factor $`d_B(k)`$ taken in the form $`d_B=(1m_\pi ^2/\mathrm{\Lambda }_B^2)/(1+k^2/\mathrm{\Lambda }_B^2)`$.
The constants $`\gamma _{NN},\gamma _\mathrm{\Delta },\gamma _{\mathrm{\Delta }\mathrm{\Delta }}`$ are
$$\gamma _{NN}=C_0g_{NN}^{}\left(\frac{\sqrt{2}f_\pi }{g_A}\right)^2,\gamma _\mathrm{\Delta }=\frac{C_0g_{N\mathrm{\Delta }}^{}}{f_{\mathrm{\Delta }/N}}\left(\frac{\sqrt{2}f_\pi }{g_A}\right)^2,\gamma _{\mathrm{\Delta }\mathrm{\Delta }}=\frac{C_0g_{\mathrm{\Delta }\mathrm{\Delta }}^{}}{f_{\mathrm{\Delta }/N}^2}\left(\frac{\sqrt{2}f_\pi }{g_A}\right)^2,$$
where $`C_0`$ is the normalization factor ($`C_0=\pi ^2/(p_Fm^{})`$) and $`g^{}s`$ are the constants of the effective quasi-particle–quasi-hole interaction in nuclear matter . The calculation results are given for the following set of parameters :
$$f_{\mathrm{\Delta }/N}2,\mathrm{\Lambda }_N=0.667GeV,\mathrm{\Lambda }_\mathrm{\Delta }=1GeV,g_A=1,f_\pi =92MeV,$$
(13)
$$g_{NN}^{}=1.0,g_{N\mathrm{\Delta }}^{}=0.2,g_{\mathrm{\Delta }\mathrm{\Delta }}^{}=0.8.$$
When the parameters $`f_{\mathrm{\Delta }/N},g^{},\mathrm{\Lambda }_B`$ vary within the experimentally acceptable limits the dispersion equation solutions change quantitatively but not qualitatively.
Later on, we can perform integration in (10) in two ways.
1. To integrate separately the first and the second terms. In this way there appears the expression for $`\mathrm{\Pi }_N^0`$ as follows:
$`\mathrm{\Pi }_N^0(\omega ,k)`$ $`=`$ $`4\left({\displaystyle \frac{g_A}{\sqrt{2}f_\pi }}\right)^2k^2\left[\mathrm{\Phi }_N(\omega ,k)+\mathrm{\Phi }_N(\omega ,k)\right]d_N^2(k),`$ (14)
$`\mathrm{\Phi }_N(\omega ,k)`$ $`=`$ $`{\displaystyle \frac{m^{}}{k}}{\displaystyle \frac{1}{4\pi ^2}}({\displaystyle \frac{\omega m^{}+kp_F}{2}}\omega m^{}\mathrm{ln}\left({\displaystyle \frac{\omega m^{}}{\omega m^{}kp_F+k^2/2}}\right)+`$ (15)
$`+`$ $`{\displaystyle \frac{(kp_F)^2(\omega m^{}k^2/2)^2}{2k^2}}\mathrm{ln}\left({\displaystyle \frac{\omega m^{}kp_Fk^2/2}{\omega m^{}kp_F+k^2/2}}\right))`$
at $`0k2p_F`$. At $`k2p_F`$ $`\mathrm{\Phi }_N(\omega ,k)`$ is Migdal’s function:
$$\mathrm{\Phi }_N(\omega ,k)=\frac{1}{4\pi ^2}\frac{m^3}{k^3}\left[\frac{a^2b^2}{2}\mathrm{ln}\left(\frac{a+b}{ab}\right)ab\right]$$
(16)
where $`a=\omega (k^2/2m^{})`$, $`b=kp_F/m^{}`$. Consider now which cuts has the polarization operator $`\mathrm{\Pi }_N^0(\omega ,k)`$ (14)–(16) in the $`\omega `$-plane. We can see that at $`k2p_F`$ there are two cuts (define them as $`I`$ and $`II`$). They are related to the first and second logarithms in (15). The cuts are situated within the intervals:
$$I:0\omega \frac{kp_F}{m^{}}\frac{k^2}{2m^{}},II:\frac{kp_F}{m^{}}\frac{k^2}{2m^{}}\omega \frac{kp_F}{m^{}}+\frac{k^2}{2m^{}}.$$
(17)
Since $`\mathrm{\Pi }_N^0`$ is symmetrical under the replacement $`\omega \omega `$, the cuts of $`\mathrm{\Phi }_N(\omega ,k)`$ are placed symmetrically on the negative semiaxis. Thus $`\mathrm{\Pi }_N^0`$ has four cuts in the complex $`\omega `$-plane, they are shown in Fig.1.
2. The other way of integration in (10) gives a well-known expression of $`\mathrm{\Pi }_N^0`$ through Migdal’s functions. To follow it, let us do a substitution in (10) $`\theta (pp_F)1\theta (p_Fp)=1n(p)`$. Then, after the integration, $`\mathrm{\Phi }_N(\omega ,k)`$ takes a well-known form (16) for all values of $`k`$ . Undoubtedly, the expression (14) is the same for both integration ways. Now $`\mathrm{\Pi }_N^0`$ has not four cuts in $`\omega `$-plane but two (overlapping). It can be seen from the expressions (10), (15), that the cut $`I`$ is the sum of two overlapping cuts. The search for dispersion equation solutions becomes more convenient and obvious when we work with one cut but not with two overlapping ones. It is difficult to follow the solution in the case 2, therefore we use equations (14)–(16) for $`\mathrm{\Pi }_N^0`$.
Now let us turn to the polarization operator $`\mathrm{\Pi }_\mathrm{\Delta }^0`$, which is the isobar–nucleon-hole loop. Since the isobar Fermi surface is absent at the nuclear densities and isobar momentum is unrestricted, there is no problem discussed above and $`\mathrm{\Pi }_\mathrm{\Delta }^0(\omega ,k)`$ reads:
$$\mathrm{\Pi }_\mathrm{\Delta }^0=\frac{16}{9}\left(\frac{g_A}{\sqrt{2}f_\pi }\right)^2f_{\mathrm{\Delta }/N}^2k^2\left[\mathrm{\Phi }_\mathrm{\Delta }(\omega ,k)+\mathrm{\Phi }_\mathrm{\Delta }(\omega ,k)\right]d_\mathrm{\Delta }^2(k).$$
(18)
The function $`\mathrm{\Phi }_\mathrm{\Delta }(\omega ,k)`$ is expressed through Migdal’s functions (16) with $`a=\omega (k^2/2m^{})\mathrm{\Delta }m`$, $`b=kp_F/m^{}`$. The mass difference, $`\mathrm{\Delta }m=m_\mathrm{\Delta }m`$, is the following: Re$`(\mathrm{\Delta }m)=292`$ MeV and Im$`(\mathrm{\Delta }m)=\mathrm{\Gamma }_\mathrm{\Delta }/2`$. The cuts of $`\mathrm{\Phi }_\mathrm{\Delta }^0(\omega ,k)`$ are shown in Fig.1. At $`\omega 0`$ the cut is in the interval
$$\frac{k^2}{2m^{}}+\mathrm{\Delta }m\frac{kp_F}{m^{}}\omega \frac{k^2}{2m^{}}+\mathrm{\Delta }m+\frac{kp_F}{m^{}}.$$
(19)
The cut is shifted into the complex plane in the value $`i\mathrm{\Gamma }_\mathrm{\Delta }/2`$.
In this paper Landau equation is used for the nucleon effective mass
$$m^{}=\frac{m}{1+(2mp_F/\pi ^2)f_1}.$$
(20)
Unknown parameter $`f_1`$ is fixed by the condition $`m^{}(p_F=p_{F0})=0.8m`$.
## 3 Solutions of the dispersion equation
In this section the solutions of the dispersion equation (1) are presented. The solution branch $`\omega _c`$ emerges on the physical sheet of the complex $`\omega `$-plane at $`p_F=283`$ MeV for the parameter values (13). The figures for the zero sound branch $`\omega _s(k)`$, pion branch $`\omega _\pi (k)`$ and isobar branch $`\omega _\mathrm{\Delta }(k)`$ are presented for $`\mathrm{\Gamma }_\mathrm{\Delta }=115`$ MeV at $`p_F=`$268 and 290 MeV (i.e. at the equilibrium density and at density slightly larger than critical one). Some special cases, with the other values of $`p_F`$ and $`\mathrm{\Gamma }_\mathrm{\Delta }`$, are considered as well.
Branch $`\omega _s(k)`$. (Fig.2a) The branch $`\omega _s(k)`$ is shown for $`p_F=268`$ and $`290`$ MeV (curves 1 and 2 correspondingly). The branch begins at $`\omega _s(k=0)=0`$, then while $`k`$ increases, moves practically along the real axis. At $`k_f=0.430m_\pi `$ for $`p_F=290`$ MeV (at $`k_f=0.436m_\pi `$ for $`p_F=268`$ MeV) goes under the cut $`II`$ (17), this corresponds to the decay of $`\omega _s`$ into real nucleon and the hole.
Branch $`\omega _\mathrm{\Delta }(k)`$. (Fig.2b) The isobar branch $`\omega _\mathrm{\Delta }(k)`$ begins at $`\omega =\mathrm{\Delta }m`$ at $`k=0`$ and ends on the isobar cut (19) at $`k_f=5.1m_\pi `$ for $`p_F=268`$ MeV ($`k_f=4.8m_\pi `$ for $`p_F=290`$ MeV).
Branch $`\omega _\pi (k)`$. (Fig.2c) The pion branch starts at $`k=0`$ in $`\omega _\pi =m_\pi ^{}`$ (see (5),(6)). The beginning of $`\omega _\pi (k=0)`$ is shifted to the smaller then $`m_\pi `$ values when GMOR is used to determine $`m_\pi ^{}`$. The pion branch ends on physical sheet under the isobar cut, this corresponds to the decay of pion into isobar and nucleon hole. It takes place at $`k_f=3.5m_\pi `$ for $`p_F=268`$ MeV $`(k_f=3.8m_\pi `$, $`p_F=290`$ MeV).
Branch $`\omega _c(k)`$. (Fig.2d, 3a,b) While the density increases one more branch of solutions, $`\omega _c(k)`$, appears on the physical sheet. It emerges at $`p_F283`$ MeV when parameter values (13) are used. In Fig.2d the pion branch $`\omega _\pi (k)`$ and condensate branch $`\omega _c(k)`$ are presented at $`p_F=290`$ MeV. The dashed piece of $`\omega _c(k)`$ belongs to the upper unphysical sheet of the logarithmic cut $`I`$ (17) (Fig.1). The branch $`\omega _c(k)`$ starts at $`k=0`$ at the same point as $`\omega _\pi (k)`$ and moves onto the unphysical sheet; at $`k=k_1=1.3m_\pi `$ the branch goes down to the physical sheet and at $`k=k_2=2.3m_\pi `$ moves back to the same unphysical sheet. In the momentum interval $`(k_1,k_2)`$ one can follow over all the branches shown in Fig.2 and check that $`\omega _c`$ does not belong to any branches studied before ($`\omega _s,\omega _\pi ,\omega _\mathrm{\Delta }`$).
The branch $`\omega _c`$ depends on the isobar width $`\mathrm{\Gamma }_\mathrm{\Delta }`$. Decreasing the isobar width we see that the isobar cuts move to the real axis and $`\omega _s`$, $`\omega _\pi `$ and $`\omega _\mathrm{\Delta }`$ have smaller imaginary parts. When $`\mathrm{\Gamma }_\mathrm{\Delta }=0`$ the branches $`\omega _s`$, $`\omega _\pi `$ and $`\omega _\mathrm{\Delta }`$ are real. On the contrary, $`\omega _c(k)`$ moves to the imaginary axis with decreasing $`\mathrm{\Gamma }_\mathrm{\Delta }`$ (Fig.3a). At $`\mathrm{\Gamma }_\mathrm{\Delta }=0`$ we have pure imaginary solutions on the physical sheet: $`\omega _c^20`$.
In Fig.3b the branch $`\omega _c(k)`$ is shown at the different densities: $`p_F=280,290,300,360`$ MeV. When $`p_F=280`$ MeV the whole branch (curve 1) is located on unphysical sheet. At $`p_F=283`$ MeV the branch touches the real axis (not shown). For $`p_F>283`$ Mev the part of $`\omega _c(k)`$ in the interval $`(k_1,k_2)`$ is placed on the physical sheet. At the critical density $`\rho =rho_c`$ and $`\mathrm{\Gamma }_\mathrm{\Delta }=0`$ $`\omega _c(k)`$ touches the real axis at the point $`\omega _c(k)=0.`$
It was shown in paper that the appearance of $`\omega _c`$ on the physical sheet results not only in pion condensation but in restoration of chiral symmetry in the nuclear matter at critical density as well.
## 4 Conclusion
In the paper the solutions of pion dispersion equation are considered in details in the complex $`\omega `$-plane. It is shown that, besides the well-known solutions with quantum numbers $`0^{}`$ (zero spin-isospin sound, pion and isobar waves), there exists the fourth branch $`\omega _c(k)`$. It is the branch which obeys the condition $`\omega _c^20`$, therefore it is responsible for instability of the ground state. Such instability can indicate the beginning of ’pion condensation’. We demonstrate that at the density less than critical one, $`\rho <\rho _c`$ the branch $`\omega _c(k)`$ is situated on unphysical sheet and at $`\rho \rho _c`$ it comes on physical one.
### Acknowledgments
I am grateful to M.G. Ryskin for the important and fruitful discussions during the work. I thanks E.G. Drukarev and E.E. Saperstein for useful discussions. This work was supported by RFFI grant 96-15-96764.
## 5 Figure captions
Fig.1. The cuts on the physical sheet of the complex $`\omega `$-plane of polarization operators $`\mathrm{\Pi }_N^0`$, $`\mathrm{\Pi }_\mathrm{\Delta }^0`$, corresponding to equations (14), (17), (18), (19). The cuts are presented at $`p_F=290`$ MeV, $`k=m_\pi `$.
Fig.2. Branches of solutions of (1) in the complex $`\omega `$-plane. Curves 1 and 2 stand for $`p_F=268,290`$ MeV, correspondingly. The dashed pieces of curves are situated on unphysical sheets. a) Zero spin-isospin sound branch $`\omega _s(k)`$. The curves are presented up to 1.6$`m_\pi `$. b) The isobar branch $`\omega _\mathrm{\Delta }`$. The horizontal dashed line is a logarithmic cut (19) for $`p_F=290`$ MeV at momentum $`k`$ when $`\omega _\mathrm{\Delta }(k)`$ is on the cut. c) The pion branch $`\omega _\pi `$. The horizontal dashed line is a logarithmic cut (19) for $`p_F=290`$ MeV at momentum $`k`$ when $`\omega _\pi `$ is on the cut. d) The total picture at $`p_F=290`$ MeV for pion branch $`\omega _\pi (k)`$ and condensate branch $`\omega _c`$(k).
Fig.3. Condensate branch $`\omega _c`$ in the complex $`\omega `$-plane. Here the dashed pieces of branches are on the physical sheet, but solid lines belong to unphysical sheet. a) The branch $`\omega _c`$ is presented at $`p_F=290`$ MeV for different values of isobar width: $`\mathrm{\Gamma }_\mathrm{\Delta }=0,10,50,115`$ MeV (curves 1,2,3,4, correspondingly); $`\omega _c(k=0)=0.744m_\pi `$. b) The branch $`\omega _c`$ is presented at $`\mathrm{\Gamma }_\mathrm{\Delta }=115`$ MeV for different values of Fermi momenta $`p_F=280,290,300,360`$ MeV (curves 1,2,3,4, correspondingly). For $`p_F`$=300 and 360 MeV only the part of the branch, which is placed on the physical sheet, is shown (the curves 3 and 4). The whole branch for $`p_F=280`$ MeV (curve 1) is on unphysical sheet. |
warning/0001/math0001162.html | ar5iv | text | # Combed 3-Manifolds with Concave Boundary, Framed Links, and Pseudo-Legendrian Links
## Introduction
This paper describes combinatorial realizations, based on the machinery of branched standard spines (see Section 1) of the following three topological categories (in which manifolds and diffeomorphisms are oriented by default):
1. Combed $`3`$-manifolds with concave boundary, that is pairs $`(M,v)`$, where $`M`$ is a compact 3-manifold (possibly with boundary), and $`v`$ is a nowhere-zero vector field on $`M`$ with simple tangency circles of concave type on $`M`$, up to diffeomorphism of $`M`$ and homotopy of $`v`$ through fields of the same sort;
2. Framed links in $`3`$-manifolds, that is pairs $`(M,L)`$, where $`M`$ is as above and $`L`$ is a framed link in $`M`$, up to diffeomorphism of $`M`$ and framed isotopy of $`L`$;
3. Pseudo-Legendrian links in combed $`3`$-manifolds, that is triples $`(M,v,L)`$, where $`(M,v)`$ is as above and $`L`$ is transversal to $`v`$, up to diffeomorphism of $`M`$ and ‘pseudo-Legendrian isotopy’ of $`(v,L)`$, i.e. simultaneous homotopy of $`v`$ and isotopy of $`L`$ through pairs $`(v,L)`$ of the same type.
We will denote these categories respectively by $`\mathrm{Comb}`$, $`\mathrm{Fram}`$ and $`\mathrm{PLeg}`$. (Regarding names, recall that a non-zero vector field up to homotopy is often called a combing, and that if $`\xi `$ is an oriented contact structure and $`L`$ is Legendrian in $`\xi `$, then $`(M,\xi ^{},L)`$ defines an element of $`\mathrm{PLeg}`$.) Our realizations are given according to the by now popular scheme in 3-dimensional topology, namely:
1. A class of combinatorial objects, each of which can be specified by a finite set of data, and a surjective reconstruction map which assigns to a combinatorial object a topological one;
2. A finite set of local combinatorial moves on objects, finite combinations of which give the equivalence relation induced by the reconstruction map.
In the definitions of the topological categories given above we have been forced to include the action of diffeomorphisms, because we use spines, which determine manifolds only up to diffeomorphism. However if a certain manifold $`M`$ is given we can restrict to spines embedded in $`M`$ (rather than abstract ones), and get formally identical combinatorial realizations of the refined categories where only diffeomorphisms of $`M`$ isotopic to the identity are considered. We will mention how to do this in Section 2 for combings, but a similar refinement could easily be stated for framed links and for pseudo-Legendrian links.
Rather than providing precise statements of our realizations, in this introduction we give some general background and motivations, starting with $`\mathrm{Comb}`$. A combinatorial realization of the subcategory $`\mathrm{Comb}^{\mathrm{cl}}`$ of $`\mathrm{Comb}`$ given by pairs $`(M,v)`$ with closed $`M`$ was given in . In Section 2 we extend the arguments of to the bounded case, and we actually refine the results proved there, by showing that some of the moves previously considered may actually be neglected. The realization of $`\mathrm{Comb}^{\mathrm{cl}}`$ in was the basis for the treatment of other refinements of the category of 3-manifolds, involving spin structures and framings. These realizations proved fruitful in connection with spin-refined Turaev-Viro invariants (see Section 8.3 in ) and G. Kuperberg’s invariants for combed and framed manifolds, of which a very constructive description is given in . Our main motivation here comes from , where we have developed a theory of Euler structures with simple boundary and their Reidemeister-Turaev torsion (see ). The surjectivity of the reconstruction map of the realization of $`\mathrm{Comb}`$ was used in to construct an explicit canonical $`H_1`$-equivariant bijection from the space of smooth Euler structures to the space of combinatorial Euler structures, and to exhibit a canonical Euler chain for the structure carried by a branched spine.
The subcategory of $`\mathrm{Fram}`$ consisting of framed links in closed manifolds was combinatorially realized by Turaev in terms of link diagrams on a given standard spine, and moves on these diagrams (including the classical framed Reidemeister moves). In Section 3 we modify the situation considered by Turaev by taking a branched standard spine of the manifold, and restricting to link diagrams which are $`\mathrm{C}^1`$ with respect to the branching. On one hand, this allows to simplify the encoding of the framing, because the field carried by the spine is automatically transverse to the link, while Turaev needs to add half-twists. On the other hand, some technical complications emerge, because only $`\mathrm{C}^1`$ moves can be used. Nevertheless, a result formally analogous to Turaev’s turns out to be true, yielding the presentation of $`\mathrm{Fram}`$ discussed in Section 3. In Section 4 we exploit the fact that if a branched spine defines a global field on a manifold, according to the scheme given for $`\mathrm{Comb}`$, then the link defined by a $`\mathrm{C}^1`$ diagram on the spine is automatically pseudo-Legendrian with respect to the field. This leads us to the presentation of $`\mathrm{PLeg}`$.
Comparing the presentations of $`\mathrm{Fram}`$ and $`\mathrm{PLeg}`$ one notices a rather remarkable feature: the former is obtained from the latter just by adding the ‘curl’ (first Reidemeister) move. This fact has two interesting interpretations:
* it is a perfect combinatorial analogue of the imitation of a framed isotopy by a Legendrian isotopy in a contact manifold;
* it allows a partial extension of the notion of winding number of a link diagram.
The imitation mentioned in (a) plays a central role in the comparison, due to Fuchs-Tabachnikov and Tchernov of framed and Legendrian finite-order invariants, and we believe that our combinatorial realizations could be of some help in the understanding of these invariants. In particular, we conjecture that the right environment in which finite-order invariants should be considered in precisely our category $`\mathrm{PLeg}`$ (we will provide an exact statement and some evidence in Section 5).
Concerning (b), recall first that if a knot $`K:S^1\text{}^3`$ is transverse to the constant vertical field $`/z`$, then its equivalence class up to isotopy transverse to $`/z`$ is determined by the framed isotopy class and by the ‘winding’ number (the degree of $`\pi K^{}`$, where $`\pi `$ is the obvious projection on the horizontal unit circle). Using our presentations of $`\mathrm{Fram}`$ and $`\mathrm{Comb}`$ we can show that a partial analogue of this fact is true in any combed manifold with concave boundary, provided one allows a homotopy of the field simultaneous with the isotopy of the knot. In the general setting, however, the winding number only exists as a relative object, and we can prove that it leads to a well-defined invariant only under the assumption that the knot is ‘good.’ The notion of ‘goodness’ for knots emerged in our study of torsion as a relative invariant of pairs of pseudo-Legendrian knots which are framed-isotopic . Many knots are good: for instance, all knots are good if the ambient manifold is a homology sphere, and most knots with hyperbolic complement are good. In Section 5 we will give some applications of the notion of relative winding number, in connection with torsion and finite-order invariants. In particular we will show the following:
###### Proposition 0.1
Let $`M`$ be a homology sphere, let $`v`$ be a field on $`M`$ and consider two pseudo-Legendrian knots in $`(M,v)`$ which are isotopic as framed knots. Then the following conditions are pairwise equivalent:
1. the knots are pseudo-Legendrian isotopic;
2. the relative winding number vanishes;
3. the knots have the same Maslov index;
4. the knots cannot be distinguished by the relative torsion invariants of ;
5. the knots are homotopic as pseudo-Legendrian immersions.
Moreover we will prove that torsion invariants cannot distinguish the pairs of framed-isotopic Legendrian knots given in , which Tchernov shows to be distinguished by finite-order invariants.
We conclude by giving another perspective of the realization of $`\mathrm{PLeg}`$. Recall that $`\mathrm{PLeg}`$ comes as a refinement of Turaev’s presentation of $`\mathrm{Fram}`$, which was the starting point of his beautiful theory of 4-dimensional shadows. We believe that the extra structure given by the branching of the spine, which underlies the presentation of $`\mathrm{PLeg}`$, should have 4-dimensional counterparts. Our intuition is that “4-dimensional branched shadows”, which are not quite defined yet, should correspond to $`\mathrm{Spin}^\mathrm{c}`$ structures on 4-manifolds and allow to treat their invariants. This intuition is supported by the fact that in dimension three branched spines indeed are a good framework to treat torsion of Euler structures (and hence, in particular, Seiberg-Witten invariants of closed 3-manifolds with $`\mathrm{Spin}^\mathrm{c}`$ structures, see ).
Acknowledgment. Section 5 owes a lot to very useful discussions we had with Vladimir Tchernov.
## 1 Branched spines and combings
This section contains many definitions used below and reviews the theory developed in .
#### Manifolds and fields
All the manifolds we will consider are 3-dimensional, oriented, and compact, with or without boundary. Using the Hauptvermutung, we will somewhat intermingle the differentiable and piecewise linear viewpoints. Maps will always respect orientations. All vector fields mentioned in this paper will be non-singular, and they will be termed just fields for the sake of brevity. A field $`v`$ on a manifold $`M`$ is called traversing if its orbits eventually intersect $`M`$ transversely in both directions (in other words, orbits are compact intervals or points). A point where $`v`$ is tangent to $`M`$ is called simple if it appears in a cross-section as in Fig. 1.
The field is called concave if it is tangent to $`M`$ only in a concave fashion, as shown on the left in the figure. Given a concave field $`v`$ on $`M`$, the boundary of $`M`$ naturally splits into the region on which $`v`$ points outside $`M`$ (which we denote by $`B`$ and call the black region), and the region on which $`v`$ points inside (denoted by $`W`$ and called white). Note that $`B=W`$ is a union of circles. The pair $`(B,W)`$, which is actually determined by any two of its elements, is called a boundary pattern on $`M`$. (This definition is a simplified version of that given in , because here we do not allow convex tangency.) Starting from the Poincaré-Hopf theorem one can show that a given boundary pattern $`𝒫=(B,W)`$ on $`M`$, i.e. a splitting of $`M`$ into two surfaces with common boundary, arises from a concave field if and only if $`\chi (W)=\chi (M)`$. See .
#### Standard spines
A simple polyhedron $`P`$ is a finite, connected, purely 2-dimensional polyhedron with singularity of stable nature (triple lines and points where six non-singular components meet). Such a $`P`$ is called standard if all the components of the natural stratification given by singularity are open cells. Depending on dimension, we will call the components vertices, edges and regions.
A standard spine of a $`3`$-manifold $`M`$ with $`M\mathrm{}`$ is a standard polyhedron $`P`$ embedded in $`\mathrm{Int}(M)`$ so that $`M`$ collapses onto $`P`$. Standard spines of oriented $`3`$-manifolds are characterized among standard polyhedra by the property of carrying an orientation, defined (see Definition 2.1.1 in ) as a “screw-orientation” along the edges (as in the left-hand-side of Fig. 2), with the obvious compatibility at vertices (as in the centre of Fig. 2).
It is the starting point of the theory of standard spines that every oriented $`3`$-manifold $`M`$ with $`M\mathrm{}`$ has an oriented standard spine, and can be reconstructed (uniquely up to equivalence) from any of its oriented standard spines. See for the non-oriented version of this result and or Proposition 2.1.2 in for the (slight) oriented refinement. We will denote by $`M(P)`$ the manifold defined by $`P`$. Note that $`M(P)\mathrm{}`$. To recover closed manifolds one considers spines $`P`$ such that $`M(P)S^2`$, and defines $`\widehat{M}(P)`$ as $`M(P)_fD^3`$ with $`f:M(P)S^2`$ a diffeomorphism. Note that this definition makes sense also when $`M(P)`$ has more than one component, but at least one is a sphere.
#### Moves for standard spines
The fundamental move for standard spines, which (in both directions) preserves the topological type of the associated manifold, is the Matveev-Piergallini MP move, see and Fig. 3.
Counting the vertices involved one is naturally led to call the positive MP a “2-to-3” move. The MP-move and its inverse are actually not sufficient to relate spines of the same manifold, because they obviously cannot apply to spines with one vertex. However, as soon as one decides to dismiss these “MP-rigid” spines (not the corresponding manifolds, which have plenty of other spines), the MP-move does become sufficient . To deal with spines with one vertex the “0-to-2” move of Fig. 4 (and its inverse)
must be added.
#### Branched spines
A branching on a standard polyhedron $`P`$ is an orientation for each region of $`P`$, such that no edge is induced the same orientation three times. See the right-hand side of Fig. 2 and Definition 3.1.1 in for the geometric meaning of this notion. An oriented standard spine $`P`$ endowed with a branching is shortly named branched spine. We will never use specific notations for the extra structures: they will be considered to be part of $`P`$. The following result, proved as Theorem 4.1.9 in , is the starting point of our constructions.
###### Proposition 1.1
To every branched spine $`P`$ there corresponds a manifold $`M(P)`$ with non-empty boundary and a concave traversing field $`v(P)`$ on $`M(P)`$. The pair $`(M(P),v(P))`$ is well-defined up to equivalence, and an embedding $`i:P\mathrm{Int}(M(P))`$ is defined with the property that $`v(P)`$ is positively transversal to $`i(P)`$.
The topological construction which underlies this proposition is actually quite simple, and it is illustrated in Fig. 5. Concerning the last
assertion of the proposition, note that the branching allows to define an oriented tangent plane at each point of $`P`$.
#### Non-traversing fields and closed manifolds
As noted above, standard spines do not directly represent closed manifolds, but one can use spines of manifolds bounded by $`S^2`$ and cap off this sphere to get a closed manifold, or, viewing things the other way around, one can remove an open ball from a given closed manifold to get a bounded one. When one is interested in a manifold equipped with a field, one can try to use branched spines, but of course one sees that they are inadequate to give a direct description both when the manifold is closed and when the field is non-traversing. This limitation is circumvented again by removing a ball, with a proviso on the field on that ball.
Let $`P`$ be a branched standard spine, and assume that in $`M(P)`$ there is only one component which is diffeomorphic to $`S^2`$ and is split by the tangency line of $`v(P)`$ to $`M(P)`$ into two discs. Such a component will be denoted by $`S_{\mathrm{triv}}^2`$. Now, notice that $`S_{\mathrm{triv}}^2`$ is also the boundary of the closed $`3`$-ball with constant vertical field, denoted by $`B_{\mathrm{triv}}^3`$. This shows that we can cap off $`S_{\mathrm{triv}}^2`$ by attaching a copy of $`B_{\mathrm{triv}}^3`$, getting a compact manifold $`\widehat{M}(P)`$ and a concave field $`\widehat{v}(P)`$ on $`\widehat{M}(P)`$. If we denote by $`\widehat{𝒫}(P)`$ the boundary pattern of $`\widehat{v}(P)`$ on $`\widehat{M}(P)`$, we easily see that the pair $`(\widehat{M}(P),\widehat{v}(P))`$ is only well-defined up to diffeomorphism of $`\widehat{M}(P)`$ and homotopy of $`\widehat{v}(P)`$ through fields compatible with $`\widehat{𝒫}(P)`$.
#### Standard sliding moves
Let $`PP^{}`$ be a positive MP-move (so, $`P^{}`$ has one vertex and one region more than $`P`$). If $`P`$ has a branching, all the regions of $`P^{}`$, except for the new one, already have an orientation, and it is a fact that the new region can always be given an orientation (sometimes not a unique one) so to get a branching on $`P^{}`$. Each of the moves on branched spines arising like this will be called a branched MP-move, and it will be called a sliding MP-move if moreover it does not modify the boundary pattern of the associated concave field. One can actually see that each sliding-MP-move can be realized within a certain pair $`(M,v)`$ as a continuous deformation through branched spines of $`M`$ transverse to $`v`$, with only one singularity at which the spine is non-standard but transversality is preserved. This deformation is shown in Fig. 6,
and it justifies the term ‘sliding’ quite clearly. Since in Fig. 6 we are showing portions of spines embedded in $`\text{}^3`$, to give a completely intrinsic description of the moves we should specify in each portion whether the screw-orientation of the spine is equal or opposite to that induced by $`\text{}^3`$, and whether the upward vertical field is positively or negatively transversal to the spine. As a result, the complete list of sliding-MP-moves contains 16 different ones, but the essential physical modifications are only those shown in Fig. 6. From this figure one also sees quite clearly that if $`\widehat{v}(P)`$ and $`\widehat{v}(P^{})`$ can be defined then they coincide (up to homotopy through fields compatible with $`\widehat{𝒫}(P)=\widehat{𝒫}(P^{})`$). Another move which obviously has the same property, and will be needed below, is the branched version of the 0-to-2 move, shown Fig. 7, and called the snake move
in the sequel. As above, if one takes orientations into account, there is another essentially different snake move, obtained by mirroring Fig. 7. Since also the snake move involves a sliding, we will call standard sliding move any sliding-MP or snake move.
## 2 A calculus for combed manifolds <br>with concave boundary
In this section we will extend and refine the main results of Chapter 5 of . The extension consists in passing from the closed to the bounded case, and the refinement comes from the shortening of the list of moves to be considered. More precisely, we will show that compact manifolds with concave combings are combinatorially described by (suitable) branched spines up to certain moves, namely the standard sliding (snake and sliding-MP) moves shown above. Moreover, we will show that spines which are rigid with respect to sliding-MP-moves can be dismissed with no harm, and that the sliding-MP-moves suffice to generate the equivalence on the remaining spines. This implies that our result is a perfect combed analogue of the Matveev-Piergallini theorem (but our proof is self-contained).
#### Definitions and statements
We will denote by $`\mathrm{Comb}`$ the set of all pairs $`(M,v)`$, where $`M`$ is a compact oriented manifold and $`v`$ is a concave field on $`M`$, viewed up to diffeomorphism of $`M`$ and homotopy of $`v`$ through concave fields. A class $`[M,v]\mathrm{Comb}`$ is called a combing on the diffeomorphism class of the manifold $`M`$. Note that the boundary pattern on $`M`$ evolves isotopically during a homotopy of $`v`$, so a pair $`(M,𝒫)`$, viewed up to diffeomorphism of $`M`$, can be associated to each $`[M,v]\mathrm{Comb}`$. In particular, $`\mathrm{Comb}`$ naturally splits as the disjoint union of subsets $`\mathrm{Comb}([M,𝒫])`$, consisting of combings on $`M`$ compatible with $`𝒫`$.
For a technical reason we actually rule out from $`\mathrm{Comb}`$ the set of those classes $`[M,v]`$ such that the corresponding boundary pattern contains components of the type $`S_{\mathrm{triv}}^2`$. This is actually not a serious restriction, because each $`S_{\mathrm{triv}}^2`$ component can be capped off by a $`B_{\mathrm{triv}}^3`$, and the result is well-defined up to homotopy. Note that we do accept pairs $`(M,v)`$ with closed $`M`$, and pairs in which $`v`$ has no tangency at all to $`M`$.
Let us denote now by $``$ the set of all branched spines $`P`$ (up to PL isomorphism) such that $`𝒫(P)`$ contains only one $`S_{\mathrm{triv}}^2`$. Such a $`P`$ being given, $`\widehat{M}(P)`$ and $`\widehat{v}(P)`$ can be considered, and the pair $`(\widehat{M}(P),\widehat{v}(P))`$ gives rise to a well-defined element of $`\mathrm{Comb}`$, which we denote by $`C(P)`$. The following will be shown below:
###### Theorem 2.1
The map $`C:\mathrm{Comb}`$ is surjective, and the equivalence relation defined by $`C`$ on $``$ is generated by sliding-MP-moves and snake moves.
###### Remark 2.2
The following interpretation of the surjectivity of $`C`$ is perhaps useful. Note first that the dynamics of a field, even a concave one, can be very complicated, whereas the dynamics of a traversing field (in particular, $`B_{\mathrm{triv}}^3`$) is simple. Surjectivity of $`C`$ means that for any (complicated) concave field there exists a sphere $`S^2`$ which splits the field into two (simple) pieces: a standard $`B_{\mathrm{triv}}^3`$ and a concave traversing field. Actually, a 1-parameter version of this statement also holds (see Remark 2.6): we will need it to show that the $`C`$-equivalence is the same as the sliding equivalence.
As announced, we state now the sliding analogue of the fact that the MP moves suffice. Let us denote by $``$ the subset of $``$ consisting of the branched spines which are “rigid” from the point of view of the sliding-MP-moves, i.e. the spines to which no such move applies. An explicit description of $``$ is given in the proof of the next result. In the statement we only emphasize the most important consequences of this description.
###### Proposition 2.3
1. For every surface $`\mathrm{\Sigma }`$ and pattern $`𝒫`$ on $`\mathrm{\Sigma }`$ there are at most two spines $`P`$ such that $`(M(P))(\mathrm{\Sigma },𝒫)`$.
2. If two elements of $``$ are related through sliding-MPmoves and snake moves, they are also related through sliding-MP-moves only.
3. Every $`P`$ is related by a snake move to an element of $``$.
This proposition shows that in the statement of Theorem 2.1 one may remove $``$ from $``$ and forget the snake move.
#### Embedding-refined calculus
We spell out in this paragraph the embedding-refined version of our calculus, which allows to neglect the action of automorphisms. Let a certain manifold $`M`$ be given, and consider the set $`\mathrm{Comb}(M)`$ of concave vector fields on $`M`$, up to homotopy. Let $`(M)`$ consist of the elements of $``$ which are smoothly embedded in $`M`$ as spines of $`M`$ minus a ball $`B_{\mathrm{triv}}^3`$. Each element $`P`$ of $`(M)`$ is viewed up to isotopy in $`M`$, and gives rise to a well-defined element $`C_M(P)`$ of $`\mathrm{Comb}(M)`$. Moreover sliding-MP-moves and snake moves are well-defined in $`(M)`$, because they can be realized as embedded moves. The embedded analogue of Theorem 2.1 states that $`C_M:(M)\mathrm{Comb}(M)`$ is surjective, and the relation it defines is generated by the embedded moves. The proof of this result is a refinement of the proof of the general statement, along the lines explained in (4.1.12, 4.1.13, 4.3.5, and 5.2.1.)
#### Normal sections of a concave field
The proof of Theorem 2.1 is an extension of the argument given in Chapter 5 of , and it is based on the following technical notion, which extends ideas originally due to Ishii . Let $`v`$ be a concave field on $`M`$. Let $`B_1,\mathrm{},B_k`$ be the black components of the splitting of $`M`$, i.e. the regions on which $`v`$ points outwards. A normal section for $`(M,v)`$ is a compact surface $`\mathrm{\Sigma }`$ with boundary, embedded in the interior of $`M`$, with the following properties:
1. $`v`$ is transverse to $`\mathrm{\Sigma }`$;
2. $`\mathrm{\Sigma }`$ has exactly $`k+1`$ components $`\mathrm{\Sigma }_0,\mathrm{},\mathrm{\Sigma }_k`$, with $`\mathrm{\Sigma }_0D^2`$;
3. For $`i>0`$, the projection of $`B_i`$ on $`\mathrm{\Sigma }`$ along the orbits of $`v`$ is well-defined and yields a diffeomorphism between $`B_i`$ and a surface $`B_i^{}`$ contained in the interior of $`\mathrm{\Sigma }_i`$, with $`\mathrm{\Sigma }_iB_i^{}`$ being a collar on $`\mathrm{\Sigma }_i`$ (so $`B_i^{}=\mathrm{\Sigma }_i`$ if $`\mathrm{\Sigma }_i=\mathrm{}`$);
4. Each positive half-orbit of $`v`$ meets either the interior of some $`B_i`$ (where it stops), or the interior of some $`\mathrm{\Sigma }_i`$;
5. $`\mathrm{\Sigma }`$ meets itself generically along $`v`$ (i.e. each orbit of $`v`$ meets $`\mathrm{\Sigma }`$ at most two consecutive times on $`\mathrm{\Sigma }`$, and, if so, transversely);
6. Let $`P_\mathrm{\Sigma }`$ be the union of $`\mathrm{\Sigma }`$ with all the orbit segments starting on $`\mathrm{\Sigma }`$ and ending on $`\mathrm{\Sigma }`$. Then $`\mathrm{\Sigma }`$, which is a simple polyhedron by the previous point, is actually standard.
The next two lemmas show that normal sections of $`(M,v)`$ correspond bijectively to spines $`P`$ such that $`C(P)=[M,v]`$. The proof of surjectivity of $`C`$ and the discussion of its non-injectivity will be based on these lemmas.
###### Lemma 2.4
If $`(M,v)`$, $`\mathrm{\Sigma }`$ and $`P_\mathrm{\Sigma }`$ are as above, then $`P_\mathrm{\Sigma }`$ can be given a structure of branched spine such that $`C([P_\mathrm{\Sigma }])=[M,v]`$.
Proof of2.4. We orient $`\mathrm{\Sigma }`$ so that $`v|^+\mathrm{\Sigma }`$ (by default $`M`$ is oriented). Every region of $`P_\mathrm{\Sigma }`$ contains some open portion of $`\mathrm{\Sigma }`$, so it can be oriented accordingly; with the obvious screw-orientation, this turns $`P_\mathrm{\Sigma }`$ into a branched spine of its regular neighbourhood in $`M`$.
We show that $`C([P_\mathrm{\Sigma }])=[M,v]`$ by embedding the abstract manifold $`M(P_\mathrm{\Sigma })`$ in $`M`$, in such a way that the field carried by $`P_\mathrm{\Sigma }`$ on $`M(P_\mathrm{\Sigma })M`$ is just the restriction of $`v`$. By construction, $`MM(P_\mathrm{\Sigma })`$ will consist of a copy of $`B_{\mathrm{triv}}^3`$, together with a collar on $`M`$ which can be parameterized as $`(M)\times [0,1]`$ in such a way that $`v`$ is constant in the $`[0,1]`$-direction. This easily implies that $`C([P_\mathrm{\Sigma }])=[M,v]`$ indeed.
We illustrate the embedding of $`M(P_\mathrm{\Sigma })`$ in $`M`$ pictorially in one dimension less. Figure 8 shows how $`\mathrm{\Sigma }_0`$ gives rise to a $`B_{\mathrm{triv}}^3`$.
In the figure we describe $`v`$ by dotted lines, $`\mathrm{\Sigma }`$ by thick lines, portions of $`P_\mathrm{\Sigma }\mathrm{\Sigma }`$ by thin lines, and $`(M(P_\mathrm{\Sigma }))`$ by a thick dashed line. Note also that the portions of $`P_\mathrm{\Sigma }\mathrm{\Sigma }`$ have been slightly modified so to become positively transversal to $`v`$, which allows us to represent the branching as usual, i.e. as a $`\mathrm{C}^1`$ structure on $`P_\mathrm{\Sigma }`$.
Figure 9 shows the collar based on a component of $`M`$.
We use the same conventions as in the previous figure, and in addition we represent the black and white components of $`M`$ by thick and thin lines respectively. This description concludes the proof. 2.4
###### Lemma 2.5
Let $`[P]`$ and $`C([P])=[M,v]\mathrm{Comb}`$, with $`P`$ embedded in $`(M,v)`$ according to the geometric description of $`C`$. Let $`\mathrm{\Sigma }`$ be obtained from $`P`$ as suggested (in one dimension less) in Fig. 10.
Then $`\mathrm{\Sigma }`$ is a normal section of $`(M,v)`$, and $`P_\mathrm{\Sigma }`$ is isomorphic to $`\mathrm{\Sigma }`$.
Proof of2.5. The construction suggested by Fig. 10 is obviously the inverse of the construction in the proof of Lemma 2.4. 2.5
#### The concave combing calculus
Using normal sections we can now show the main result of this section.
Proof of2.1. We start with the proof of surjectivity. So, let us consider a combed manifold $`(M,v)`$, subject to the usual restrictions. By Lemma 2.4 it is natural to try and construct a normal section for $`(M,v)`$. Let $`B_1,\mathrm{},B_k`$ be the black regions in $`M`$. Slightly translate each $`B_i`$ along $`v`$, getting $`B_i^{}`$. Add to each $`B_i^{}`$ a small collar normal to $`v`$, getting $`\mathrm{\Sigma }_i`$ (if $`B_i=\mathrm{}`$, we set $`\mathrm{\Sigma }_i=B_i^{}`$). Select finitely many discs $`\{D_n\}`$ disjoint from each other and from all the $`\mathrm{\Sigma }_i`$’s, such that all positive orbits of $`v`$, except for the small segments between $`B_i^{}`$ and $`B_i`$, meet $`(_{i1}\mathrm{\Sigma }_i)(D_n)`$ in some interior point. Connect the $`D_n`$’s together by strips transversal to $`v`$ and disjoint from $`_{i1}\mathrm{\Sigma }_i`$, getting a disc $`\mathrm{\Sigma }_0`$. Up to a generic small perturbation, the surface $`\mathrm{\Sigma }=_{i0}\mathrm{\Sigma }_i`$ satisfies all axioms of a normal section for $`(M,v)`$, except axiom 6.
Now, even if it is not standard, $`P_\mathrm{\Sigma }`$ can be defined, and the proof of Lemma 2.4 shows that it is a simple branched spine of $`(MB^3,v)`$. In particular, $`P_\mathrm{\Sigma }`$ is connected and its singular locus is non-empty. We recall now that in Chapter 4 of we have considered a set of local moves on simple branched spines, called ‘simple sliding moves’, which preserve the transversal field (and hence the splitting of the boundary), but do not require or preserve the cellularity condition. Knowing that $`P_\mathrm{\Sigma }`$ is connected and $`S(P_\mathrm{\Sigma })\mathrm{}`$, it is not too hard to see that there exists a sequence of (abstract) simple sliding moves which turns $`P_\mathrm{\Sigma }`$ into a standard spine (see , Section 4.4). If we physically realize these moves within $`M`$, preserving transversality to $`v`$, the result is a standard branched spine $`P`$ such that $`C([P])=[M,v]`$.
We are left to show that if $`C([P_0])=C([P_1])`$ then $`P_0`$ and $`P_1`$ are related by sliding-MP-moves and snake moves (‘sliding-equivalent’ for short). By the definition of $`\mathrm{Comb}`$ and $`C`$, using also the above lemmas, there exists a manifold $`M`$ and a homotopy $`(v_t)`$ of concave fields on $`M`$, such that $`P_0`$ and $`P_1`$ are defined by normal sections $`\mathrm{\Sigma }^{(0)}`$ and $`\mathrm{\Sigma }^{(1)}`$ of $`(M,v_0)`$ and $`(M,v_1)`$ respectively.
We prove that $`P_0`$ and $`P_1`$ are sliding-equivalent first in the special case where $`v_0=v_1=v`$. The general case will be an easy consequence. For $`j=0,1`$, let $`\mathrm{\Sigma }^{(j)}=_{i0}\mathrm{\Sigma }_i^{(j)}`$. Proceeding as in the above proof of surjectivity, for each black region $`B_i`$ of $`M`$, we consider a collared negative translate $`\overline{\mathrm{\Sigma }}_i`$ of $`B_i`$. We choose $`\overline{\mathrm{\Sigma }}_i`$ so close to $`B_i`$ that $`\overline{\mathrm{\Sigma }}_i\mathrm{\Sigma }^{(j)}=\mathrm{}`$, and the negative integration of $`v`$ yields a diffeomorphism from $`\overline{\mathrm{\Sigma }}_i`$ to a subset of $`\mathrm{\Sigma }_i^{(j)}`$.
Step I. For $`j=0,1`$, there exists a disc $`D_j`$ such that $`D_j(_{i1}\overline{\mathrm{\Sigma }}_i)`$ is a normal section of $`(M,v)`$, and the associated branched spine is sliding-equivalent to $`P_j`$. To prove this, we temporarily drop the index $`j`$. We first isotope each $`\mathrm{\Sigma }_i`$, without changing the associated spine, until it contains $`\overline{\mathrm{\Sigma }}_i`$, as suggested in Fig. 11.
Note that if $`B_i=\mathrm{}`$ we automatically have $`\mathrm{\Sigma }_i=\overline{\mathrm{\Sigma }}_i`$. Otherwise, we concentrate on one of the annuli $`A`$ of which $`\mathrm{\Sigma }_i\overline{\mathrm{\Sigma }}_i`$ consists. Note that we cannot just shrink $`A`$ leaving the rest of the section unchanged, because we could spoil axiom 4 of the definition of normal section. To actually shrink $`A`$ we first need to “insulate” it, toward the positive direction of $`v`$, by adding to the disc $`\mathrm{\Sigma }_0`$ a strip normal to $`v`$. Figure 12 suggests how to do this.
As we modify $`\mathrm{\Sigma }_0`$ as suggested, it is clear that we keep having a “quasi-normal” section, i.e. all axioms except 6 hold. Moreover the corresponding simple branched spines are obtained from each other by the simple sliding moves already mentioned above. To conclude we apply, as above, the fact that a simple branched spine can be transformed via simple sliding moves to a standard one, and the technical result established in , Proposition 4.5.6, according to which standard spines which are equivalent under simple sliding moves are also sliding-equivalent. This proves Step I.
The conclusion will now follow quite closely the argument in , Theorem 5.2.1.
Step II. There exist discs $`D_0^{}`$ and $`D_1^{}`$ such that $`D_j^{}(_{i1}\overline{\mathrm{\Sigma }}_i)`$ is a normal section of $`(M,v)`$ for $`j=0,1`$, and $`D_0D_0^{}=D_0^{}D_1^{}=D_1^{}D_1=\mathrm{}`$. Choosing a metric on $`M`$, one can construct $`D_0^{}`$ and $`D_1^{}`$ by first taking many very small discs almost orthogonal to $`v`$, and then connecting these discs by strips transversal to $`v`$.
Step III. Conclusion in the case $`v_0=v_1`$. If we connect $`D_0`$ and $`D_0^{}`$ by a strip orthogonal to $`v`$, we get a bigger disc $`\stackrel{~}{D}_0`$ such that $`\stackrel{~}{D}_0(_{i1}\overline{\mathrm{\Sigma }}_i)`$ is still a normal section of $`(M,v)`$. We can actually imagine a dynamical process, in which $`D_0`$ is first enlarged to $`\stackrel{~}{D}_0`$, and then is reduced to $`D_0^{}`$, as in Fig. 13.
If the transformation is chosen generic enough, at all times axioms 123 and 4 will hold, and axiom 5 will hold at all but finitely many times. This means that the corresponding branched spines are related by simple sliding moves. Similarly, we can replace $`D_0^{}`$ first by $`D_1^{}`$ and then by $`D_1`$. Using the facts quoted above, the conclusion follows.
We are left to deal with the general case, where $`(v_t)`$ is a non-constant homotopy. It is then sufficient to take a partition $`0=t_0<t_1<\mathrm{}<t_n=1`$ of $`[0,1]`$, fine enough that $`(M,v_{t_{k1}})`$ and $`(M,v_{t_k})`$ admit a common normal section which gives rise to isomorphic branched spines. 2.1
###### Remark 2.6
Along the lines of the previous proof we have established the following topological fact, whose statement does not involve spines. Let $`(v_t)`$ be a homotopy of concave fields on $`M`$, let $`B_0,B_1M`$ be balls with $`(B_j,v_j)B_{\mathrm{triv}}^3`$ and $`v_j`$ traversing on $`MB_j`$ for $`j=0,1`$. Then there exist another homotopy $`(v_t^{})`$ between $`v_0`$ and $`v_1`$ and an isotopy $`(B_t)`$ with $`(B_t,v_t)B_{\mathrm{triv}}^3`$ and $`v_t`$ traversing on $`MB_t`$ for all $`t`$.
#### Sufficiency of the sliding-MP-moves
To show Proposition 2.3 we will find it convenient to use the graphic representation of branched spines introduced in , Section 3.2, but we do not reproduce here the technicalities needed to introduce this representation.
Proof of2.3. We start by listing rigid spines. Note first that if a negative sliding-MP-move applies to a spine then also a positive one does, so we only need to consider positive rigidity. The spines with one vertex, shown in Fig. 14, are of course rigid.
Using , Proposition 3.3.5, one easily checks that $`(M(P))`$ is $`S_{\mathrm{triv}}^2`$ for the first two spines, and $`S_{\mathrm{triv}}^2S_{\mathrm{triv}}^2`$ for the other two.
Now we turn to rigid spines with more than one vertex. Rigidity implies that all edges with distinct endpoints should appear as on the left in Fig. 15.
It is not hard to deduce that rigid spines come in a sequence $`P_1^{\mathrm{rig}}`$, $`P_2^{\mathrm{rig}}`$, $`\mathrm{}`$ as shown in the rest of Fig. 15, where $`P_k^{\mathrm{rig}}`$ has $`2k`$ vertices, and $`(M(P_k^{\mathrm{rig}}))`$ is the union of $`S_{\mathrm{triv}}^2`$ together with $`k`$ copies of $`S_{\mathrm{white}}^2`$ and $`k`$ copies of $`S_{\mathrm{black}}^2`$. This classification proves (i).
To show (ii) we must prove that:
1. Sequences which contain rigid spines can be replaced by sequences which do not.
2. If two non-rigid spines are related by one snake move then they are also related by a sequence of sliding-MP-moves.
For (ii-a), we note that the result of a positive snake move is never rigid. So if a rigid spine $`P`$ appears in a sequence of moves then $`P`$ is the result of a negative snake move $`\mu _1^1:P_1P`$, and a positive snake move $`\mu _2:PP_2`$ is applied to $`P`$. Since all edges of a spine survive through a snake move, there is a version $`\stackrel{~}{\mu }_2`$ of $`\mu _2`$ which applies to $`P_1`$ and a version $`\stackrel{~}{\mu }_1`$ of $`\mu _1`$ which applies to $`P_2`$, and the result $`\stackrel{~}{P}`$ is the same. So can replace the segment $`(P_1,P,P_2)`$ by $`(P_1,\stackrel{~}{P},P_2)`$, and now all the spines involved are non-rigid.
Let us turn to (ii-b). The proof results from three steps, to describe which we introduce in Figure 16
another move, called sliding-vertex move, whose unbranched version was already considered in and . Again, taking into account orientations, there are two versions of the move (for each vertex type), but we will ignore this detail.
Step 1: if $`v`$ is a vertex of a branched spine $`P`$, $`e`$ is any one of the edges incident to $`v`$, $`P_v`$ is obtained from $`P`$ via the sliding-vertex move at $`v`$, and $`P_e`$ is obtained from $`P`$ via the snake move on $`e`$, then $`P_v`$ and $`P_e`$ are related by sliding-MP-moves. This is proved by an easy case-by-case analysis. It turns out that two MP-moves (a positive and a negative one) are always sufficient.
Step 2: let $`v`$, $`P`$ and $`P_v`$ be as above. If $`P`$ and $`P_v`$ are related by sliding-MP-moves, the same is true for $`P`$ and any spine obtained from $`P`$ by a snake move. To see this, use step 1 to successively transform sliding-vertex moves into snake moves and conversely, until the desired snake move is reached.
Step 3: if $`P`$ is non-rigid then there exists a vertex $`v`$ such that $`P`$ and $`P_v`$ are related by MP-moves. The vertex $`v`$ is chosen to be an endpoint of an edge to which the positive MP-move applies. The argument is again a long case-by-case one, which refines in a branched context the argument given by Piergallini in . The sequence always consists of three positive moves followed by a negative one. This concludes the proof of (ii), whereas (iii) is evident. 2.3
## 3 A calculus for framed links
We fix in this section a compact manifold $`M`$ and consider the set $`\mathrm{Fram}(M)`$ of isotopy classes of framed links in $`M`$. Since a link isotopy generically avoids a fixed 3-ball, $`\mathrm{Fram}(M)`$ and $`\mathrm{Fram}(\widehat{M})`$ are canonically isomorphic when $`M=S^2`$, so we can restrict to non-closed $`M`$’s and include the closed case as usual.
#### Statement
Let us fix a branched standard spine $`P`$ of $`M`$. The fact that such a spine always exists was proved as Theorem 3.4.9 in . We call $`\mathrm{C}^1`$ link diagram on $`P`$ an immersion of a disjoint union of circles into $`P`$, with generic intersection with $`S(P)`$ appearing as in Fig. 17,
generic self-intersections (crossings), and the usual under-over marking at crossings (as shown in the same figure; here ‘under’ and ‘over’ refer to the field positively transversal to $`P`$). The set of all $`\mathrm{C}^1`$ link diagrams on $`P`$ will be denoted by $`𝒟(P)`$. An element of $`𝒟(P)`$ obviously defines a link. Moreover $`v(P)`$ is transversal to this link, so it defines a framing, and we get an (obviously well-defined) map $`F_{(P,M)}:𝒟(P)\mathrm{Fram}(M)`$. Besides isotopy on $`P`$ through immersions having the same configuration of crossings and intersections with $`S(P)`$, there are several combinatorial moves which of course do not modify the isotopy class of the framed link defined by a diagram. We show a list of moves having this property in Fig. 18.
For a reason to be given below, which also explains the apparently weird notation, we call these moves $`\mathrm{C}^1`$-Turaev moves. As we did when we described the sliding-MP-moves in Fig. 6, we are showing in Fig. 18 only the essential physical modifications, without specifying the screw-orientation of the spine and the orientation of its regions.
###### Theorem 3.1
The map $`F_{(P,M)}:𝒟(P)\mathrm{Fram}(M)`$ is surjective, and the equivalence relation it defines is generated by $`\mathrm{C}^1`$-Turaev moves.
Our argument, after the easy proof of surjectivity, goes along the following lines:
1. We state the analogue of Theorem 3.1 for non-branched spines, due to Turaev (we include a quick proof for the sake of completeness);
2. We modify Turaev’s result to the case of a branched spine, but allowing non-$`\mathrm{C}^1`$ diagrams;
3. We prove our theorem, showing how to canonically replace each non-$`\mathrm{C}^1`$ diagram by a $`\mathrm{C}^1`$ one along a sequence of modifications.
#### Surjectivity
Since $`M`$ and $`P`$ are fixed, we write $`F`$ for short. Given a framed link $`L`$ in $`M`$, we can prove that it is contained in the image of $`F`$ as follows:
* First, forget the framing and take a generic projection on $`P`$, recalling that $`MPM\times (0,1]`$;
* Next, eliminate non-$`\mathrm{C}^1`$ intersections with $`S(P)`$ as shown in Fig. 19 (left).
* Finally, give the resulting projection the right framing by adding the necessary numbers of curls. (Here we use the fact that two framings on a given knot differ at most by a finite number of full rotations.)
It may be noted that surjectivity of $`F`$ is preserved by restriction to the set of diagrams without crossings. This follows quite easily from the fact that all the regions of $`P`$ have non-empty boundary, as suggested in Fig. 19 (right). This property will not be used below.
#### Turaev moves on standard spines
The ideas and results of this paragraph are due to Turaev . We temporarily allow $`P`$ to be any standard spine of $`M`$, not a branched one. If each region of $`P`$ is given an arbitrary transverse orientation, the definition of a link diagram $`D`$ makes sense also on $`P`$, but $`D`$ may not define a framing on the associated link, because the strip which runs along a component of $`D`$ on $`P`$ need not be a cylinder, it may be a Möbius strip. So we attach to each component $`D_i`$ of $`D`$ a full or half-integer $`a_i/2`$, depending on the topology of the strip, and we define the framing by giving $`a_i`$ positive half-twists to the strip (recall that $`M`$ is oriented). By diagram on $`P`$ we will actually mean one such pair $`(\{D_i\},\{a_i/2\})`$.
We call Turaev moves those shown in Fig. 20,
together with the $`\mathrm{R}_{\mathrm{I}\mathrm{I}}`$ and $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$ already shown above. In Fig. 20, for $`\mathrm{R}_\mathrm{I}^{}`$ and $`\mathrm{T}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$, the local orientation must be that of $`\text{}^3`$.
###### Theorem 3.2
Every isotopy class of framed link in $`M`$ is defined by some diagram on $`P`$, and two diagrams define the same class if and only if they are obtained from each other by a sequence of Turaev moves.
Proof of3.2. Recall that a framed link can be thought of as an embedded cylinder. Moreover $`M`$ projects onto $`P`$, and the projection of a cylinder generically appears as in Fig. 21.
Such a projection easily defines a diagram, the half-integers being sums of $`\pm 1/2`$’s corresponding to the bends of the projection. Moreover, the elementary catastrophes along an isotopy of a projection translate into the moves of Fig. 20, or simple combinations of them. 3.2
#### Turaev moves on branched spines
Going back to the case where $`P`$ is branched, we can still apply Theorem 3.2, but now the list of moves becomes slightly longer, if we want to take the branching into account.
###### Proposition 3.3
If $`P`$ is a branched spine then any Turaev move for a diagram on $`P`$ can be expressed as a combination (including inverses) of the moves $`\mathrm{R}_\mathrm{I}^{}`$, $`\mathrm{R}_{\mathrm{I}\mathrm{I}}`$, $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$, $`\mathrm{T}_\mathrm{I}^{}`$, $`\mathrm{T}_{\mathrm{I}^{\prime \prime }}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}^{}}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{V}^{}}`$, $`\mathrm{T}_\mathrm{V}^{}`$ shown above, together with the moves $`\mathrm{T}_{\mathrm{I}^{\prime \prime \prime }}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}^{\prime \prime }}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}\mathrm{I}^{}}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{V}^{\prime \prime }}`$, $`\mathrm{T}_{\mathrm{V}^{\prime \prime }}`$ shown in Fig. 22
Proof of3.3. The branching can be interpreted as a loss of symmetry of a spine, so each of Turaev’s moves, when viewed as a move on a branched spine, generates many different ones according to the position of the diagram with respect to the branching. The result is a list much longer than that given in the statement, but one can show that all the moves omitted from the statement
are generated by the moves included. Two examples are provided in Figg. 23
and 24. 3.3
#### $`𝐂^\mathrm{𝟏}`$ moves
We can now conclude the proof of Theorem 3.1 (surjectivity having been shown above). We are left to show that if $`D`$ and $`D^{}`$ are $`\mathrm{C}^1`$-diagrams on $`P`$ which define the same framed link, then they are related by a sequence of $`\mathrm{C}^1`$ Turaev moves. By Theorem 3.2 and Proposition 3.3, there exists a sequence $`D=D_0D_1\mathrm{}D_{n1}D_n=D^{}`$ where each move $`D_{ii}D_i`$ is one of those listed in Proposition 3.3. In particular the $`D_i`$’s with $`0<i<n`$ can be non-$`\mathrm{C}^1`$ and can have a non-zero half-integer attached to them. We will now show how to construct a modified sequence $`D=\stackrel{~}{D}_0\stackrel{~}{D}_1\mathrm{}\stackrel{~}{D}_{n1}\stackrel{~}{D}_n`$ with the following properties:
1. each $`\stackrel{~}{D}_i`$ is a $`\mathrm{C}^1`$ diagram with number 0 attached;
2. each $`\stackrel{~}{D}_i`$ is obtained from $`\stackrel{~}{D}_{i1}`$ by a sequence of $`\mathrm{C}^1`$ Turaev moves;
3. each $`\stackrel{~}{D}_i`$ differs from $`D_i`$ for the presence of some extra curls; in particular each component $`\stackrel{~}{D}_i^{(j)}`$ of $`\stackrel{~}{D}_i`$ has a natural companion $`D_i^{(j)}`$ in $`D_i`$, with the property that, as unframed knots, the knots associated to $`\stackrel{~}{D}_i^{(j)}`$ and $`D_i^{(j)}`$ are both contained in a solid torus $`T_i^{(j)}`$ and parallel to the core of the torus;
4. the framed knots associated to $`\stackrel{~}{D}_i^{(j)}`$ and $`D_i^{(j)}`$ are framed-isotopic within $`T^{(j)}`$.
Requirement 4 is the crucial technical point of our proof. To verify that the requirement is stronger than just framed isotopy, note that in $`D^2\times S^1`$ the framings on the core $`\{0\}\times S^1`$ are parameterized by the integers, but, when $`D^2\times S^1`$ is mirrored in its boundary to get $`S^2\times S^1`$, only two inequivalent framings remain (corresponding to even and odd integers).
We assume for a moment the sequence $`\stackrel{~}{D}_i`$ to exist, and we show how to conclude. The transformation from $`D`$ to $`\stackrel{~}{D}_n`$ is made with $`\mathrm{C}^1`$ Turaev moves, so we only need to compare $`\stackrel{~}{D}_n`$ and $`D^{}=D_n`$, which by assumption differ for some curls. Using $`\mathrm{C}^1`$ Turaev moves we can easily make all these curls slide until they are consecutive on the diagram. We recall now that there are four local pictures for a curl, depending on its local contributions $`\pm `$ to the framing and to the winding number . Assumption 4 now implies that the algebraic sum of local contributions to the framing vanishes. Therefore we can cancel the curls in pairs, either by moves $`\mathrm{R}_\mathrm{I}^{}`$ (when the local contributions to the winding number are the same), or by a combination of moves $`\mathrm{R}_{\mathrm{I}\mathrm{I}}`$ and $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$ (when the contributions cancel). This shows the conclusion.
We are left to define the sequence $`\stackrel{~}{D}_i`$. The idea is simply not to perform the moves which change the half-integer or introduce cusps, and show that the sequence of moves can be followed anyway. While doing this we need to keep track of the portions where the new diagram $`\stackrel{~}{D}_i`$ differs from $`D_i`$, which we do by marking a neighbourhood of the portion as a shadowed box.
The moves which change the colour or introduce cusps are $`\mathrm{R}_\mathrm{I}^{}`$ (in both directions), $`\mathrm{T}_{\mathrm{I}^{\prime \prime \prime }}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}\mathrm{I}^{}}`$, and $`\mathrm{T}_{\mathrm{V}^{\prime \prime }}`$, and we show in Figg. 25 and 26
what we replace them with. For move $`\mathrm{R}_\mathrm{I}^{}`$ it has been necessary to be more
specific because, in the original definition of the move, two different ones were actually defined at the same time.
To show that $`\{\stackrel{~}{D}_i\}`$ can indeed be constructed we must now show that after performing the construction up to some level $`k`$ we can still still follow the rest of the sequence and go on with the replacements of Figg. 25 and 26. By construction $`\stackrel{~}{D}_k`$ differs from $`D_k`$ only within some shadowed boxes. We denote by $`\mu _k`$ the move $`D_kD_{k+1}`$ and explain how to lift it to a move on $`\stackrel{~}{D}_k`$. First, note that a shadowed box lying within a region of $`P`$ and containing a curl does not interfere with $`\mu _k`$ whatever its type (but it may be necessary to add some $`\mathrm{R}_{\mathrm{I}\mathrm{I}}`$’s and $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$’s to replace isotopy supported within the region).
We fix now our attention on a shadowed box $`B`$ which lies on $`S(P)`$ and contains a smoothed cusp, and examine the various instances for $`\mu _k`$ with respect to $`B`$. If $`\mu _k`$ is of type $`\mathrm{R}_{\mathrm{I}\mathrm{I}}^{\pm 1}`$, $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}^{\pm 1}`$, $`\mathrm{T}_\mathrm{I}^{}^{\pm 1}`$, $`\mathrm{T}_{\mathrm{I}^{\prime \prime }}^{\pm 1}`$ $`\mathrm{T}_{\mathrm{I}\mathrm{I}^{}}^{\pm 1}`$, or $`\mathrm{T}_{\mathrm{I}\mathrm{V}^{}}^{\pm 1}`$, then it obviously does not interfere with $`B`$, so we can just perform $`\mu _k`$ on $`\stackrel{~}{D}_k`$. If $`\mu _k`$ is one of the moves $`\mathrm{R}_\mathrm{I}^{\pm 1}`$, $`\mathrm{T}_{\mathrm{I}^{\prime \prime \prime }}`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}\mathrm{I}^{}}`$ or $`\mathrm{T}_{\mathrm{V}^{\prime \prime }}`$ then again it does not interfere with $`B`$, and we can perform the appropriate replacement from Figg. 25 or 26, getting a move from $`\stackrel{~}{D}_k`$ to $`\stackrel{~}{D}_{k+1}`$. If $`\mu _k`$ is a move of type $`\mathrm{T}_{\mathrm{I}\mathrm{I}^{\prime \prime }}^{\pm 1}`$ or $`\mathrm{T}_{\mathrm{I}\mathrm{V}^{\prime \prime }}^{\pm 1}`$ then it may interfere with $`B`$. However, since it does not create or destroy cusps, $`\mu _k`$ can be translated on $`\stackrel{~}{D}_k`$ as a combination of moves which do not involve cusps, i.e. allowed from the statement. We are only left to deal with the case where $`\mu _k`$ is one of the moves $`\mathrm{T}_{\mathrm{I}^{\prime \prime \prime }}^1`$, $`\mathrm{T}_{\mathrm{I}\mathrm{I}\mathrm{I}^{}}^1`$ or $`\mathrm{T}_{\mathrm{V}^{\prime \prime }}^1`$ which destroy cusps. By construction the cusp(s) to be destroyed still appear in $`\stackrel{~}{D}_k`$ as smoothed cusps within shadowed boxes, so we can destroy them also from $`\stackrel{~}{D}_k`$ by means of allowed moves.
The crucial properties 3 and 4 hold by construction, and the conclusion eventually follows. 3.1
## 4 A calculus for pseudo-Legendrian links <br>in combed manifolds
We will deal in this section with the set $`\mathrm{PLeg}`$ of equivalence classes of triples $`(M,v,L)`$ already described in the introduction. Its combinatorial counterpart will be given by the set
$$=\{(P,D):P,D𝒟(P)\}$$
where $``$ is as in Section 2 and $`𝒟(P)`$ is as in Section 3. The reconstruction map $`(P,D)L(P,D)`$ is here defined by noting that $`D`$ defines a link transversal to $`v(P)`$ in $`M(P)`$, and hence also a link transversal to $`\widehat{v}(P)`$ in $`\widehat{M}(P)`$. According to what we stated after Proposition 2.3 we will actually drop from $``$ the sliding-MP-rigid spines, and ignore the snake move.
If for a fixed $`P`$ we consider the effect on $`L(P,D)`$ of the $`\mathrm{C}^1`$-Turaev moves on $`D`$, we see that the class of $`L(P,D)`$ is in general modified by the first Reidemeister move $`\mathrm{R}_\mathrm{I}^{}`$ of Fig. 18 (see also Section 5), but not by the other moves, which we will therefore call pseudo-Legendrian Turaev moves. Other moves which obviously do not change $`L(P,D)`$ up to equivalence are the sliding-MP-moves on $`P`$ which do not involve $`D`$ (these moves permit to follow $`D`$ along the modification of $`P`$, so they are well-defined for pairs). It is not hard to see that before performing a sliding-MP-move on $`P`$ it is always possible to modify $`D`$ by pseudo-Legendrian Turaev moves to a diagram which is not involved in the sliding-MP-move, and the diagram after the sliding-MP-move is well-defined up to pseudo-Legendrian Turaev moves. For this reason we will freely speak of sliding-MP-moves also for pairs. The following will be established below.
###### Theorem 4.1
The map $`L:\mathrm{PLeg}`$ is surjective, and the equivalence relation defined by $`L`$ is generated by pseudo-Legendrian Turaev moves and sliding-MP-moves.
#### Fixed-spine statement
Recall from the definition that the map $`L`$ of the statement of Theorem 4.1 involves the passage from $`P`$ to $`(M(P),v(P))`$ and then to $`(\widehat{M}(P),\widehat{v}(P))`$. As already pointed out, this is necessary if one wants to be able to deal with non-traversing fields. However, if one happens to have a concave traversing field, one can directly encode this field by a spine, without first removing a ball, and one can investigate how isotopy of links transversal to the field reflects on link diagrams on the spine. The following is shown below:
###### Proposition 4.2
Let $`P`$ be a branched spine. Fix a representative of $`(M(P),v(P))`$ and an embedding of $`P`$ in $`M(P)`$ transversal to $`v(P)`$. Then every link transversal to $`v(P)`$ is represented by a $`\mathrm{C}^1`$ diagram on $`P`$. Moreover two $`\mathrm{C}^1`$ diagrams define the same link up to isotopy through links transversal to $`v(P)`$ if and only if they are related by pseudo-Legendrian Turaev moves.
#### From fixed to variable spine
We show in this paragraph how to deduce Theorem 4.1 from Proposition 4.2. First of all, to prove surjectivity, we consider a triple $`(M,v,L)`$ representing an element of $`\mathrm{PLeg}`$. Using a normal section as in Proposition 2.4, we can obtain a spine $`P`$ which encodes the equivalence class of $`(M,v)`$ in the sense of Theorem 2.1. Moreover $`P`$ comes with an embedding in $`M`$ transversal to $`v`$. Now, a neighbourhood of $`P`$ can be identified to $`M(P)`$ and its complement is isomorphic to $`B_{\mathrm{triv}}^3`$. Using the flow generated by $`v`$ in this ball we can now isotope $`L`$ through links transversal to $`v`$ to a link which lies in $`M(P)`$, and the first assertion of Proposition 4.2 implies that $`L`$ is represented by a diagram $`D`$ on $`P`$. Summing up, we see that $`(M,v,L)`$ is represented by $`(P,D)`$, and surjectivity of $`L`$ is proved.
To conclude we must now show that two pairs $`(P_0,D_0)`$ and $`(P_1,D_1)`$ are equivalent via pseudo-Legendrian Turaev moves when $`L(P_0,D_0)=L(P_1,D_1)`$. Spelling out the relation of pseudo-Legendrian isotopy, which defines $`\mathrm{PLeg}`$, we assume that $`P_0`$ and $`P_1`$ embed in the same manifold $`M`$ and that there exist a field homotopy $`(v_t)_{t[0,1]}`$ and a link isotopy $`(L_t)_{t[0,1]}`$ on $`M`$ such that:
1. for $`i=0,1`$, the link $`L_i`$ is the one defined by $`D_i`$, and the field $`v_i`$ is positively transversal to $`P_i`$ and restricts to $`B_{\mathrm{triv}}^3`$ on the complement of $`P_i`$;
2. $`L_t`$ is transversal to $`v_t`$ for all $`t`$.
Given $`t[0,1]`$, we note that for $`|ts|1`$ the link $`L_s`$ is transversal to $`v_t`$, and that a branched spine for $`v_t`$ (in the sense repeatedly used above) is a branched spine also for $`v_s`$. So we can subdivide $`[0,1]`$ into subintervals $`[t_{i1},t_i]`$ so that:
1. $`L_{t_{i1}}`$ is isotopic to $`L_{t_i}`$ through links transversal to $`v_{t_i}`$;
2. $`v_{t_i}`$ has a spine $`P_{t_i}`$ which is also a spine for $`v_{t_{i1}}`$.
Now let $`D_{t_i}`$ be a diagram for $`L_{t_i}`$ on $`P_{t_i}`$. Since both $`P_{t_{i1}}`$ and $`P_{t_i}`$ are spines for $`v_{t_i}`$, we can transform $`P_{t_{i1}}`$ into $`P_{t_i}`$ via sliding-MP-moves. Using pseudo-Legendrian Turaev moves we can now follow $`D_{t_{i1}}`$ along this sequence of moves, getting a diagram $`D_{t_i}^{}`$ on $`P_{t_i}`$. Since the sequence of sliding-MP-moves can be realized in $`M`$ so that each spine of the sequence is a branched spine for $`v_{t_i}`$, we deduce that the link defined by $`D_{t_i}^{}`$ is isotopic to $`L_{t_{i1}}`$, and hence to $`L_{t_i}`$, through links transversal to $`v_{t_i}`$. Proposition 4.2 now implies that $`D_{t_i}^{}`$ and $`D_{t_i}`$ are related by pseudo-Legendrian Turaev moves on $`P_{t_i}`$. This shows that $`(P_{t_i},D_{t_i})`$ is obtained from $`(P_{t_{i1}},D_{t_{i1}})`$ via the moves of the statement, and the conclusion follows by iteration. 4.1
#### Fixed-spine proof
We will establish now Proposition 4.2, writing just $`M`$ and $`v`$ for $`M(P)`$ and $`v(P)`$. For the sake of simplicity we will assume that $`v`$ is tangent to $`M`$ along only one curve (denoted by $`\gamma `$), but our arguments extends almost verbatim to the general case of more than one curve.
We fix in $`M`$ an annulus $`A`$ which connects $`\gamma `$ to $`S(P)`$ as shown in a cross-section in Fig. 27.
Note that $`A`$ is almost but not quite embedded: it has double point at the vertices of $`P`$. Since we will only need to consider $`A`$ locally and away from vertices of $`P`$, this fact will not disturb us. We choose coordinates $`(\rho ,\theta )[0,1]\times [0,2\pi ]`$ on $`A`$, where $`\rho =0`$ corresponds to $`S(P)`$ and $`\rho =1`$ to $`\gamma `$. Near $`A`$ we can also define a coordinate $`z[\epsilon ,\epsilon ]`$ by integrating $`v`$.
Now let $`L`$ be transversal to $`v`$, and assume by general position that $`L`$ intersects $`A`$ only at points with $`0<\rho <1`$, that no two such intersections have the same coordinate $`\theta `$, and that at all the intersections the tangent direction to $`L`$ has non-zero components in all three coordinates $`\rho ,\theta ,z`$. Depending on the sign of these components, we can divide the points of $`LA`$ into four types, shown in Fig. 28
($`L`$ is dashed when it lies over $`A`$ and dotted when it lies under $`A`$). We consider now the projection $`\pi `$ of $`MA`$ onto $`P`$ along the orbits of $`v`$, as shown in Fig. 29.
Of course $`LA`$ locally projects to a $`\mathrm{C}^1`$-strand on $`P`$, and by general position we can assume that $`\pi (LA)`$ locally appears as a $`\mathrm{C}^1`$-diagram. We are only left to extend the diagram at the points of $`LA`$, which we do locally in Fig. 30.
The top part of this figure actually refers to a simplified situation, because other strands of $`L`$ already projected on $`P`$ may locally interfere. We show at the bottom of the same figure in one example how to deal with this fact. The resulting diagram of course represents $`L`$, and we have proved the first assertion in Proposition 4.2. To prove the second assertion we must now examine an isotopy of $`L`$ through links transversal to $`v`$, and hence examine first-order violations of genericity of $`L`$ with respect to $`A`$ and $`\pi `$. All the elementary accidents which do not involve $`A`$ of course correspond to pseudo-Legendrian Turaev moves. We are left to deal with the following accidents:
1. $`L`$ intersects $`A`$ at a point of $`S(P)A`$;
2. at a point of $`LA`$, the tangent direction to $`L`$ has vanishing $`\rho `$-coordinate;
3. similarly, with the $`\theta `$-coordinate;
4. similarly, with the $`z`$-coordinate.
In Fig. 31 we show the situation just before and just after each of these accidents, and
we analyze the corresponding transformations of the diagrams constructed as in Fig. 30. In all cases one easily sees that indeed the transformation is generated by pseudo-Legendrian Turaev moves: the number of moves needed is respectively one, zero (isotopy within regions), three and two. By simplicity in Fig. 31 we have ignored the possible interference of other strands of $`L`$ projected on $`P`$, but the conclusion is valid anyway (some Reidemeister moves must be added in the general case). 4.2
## 5 Applications and speculations
In this section we will discuss some consequences of the calculi described above, and mention some natural questions and problems which we put forward for further investigation. The section is split into two subsections.
### 5.1 Winding number, torsion, and finite-order invariants
In this section we employ our realizations of $`\mathrm{Fram}`$ and $`\mathrm{PLeg}`$ in connection with winding number, Maslov index, torsion, and finite-order invariants of pseudo-Legendrian knots.
#### Relative winding number
We spell out in this paragraph the analogue of Trace’s result on knot diagrams in $`\text{}^3`$. We confine ourselves to knots for the sake of simplicity, but essentially the same holds for links.
###### Proposition 5.1
Let $`(v_0,K_0)`$ and $`(v_1,K_1)`$ be pseudo-Legendrian pairs in a manifold $`M`$, where $`K_0`$ and $`K_1`$ are oriented knots. Assume that $`v_0`$ and $`v_1`$ are homotopic relatively to $`M`$, and that $`K_0`$ and $`K_1`$ are isotopic as oriented framed knots. Then, up to pseudo-Legendrian isotopy on $`(v_1,K_1)`$, we can assume that $`v_1=v_0`$ and that $`K_1`$ differs from $`K_0`$ only within a region of $`(M,v_0)`$ isomorphic to $`(\text{}^3,/z)`$, where $`K_0`$ is a straight horizontal line and $`K_1`$ has either some positive or some negative double curls (shown in Fig. 32).
Proof of5.1. Choose branched spines $`P_0`$ and $`P_1`$ of $`v_0`$ and $`v_1`$ according to Proposition 2.4, and use Proposition 4.2 to represent $`K_0`$ and $`K_1`$ on $`P_0`$ and $`P_1`$ respectively by $`\mathrm{C}^1`$-diagrams $`D_0`$ and $`D_1`$. Since $`v_0`$ and $`v_1`$ are homotopic, a sequence of sliding-MP moves connects $`P_1`$ to $`P_0`$. Following $`D_1`$ along this sequence of moves we get a pseudo-Legendrian isotopy, so we can assume that $`v_1=v_0`$ and $`P_1=P_0`$. Now $`D_0`$ and $`D_1`$ define on $`P_0`$ framed-isotopic oriented knots, so by Theorem 3.1 they are related by $`\mathrm{C}^1`$-Turaev moves. If along the sequence of moves there is no $`\mathrm{R}_\mathrm{I}^{}`$, we deduce pseudo-Legendrian equivalence of $`K_0`$ and $`K_1`$. If however there is some $`\mathrm{R}_\mathrm{I}^{}`$, we replace it by a pseudo-Legendrian move as shown in Fig. 33.
One easily sees that this replacement can be done consistently along the sequence of moves. The result is a pseudo-Legendrian isotopy between $`D_1`$ and a diagram which differs from $`D_0`$ only for some double curls. These double curls can of course be slid to be consecutive. Now, it is precisely the content of that up to moves $`\mathrm{R}_{\mathrm{I}\mathrm{I}}`$ and $`\mathrm{R}_{\mathrm{I}\mathrm{I}\mathrm{I}}`$ there are only the types of double curls shown in Fig. 32, and that a positive and a negative double curl cancel out. 5.1
Under the assumptions of the previous proposition one may be tempted to define a relative winding number $`w(K_1,K_0)`$ as the (algebraic) number of double curls by which the diagram of $`K_1`$ differs from the diagram of $`K_0`$. This number is however not well-defined in general, as one easily sees in $`S^2\times S^1`$ with vector field parallel to the $`S^1`$-factor, because in this case a double curl on a diagram contained in $`\text{}^2=S^2\{\mathrm{}\}`$ can always be removed by isotoping the diagram through $`\mathrm{}`$. This seems to suggest that not $`w(K_1,K_0)`$, but maybe $`w(K_1,K_0)[\mu _{K_0}]H_1(E(K_0);\text{})`$ is well-defined, where $`E(K_0)`$ is the exterior of $`K_0`$ and $`\mu _{K_0}`$ is the meridian. We will show this fact under the additional assumption that $`K_0`$ is ‘good’ (see and below for explanations).
#### Torsion invariants and good knots
In we have defined the Reidemeister-Turaev torsion of an Euler structure with simple boundary, and we have applied this notion to define the torsion of pseudo-Legendrian knots. As an absolute invariant torsion contains a sort of lift of the classical Alexander invariant. We will discuss in this paragraph the information carried by torsion as a relative invariant of two pseudo-Legendrian framed-isotopic knots $`K_0,K_1`$ in the same concave combed manifold $`(M,v)`$. We recall from that this information is most easily expressed when $`K_0`$ has the property of being good. Goodness depends only on the isotopy class of $`K_0`$ as a framed knot, and it means that a certain quotient of the mapping class group of $`E(K_0)`$ acts trivially on the space of Euler structures on $`E(K_0)`$. We omit the precise definition here, but we recall that many knots indeed are good (for instance, all are good if $`M`$ is a homology sphere, and most hyperbolic knots are good).
When $`K_0`$ is good, the information carried by torsion as a relative invariant depends only on $`\alpha (v|_{E(K_0)},f(v|_{E(K_1)})H_1(E(K_0);\text{})`$, where $`f\mathrm{Diff}_0(M)`$ maps $`K_1`$ to $`K_0`$ as framed knots, and $`\alpha `$ is the first obstruction for two vector fields to be homotopic relative to the boundary. So the next result means that for good knots the relative winding number gives a well-defined invariant, and all the information torsion can capture is contained in the relative winding number. The statement involves all the assumptions and notations of the present and previous paragraph.
###### Proposition 5.2
$`\alpha (v|_{E(K_0)},f(v|_{E(K_1)})=w(K_1,K_0)[\mu _{K_0}]H_1(E(K_0);\text{}).`$
Proof of5.2. We note first that $`w(K_1,K_0)`$ and $`[\mu _{K_0}]`$ depend on the choice of an orientation on $`K_0`$, but their product does not, so the statement makes sense. For the proof, note that by goodness we can just assume that $`K_0`$ and $`K_1`$ differ as in the statement of Proposition 5.1, and that $`f`$ is supported on a neighbourhood of the region where $`K_0`$ and $`K_1`$ differ. The conclusion then follows directly from Proposition 2.17 of . 5.2
Even if we have not discussed finite-order invariants yet, we note here that Proposition 5.2 implies that torsion is a weaker invariant than the finite-order ones for Legendrian knots in a given homotopy class of Legendrian immersions. To our knowledge the only known examples of framed-isotopic knots distinguished by such invariants are those due to Tchernov , and we believe that they are all good (at least, they certainly are good when the ambient manifold is $`S^2\times S^1`$). Now one sees that in all of Tchernov examples $`w(K_1,K_0)[\mu _{K_0}]=0`$, so torsion definitely cannot distinguish. On the other hand, the definition of torsion does not require fixing a homotopy class of Legendrian immersions, so torsion and finite-order invariants are in some sense complementary.
We will state in the rest of this paragraph some interesting consequences of Proposition 5.2, always assuming the knots involved to be good. For simplicity, as in Proposition 5.2, we stick to knots transverse to a given field $`v`$ on a given $`M`$, but we remind that the relation of pseudo-Legendrian isotopy also involves a homotopy of $`v`$.
###### Corollary 5.3
Under the same assumptions as in Proposition 5.2, suppose furthermore that $`[\mu _{K_0}]`$ has infinite order in $`H_1(E(K_0);\text{})`$, so $`w(K_1,K_0)\text{}`$ is well-defined. Then the following facts are pairwise equivalent:
1. $`w(K_1,K_0)=0`$;
2. $`K_0`$ and $`K_1`$ have trivial relative torsion invariants;
3. $`K_0`$ and $`K_1`$ are pseudo-Legendrian isotopic.
Proof of5.3. Equivalence of (1) and (3) follows from the definition of $`w(K_1,K_0)`$. Implication (1)$``$(2) follows from Proposition 5.2, and the opposite implication follows by taking the torsion associated to a representation $`\phi `$ of $`H_1(E(K_0);\text{})`$ such that $`\phi ([\mu _{K_0}])`$ has infinite order (see for details). 5.3
If $`M`$ is a homology sphere and $`K`$ is a pseudo-Legendrian knot in $`(M,v)`$ we have shown in that the rotation number $`\mathrm{rot}_v(K)`$, also called Maslov index, can be defined just as in the case where $`K`$ is Legendrian in a contact structure. Now:
###### Lemma 5.4
If $`M`$ is a homology sphere, $`K_0`$ and $`K_1`$ are pseudo-Legendrian in $`(M,v)`$ and framed isotopic, then $`w(K_1,K_0)=\frac{1}{2}(\mathrm{rot}_v(K_1)\mathrm{rot}_v(K_0))`$.
(Concerning the statement, note that $`\mathrm{rot}_v(K_1)\mathrm{rot}_v(K_0)`$ must be even if $`K_0`$ and $`K_1`$ are framed-isotopic, otherwise one of $`\{K_0,K_1\}`$ would lift to a closed path in a spin structure on $`M`$, and the other one would not: a contradiction. A proof is easily obtained by isotoping $`K_1`$ as stated in Proposition 5.1.)
Lemma 5.4 gives another proof of the fact that $`w(K_1,K_0)\text{}`$ can be defined when $`M`$ is a homology sphere. Moreover, it could be used to show goodness of knots in a homology sphere by a more direct argument than that given in . We conclude this paragraph by showing the result stated in the introduction and asking a question which naturally arises from it.
Proof of0.1. Equivalence of (1) and (2) comes from Corollary 5.3. Equivalence of (2) and (3) comes from Lemma 5.4. Equivalence of (3) and (4) follows from Corollary5.3 and the fact that the first homology group of the complement of a knot in $`M`$ is infinite cyclic and generated by a meridian. Equivalence of (3) and (5) is an application of Gromov’s $`h`$-principle (see ). 0.1
###### Question 5.5
Let $`(M,v)`$ be an arbitrary combed manifold, let $`K_0`$ and $`K_1`$ be pseudo-Legendrian in $`(M,v)`$ and framed-isotopic, and assume that they are homotopic through pseudo-Legendrian immersions. Does this imply that $`w(K_1,K_0)[\mu _{K_0}]=0`$? (We do not think that the opposite implication can be true in general, in particular when $`[\mu _{K_0}]`$ has finite order.)
#### Absolute winding number
We concentrate in this paragraph on fields $`v`$ such that $`(v^{})=0`$, where $``$ denotes the Euler class and the choice of the metric is of course immaterial. Condition $`(v^{})=0`$ is equivalent to the existence of another non-vanishing field $`x`$ always transversal to $`v`$. Since the ambient manifold is oriented, this is also equivalent to the fact that $`v`$ extends to a framing $`(v,x,y)`$, i.e. a global trivialization of the tangent bundle to $`M`$. Assume now that $`K`$ is an oriented knot transversal to $`v`$. Then, taking the projection of the tangent vector to $`K`$ on the unit sphere of the $`(x,y)`$-plane, and computing the degree, we can define a rotation number $`\mathrm{rot}_{(v,x)}(K)`$.
###### Remark 5.6
$`\mathrm{rot}_{(v,x)}(K)`$ is invariant under simultaneous homotopy $`(v_t,x_t)`$ and isotopy $`(K_t)`$ such that $`x_t`$ and $`v_t`$ are transversal to $`v_t`$ for all $`t`$. Moreover $`\mathrm{rot}_{(v,x)}(K)`$ is independent of $`x`$ when $`M`$ is a homology sphere, and it equals the Maslov index already discussed above.
Assume now that $`K_0`$ and $`K_1`$ are both transversal to $`v`$. Within the proof of Proposition 5.1 we have shown that $`K_1`$ can be isotoped through knots transversal to $`v`$ to a knots which differs from $`K_0`$ by double curls only.
###### Remark 5.7
$`\mathrm{rot}_{(v,x)}(K_1)\mathrm{rot}_{(v,x)}(K_0)`$ is independent of $`x`$ and equals twice the number of double curls by which $`K_0`$ and $`K_1`$ differ, up to isotopy transversal to $`v`$.
The previous remark shows that the relative winding number is well-defined (without any assumption on the knots) if one restricts to knots transversal to a given $`v`$ with $`(v^{})=0`$, and one views the knots up to isotopy transversal to $`v`$ (as opposed to pseudo-Legendrian isotopy, which involves also a homotopy of $`v`$). More on the difference between transversal isotopy and pesudo-Legendrian isotopy will be said below.
###### Remark 5.8
Combining the previous two remarks one gets yet another proof that the relative winding number is well-defined in up to pseudo-Legendrian isotopy in a homology sphere.
#### Finite-order invariants
We formally state and motivate in this paragraph the conjecture announced in the introduction. Let $`\xi `$ be an oriented contact structure on $`M`$ (which we assume to be closed by simplicity), and let $`v`$ be a field positively transversal to $`\xi `$. Consider the spaces $`\mathrm{Leg}(M,\xi )`$, $`\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$ and $`\mathrm{Fram}(M)`$ of $`\xi `$-Legendrian, $`v`$-transverse, and framed knots in $`M`$, with the appropriate equivalence relations (namely $`\xi `$-Legendrian, pseudo-Legendrian, and framed isotopy). Enlarge these spaces by allowing immersions of $`S^1`$ rather than embeddings, and take path-connected components l, p and f, with $`\text{l}\text{p}\text{f}`$. (Concerning $`\mathrm{PLeg}`$, note that a path is a family $`(K_t,v_t)_{t[0,1]}`$ with $`v_0=v_1=v`$.) Given an Abelian group $`A`$ one can define, using the customary Vassiliev-Goussarov skein relations, the spaces $`V_\text{l}^n(A)`$, $`V_\text{p}^n(A)`$ and $`V_\text{f}^n(A)`$ of $`A`$-valued order-$`n`$ invariants under Legendrian, pseudo-Legendrian and framed isotopy respectively. Since a Legendrian isotopy is pseudo-Legendrian, and a pseudo-Legendrian isotopy is framed, using restrictions we get a commutative diagram of homomorphisms:
$$\begin{array}{ccc}V_\text{f}^n(A)& \stackrel{\varphi _{\text{f},\text{p}}^n}{}& V_\text{p}^n(A)\\ & & \\ & \varphi _{\text{f},\text{l}}^n& \varphi _{\text{p},\text{l}}^n\\ & & \\ & & V_\text{l}^n(A).\end{array}$$
Tchernov’s arguments imply that all three $`\varphi `$’s are always injective, and our conjecture is that $`\varphi _{\text{p},\text{l}}^n`$ is always an isomorphism. By again, the conjecture is equivalent to showing that every finite-order Legendrian invariant is automatically invariant also under pseudo-Legendrian isotopy. The generalized Fuchs-Tabachnikov theorem (see and ) states that $`\varphi _{\text{f},\text{l}}^n`$ is an isomorphism in many cases (e.g. if $`M`$ is a homology sphere), so $`\varphi _{\text{p},\text{l}}^n`$ is also an isomorphism in these cases. Tchernov has provided the only known examples in which $`\varphi _{\text{f},\text{l}}^n`$ is not an isomorphisms, namely he has exhibited elements of $`V_\text{l}^n(A)`$ which do not lift to $`V_\text{f}^n(A)`$. Our impression is that these elements do lift to $`V_\text{p}^n(A)`$, which would imply that $`\varphi _{\text{p},\text{l}}^n`$ is indeed an isomorphism in all known cases. Truthness of our conjecture would imply that Legendrian finite-order invariants are only sensitive to the homotopy class of a contact structure, and in particular that they cannot capture tightness.
### 5.2 Pseudo-Legendrian vs. Legendrian knots
After the work of Eliashberg , we know that on a closed manifold an overtwisted contact structure is determined up to isotopy by its homotopy class as a plane field. We discuss in this section the extent to which this fact extends in presence of a pseudo-Legendrian link. We start with an open question which arises from the results of Section 4 and will lead us to the connection with overtwisted contact structures.
#### Fixed vs. variable spine for pseudo-Legendrian links
We will adopt in this paragraph the viewpoint which allows to dismiss automorphisms of manifolds, fixing $`M`$ and considering spines and moves embedded in $`M`$, as explained after the statement of Proposition 2.3.
Theorems 2.1 and 4.1 and Proposition 4.2 leave the following question open: given an embedded spine $`PM`$ representing a concave combing on $`M`$, what intrinsic topological object is represented by $`\mathrm{C}^1`$-diagrams on $`P`$ up to pseudo-Legendrian Turaev moves? Let us introduce some notation to formalize the situation. We denote by $`𝒟^{\mathrm{PLeg}}(P)`$ the set of equivalence classes of $`\mathrm{C}^1`$-diagrams on $`P`$ up to pseudo-Legendrian Turaev moves. We also fix a representative $`v`$ of the combing carried by $`P`$ (so, $`v`$ is positively transversal to $`P`$ and restricts to $`B_{\mathrm{triv}}^3`$ on the complement of $`P`$). We consider now the set of links in $`M`$ transversal to $`v`$, and we denote by $`\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$ the quotient space under the relation of existence of a pseudo-Legendrian isotopy, i.e. a path $`(L_t,v_t)_{t[0,1]}`$ as usual, with $`v_0=v_1=v`$. We also denote by $`\mathrm{PLeg}^{\mathrm{strong}}(M,v)`$ the (bigger) quotient obtained by forcing $`(v_t)`$ to be constant. So $`\mathrm{PLeg}^{\mathrm{strong}}(M,v)`$ is just the set of equivalence classes of links transversal to $`v`$. Using Proposition 4.2 one sees that the operation of turning a diagram into a link defines a bijection
$$\psi ^{\mathrm{strong}}:𝒟^{\mathrm{PLeg}}(P)\mathrm{PLeg}^{\mathrm{strong}}(M,v).$$
(This is not quite the content of Proposition 4.2, because here $`(M,v)`$ is $`(\widehat{M}(P),\widehat{v}(P))`$ rather that $`(M(P),v(P))`$, but a link isotopy can be modified to avoid a $`B_{\mathrm{triv}}^3`$, and the conclusion follows.)
Bijectivity of $`\psi ^{\mathrm{strong}}`$ is significant if one imagines to have started with the pair $`(M,v)`$, and to have constructed $`P`$ from a normal section of $`v`$, as in Proposition 2.4. It is however less significant if one assumes only $`P`$ to be given from the beginning, because in this case $`v`$ is actually well-defined only up to homotopy, and fixing a representative looks artificial. The natural map to consider is in this case
$$\psi ^{\mathrm{weak}}:𝒟^{\mathrm{PLeg}}(P)\mathrm{PLeg}^{\mathrm{weak}}(M,v),$$
obtained by composition with the projection $`\mathrm{PLeg}^{\mathrm{strong}}(M,v)\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$. This map is of course surjective, and one can ask whether it is injective or not. Some remarks are in order:
1. Theorem 4.1 implies that if $`\psi ^{\mathrm{weak}}(D)=\psi ^{\mathrm{weak}}(D^{})`$ then there exists a circular sequence $`P=P_0P_1\mathrm{}P_n=P`$ of sliding-MP-moves and diagrams $`D_i,D_i^{}𝒟(P_i)`$ with $`D_0=D`$, $`D_n^{}=D^{}`$, $`D_iD_i^{}`$ a pseudo-Legendrian Turaev move, and $`D_{i+1}`$ the companion of $`D_i^{}`$ through $`P_iP_{i+1}`$. Checking the injectivity of $`\psi ^{\mathrm{weak}}`$ corresponds to the (purely combinatorial) question whether such a sequence $`(P_i,D_i)`$ can be replaced by one with constant $`P_i`$.
2. Using Theorem 2.1 and the fact that $`\mathrm{C}^1`$ diagrams can be followed through sliding-MP-moves, one sees quite easily that injectivity of $`\psi ^{\mathrm{weak}}`$ actually depends only on the combing carried by $`P`$, not on $`P`$ itself.
3. Injectivity of $`\psi ^{\mathrm{weak}}`$ is equivalent to injectivity of the projection $`\mathrm{PLeg}^{\mathrm{strong}}(M,v)\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$, a purely topological question. Injectivity of this projection may appear very unlikely at first sight, since it basically corresponds to the fact that a homotopy can be replaced by an isotopy. However one can remark that injectivity of projection depends only on the homotopy class of $`v`$, rather than $`v`$ itself, so one can assume that $`v`$ is transversal to an overtwisted contact structure. For overtwisted structures, after the work of Eliashberg , it is indeed true that homotopy implies isotopy, but the presence of the link of course somewhat modifies the situation. We will expound this theme in the next paragraph.
#### Overtwisted structures and overtwisted knot complements
We fix in this paragraph an overtwisted contact structure $`\xi `$ on $`M`$ (which we assume to be closed by simplicity) and a field $`v`$ positively transversal to $`\xi `$. We will denote by $`\mathrm{Leg}(M,\xi )`$ the space of Legendrian links in $`(M,\xi )`$ up to Legendrian isotopy. In , having also in mind the facts mentioned in the previous paragraph, we put forward the question of whether the natural map
$$\mathrm{Leg}(M,\xi )\mathrm{PLeg}^{\mathrm{weak}}(M,v)$$
is a bijection. A fact implying that this map is not injective in some cases was recently communicated to us by E. Giroux . He was able to construct triples $`(M,\xi ,K)`$ where $`\xi `$ is overtwisted, $`K`$ is $`\xi `$-Legendrian, and $`\xi |_{MK}`$ is tight. Let us apply a Lutz twist away from $`K`$ to get a new structure $`\xi ^{}`$ such that $`\xi ^{}`$ is homotopic to $`\xi `$ as a plane field on $`M`$, and $`\xi ^{}|_{MK}`$ is overtwisted. Using Eliashberg’s classification we consider $`\phi \mathrm{Diff}_0(M)`$ such that $`\xi ^{}=\phi ^{}(\xi )`$, and define $`K^{}=\phi (K)`$. By construction $`K`$ and $`K^{}`$ have the same image in $`\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$, but of course they are inequivalent in $`\mathrm{Leg}(M,\xi )`$.
To avoid the phenomenon discovered by Giroux we consider in $`\mathrm{Leg}(M,\xi )`$ the subset $`\mathrm{Leg}^{\mathrm{OT}}(M,\xi )`$ given by links whose complement is overtwisted. We start by showing:
###### Proposition 5.9
The natural map $`\mathrm{Leg}^{\mathrm{OT}}(M,\xi )\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$ is surjective.
Proof of5.9. Let $`L`$ be transversal to $`v`$, and fix a metric on $`M`$. Let $`\eta `$ be a positive contact structure near $`L`$ with $`\eta =v^{}`$ on $`K`$ (such an $`\eta `$ is unique up to isomorphism). Extend $`\eta `$ to any plane field homotopic to $`v^{}`$ (and hence to $`\xi `$) on $`M`$. So $`\eta `$ is a plane distribution which has a contact zone, and $`L`$ lies in this contact zone. The technique of Eliashberg now allows to homotope $`\eta `$ away from its contact zone to an overtwisted contact structure $`\xi ^{}`$. The resulting $`\xi ^{}`$ is now isotopic to $`\xi `$, again by Eliashberg’s result. If $`\phi \mathrm{Diff}_0(M)`$ and $`\xi ^{}=\phi ^{}(\xi )`$ we define $`L^{}=\phi (L)`$. By construction $`(L^{},v)`$ is pseudo-Legendrian isotopic to $`(L,v)`$, and surjectivity is proved. 5.9
We cannot presently state whether the map $`\mathrm{Leg}^{\mathrm{OT}}(M,\xi )\mathrm{PLeg}^{\mathrm{weak}}(M,v)`$ is injective or not in general. We only give a partial argument (based on the techniques of Eliashberg again), and mention where the difficulty arises. Assume that $`L_0`$ and $`L_1`$ are $`\xi `$-Legendrian with overtwisted complements and define equivalent pseudo-Legendrian links. Then there exists a continuous family $`(L_t,\xi _t)_{t[0,1]}`$, where $`\xi _0=\xi _1=\xi `$ but $`\xi _t`$ is only a plane field for $`t0,1`$. Eliashberg’s contactization methods for homotopies, together with the uniqueness of contact structures in the neighbourhood of Legendrian links, should in our opinion allow to replace such a $`(\xi _t)_{t[0,1]}`$ by another one in which each $`\xi _t`$ is a contact structure (and still contains $`L_t`$ as a Legendrian link). Applying Gray’s theorem we get an isotopy $`(\phi _t)_{t[0,1]}`$ such that $`\xi _t=\phi _t^{}(\xi _0)`$. Setting $`\stackrel{~}{L}_t=\phi _t(L_t)`$ we get a Legendrian isotopy between $`L_0`$ and $`\phi _1(L_1)`$. The question whether $`\phi _1(L_1)`$ is automatically Legendrian isotopic to $`L_1`$, at least for some classes of manifolds, now depends on the analysis of the group $`\mathrm{Aut}(M,\xi )\mathrm{Diff}_0(M)`$, which we leave unsettled for the time being.
benedett@dm.unipi.it
petronio@dm.unipi.it
Dipartimento di Matematica
Via F. Buonarroti, 2
I-56127, PISA (Italy) |
warning/0001/math0001081.html | ar5iv | text | # Trace Functionals of the Kontsevich Quantization
## 1 Introduction
The gist of deformation quantization of a Poisson manifold lies in modifying the usual (commutative) product of functions to obtain a new operation (star product) in such a way that the first approximation is given by the Poisson bracket. This approach was first proposed in \[BFFLS\], and is at the moment one of the most developed methods to relate ”classical” data represented by a Poisson manifold to a ”quantum” associative algebra. With the advent of the Kontsevich star product in 1997 (see \[Ko\]), when the problem of existence of these algebras for arbitrary Poisson structures has been solved in full generality, it became abundantly clear that the utility of deformation quantization hinges on whether the resulting algebras are amenable to the standard operator theory tools, notwithstanding the fact that no actual operators are present. If this universal version of quantization is to be useful, some ”spectral” properties must transpire. Sure enough, on the infinitesimal level there is the notion of Poisson trace, whose existence is intimately linked with the modular class (or, in mundane terms, with the existence of a unimodular measure - one invariant with respect to all Hamiltonian flows) as demonstrated by Weinstein in \[W2\]. Another bit of information supporting the idea of the Kontsevich quantization algebras having spectral properties is the fact that on $`^d`$ any constant coefficient Poisson structure induces a Kontsevich quantized algebra with a trace given by integration.
On the other hand, considering all possible star products on symplectic manifolds, Connes et al. \[CFS\] were able to identify obstructions to the existence of a trace on a fixed star-product algebra. These turned out to be certain cocycles in the cyclic cohomology of the Hochshild complex. Moreover, the authors constructed a simple example where integration does not yield a trace functional. Combined with the fact that the Liouville measure on a symplectic manifold is unimodular par excellence, the results of Connes et al. permit one to conclude that unimodularity alone is not sufficient even in the nicest possible (symplectic) setting. Although later Fedosov \[F\] showed that on an arbitrary symplectic manifold there is a star-product quantization equipped with a trace given by integration, the following question remains: Is there a trace functional on the Kontsevich quantization algebra of an arbitrary Poisson manifold? The present paper aims to single out a class of Poisson manifolds that possess a trace functional on the Kontsevich quantization algebra. In the process we first briefly recall the essentials in Section 2, then study the Kontsevich quantization of the dual of nilpotent Lie algebras endowed with the standard Lie-Poisson structure (Section 3), go on to describe the properties of Morita equivalent Poisson manifolds in Section 4, and finally get everything together in Section 5.
Aknowledgments. We would like to thank Joseph Donin for many valuable suggestions, and Alan Weinstein for useful comments and a careful reading of the manuscript.
## 2 Star Products
Given a smooth ($`𝒞^{\mathrm{}}`$) Poisson manifold $`(P,\pi )`$, the set of smooth functions $`𝒞^{\mathrm{}}(P)`$ can be viewed as a commutative algebra. The star product on $`𝒞^{\mathrm{}}(P)`$ (c.f. \[BFFLS\]) is an associative $`[[\mathrm{}]]`$-linear product on $`𝒞^{\mathrm{}}(P)[[\mathrm{}]]=(P)`$ expressed by the following formula for $`f,g𝒞^{\mathrm{}}(P)(P)`$:
$$(f,g)fg=fg+\mathrm{}\pi (f,g)+\mathrm{}^2B^2(f,g)+\mathrm{}(P),$$
where $`\mathrm{}`$ is a formal variable, and $`B^i`$ are bidifferential operators (bilinear maps $`(P)\times (P)(P)`$ which are differential operators of globally bounded order with respect to each argument). The product of arbitrary elements of $`(P)`$ is defined by the condition of linearity over $``$ and $`\mathrm{}`$-adic continuity:
$$\left(\underset{n0}{}\mathrm{}^nf_n\right)\left(\underset{n0}{}\mathrm{}^ng_n\right):=\underset{k,l0}{}\mathrm{}^{k+l}f_kg_l+\underset{k,l0,m1}{}\mathrm{}^{k+l+m}B^m(f_k,g_l).$$
Now we proceed to give a brief account of the universal deformation quantization of $`^d`$ or a domain thereof due to Kontsevich \[Ko\]. From this point on, the symbol $``$ refers to the canonical Kontsevich star product. In order to describe the terms proportional to $`\mathrm{}^n`$ for any integer $`n1`$, Kontsevich introduced a special class $`G_n`$ of oriented labeled graphs called admissible graphs.
###### Definition 1
An oriented graph $`\mathrm{\Gamma }`$ is admissible ($`\mathrm{\Gamma }G_n`$, $`(n1)`$) if:
1. $`\mathrm{\Gamma }`$ has $`n+2`$ vertices labeled $`\{1,2,\mathrm{}n,L,R\}`$ where L and R stand for Left and Right respectively, and $`\mathrm{\Gamma }`$ has $`2n`$ oriented edges labeled $`\{i_1,j_1,i_2,j_2,\mathrm{},i_n,j_n\}`$;
2. The pair of edges $`\{i_m,j_m\},1mn`$ starts at the vertex m;
3. $`\mathrm{\Gamma }`$ has no loops (edges starting at some vertex and ending at that same vertex) and no parallel multiple edges (edges sharing the same starting and ending vertices).
The class $`G_n`$ is finite. For $`n1`$ the first edge $`i_k`$ starting at the vertex $`k`$ has $`n+1`$ possible ending vertices since there are no loops, while the second edge $`j_k`$ has only $`n`$ possible ”landing sites” since there is no parallel multiple edges. Thus there are $`n(n+1)`$ ways to draw the pair of edges starting at some vertex and therefore $`G_n`$ has $`n(n+1)^n`$ elements. For $`n=0`$, $`G_0`$ has only one element: the graph having $`\{L,R\}`$ as set of vertices and no edges. Of all admissible graphs, there is an important particular subclass, which we call the class of unions of subgraphs. This notion will be used in the sequel. Here is how we define the elements of this subclass:
###### Definition 2
A graph $`\mathrm{\Gamma }G_r`$ is the union of two subgraphs $`\mathrm{\Gamma }_1`$ and $`\mathrm{\Gamma }_2`$ with $`r_1+r_2=r`$, if the subset $`(1,\mathrm{},r)`$ of the set of vertices of $`\mathrm{\Gamma }`$ can be split into two parts $`(a_1,\mathrm{},a_{r_1})`$ and $`(b_1,\mathrm{},b_{r_2})`$ such that there is no edge between these two subsets of vertices.
A bidifferential operator $`(f,g)B_\mathrm{\Gamma }(f,g),f,g𝒞^{\mathrm{}}(^d)`$ is associated to each graph $`\mathrm{\Gamma }G_n,n1`$. To each vertex $`k`$, $`1kn`$ one associates the components $`\pi _{i_kj_k}`$ of the Poisson structure, $`f`$ is associated to the vertex $`L`$ and $`g`$ to the vertex $`R`$. Each edge such as $`i_k`$ acts by partial differentiation with respect to $`x_{i_k}`$ on its ending vertex. For $`n=0`$ we simply have the usual product of $`f`$ and $`g`$. In what follows we use the words ”graph” and ”bidifferential operator encoded by the graph” interchangeably.
Now we go on to describe the weights $`w_\mathrm{\Gamma }`$. Let $`=\{z|\text{Im}(z)>0\}`$ be the upper half-plane. $`_n`$ will denote the configuration space $`\{z_1,\mathrm{},z_n|z_iz_j\text{for}ij\}`$. $`_n`$ is an open submanifold of $`^n`$. Let $`\varphi :_2/2\pi `$ be the function:
$$\varphi (z_1,z_2)=\frac{1}{2\sqrt{1}}\text{Log}\frac{(z_2z_1)(\overline{z}_2z_1)}{(z_2\overline{z}_1)(\overline{z}_2\overline{z}_1)}.$$
$`\varphi (z_1,z_2)`$ is extended by continuity for $`z_1,z_2,z_1z_2`$.
For a graph $`\mathrm{\Gamma }G_n`$, the vertex $`k`$, $`1kn`$ is associated with the variable $`z_k`$, the vertex $`L`$ with $`0`$, and the vertex $`R`$ with $`1`$.
The weight $`w_\mathrm{\Gamma }`$ is defined by integrating an $`2n`$-form over $`_n`$:
$$w_\mathrm{\Gamma }=\frac{1}{n!(2\pi )^{2n}}__n\underset{1kn}{}(d\varphi (z_k,I_k)d\varphi (z_k,J_k)),$$
where $`I_k`$ ($`J_k`$) denotes the variable or real number associated with the ending vertex of the edge $`i_k`$ ($`j_k`$). It is showed in \[Ko\] that this integral is absolutely convergent. As one can see, the weights do not depend on the Poisson structure or the dimension of the underlying manifold. Combining these constructions Kontsevich proved
###### Background Theorem 1 (Kontsevich)
Let $`\pi `$ be a Poisson structure in a domain of $`^d`$. The formula
$$fg:=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{}^n\underset{\mathrm{\Gamma }G_n}{}w_\mathrm{\Gamma }B_{\mathrm{\Gamma },\pi }(f,g)$$
defines an associative product up to gauge equivalence.
Furthermore, all the machinery above as well as Background Theorem 1 generalize to arbitrary Poisson manifolds using localization. There are no cohomological obstructions.
Finally, we deal with trace functionals on star-product algebras. Consider the linear map $`𝖳r`$ on the compactly supported Kontsevich quantization algebra $`_c(P)`$ with values in formal Laurent power series $`\mathrm{}^{[d/2]}[[\mathrm{}]]`$:
$$𝖳r_{P,\mu }(F)=\mathrm{}^{[d/2]}_PF\mu ,$$
where \[ \] denotes the integral part of a number and $`d`$ is the dimension of $`P`$. When $`P`$ is symplectic this functional (up to a contant) coincides with the usual one as defined in \[CFS\], provided $`\mu `$ is the Liouville volume form. In general, we may have many meaningful choices for the volume form, which is emphasized by the subscript. Following Fedosov \[F\] we declare that the functional $`𝖳r`$ satisfies the trace property with respect to the Kontsevich star product if
$$𝖳r(FG)=𝖳r(GF).$$
Consequently, the problem of defining a trace functional amounts to proving that a linear functional $`𝖳r`$ satisfies the trace property. Formally, we have
###### Definition 3
The linear functional $`𝖳r`$ on $`_c(P)`$ is called a trace functional if it satisfies the trace property.
## 3 Quantization of $`𝔤^{}`$
Given a finite-dimensional Lie group $`G`$, we denote its Lie algebra by $`𝔤`$, the dual of $`𝔤`$ by $`𝔤^{}`$. The structure constants of $`𝔤`$ are determined by the set of relations with respect to a basis of $`𝔤`$:
$$[X_i,X_j]=X_kC_{ij}^k,X_i𝔤i\{1,\mathrm{}.,d\}.$$
Now we let $`(x_1,\mathrm{}.,x_d)`$ denote linear coordinates on $`𝔤^{}`$. Then the natural Lie-Poisson structure
$$\pi =\underset{i,j}{}\pi _{ij}_i_j,$$
is expressed in terms of coordinates via
$$\pi _{ij}=x_kC_{ij}^k.$$
We begin by looking at the role of unimodularity in vanishing of the Poisson trace.
###### Proposition 1
For any $`f,g𝒞_c^{\mathrm{}}(𝔤^{})`$,
$$_^d\pi (f,g)𝑑x_1\mathrm{}dx_d=\mathrm{\hspace{0.33em}0}$$
if and only if
$$\underset{i}{}C_{ij}^i=\mathrm{\hspace{0.33em}0}j\{1,\mathrm{}.,d\}.$$
Proof. We proceed by making a simple observation: $`C_{ij}^k_i_jf`$ adds up to zero by virtue of $`C_{ij}^k`$ being antisymmetric. Now straightforward integration by parts of a single component $`\pi _{ij}_if_jg`$ with subsequent summation does the trick. Q.E.D.
Remark. Proposition 1 is basically a rephrasing of the statement made in \[W2\], according to which ”on the dual of a Lie algebra $`𝔤`$, with its Lie-Poisson structure, the modular vector field with respect to any translation-invariant density is the constant vector field with value $`\mathrm{𝖳𝗋}ad`$, the trace of the adjoint representation”.
Now we are in a position to state the main theorem of this section.
###### Theorem 1
For all $`F,G_c(𝔤^{})`$ the equality
$$_^dFG𝑑x_1\mathrm{}dx_d=_^dGF𝑑x_1\mathrm{}dx_d$$
holds for $`𝔤`$ nilpotent.
Proof. First of all we note that this Theorem will follow if we manage to show the equality above for $`f,g𝒞_c^{\mathrm{}}(𝔤^{})`$. Next we invoke the formula (due to Cattaneo and Felder \[CaFe\]) relating the weight of $`B^n(f,g)`$ to that of $`B^n(g,f)`$. Namely, $`w_{B^n(f,g)}=(1)^nw_{B^n(g,f)}`$. This signifies that for all even $`n`$ the corresponding integrals do not affect the (presumed) trace property. From now on we work only with the graphs involving odd $`n`$. To manipulate the operators effectively, we introduce the following notation: $`B_k^n(f,g)`$ stands for a bidifferential operator that has order $`n`$ in the star product formula, and, in addition, has exactly $`k`$ edges acting summarily on $`f`$ AND $`g`$. Nothing else is specified, but below we take one particular operator at a time so that this apparent ambiguity does not matter. Now, by the linearity of $`\pi `$, it follows that $`nk2n`$. An important fact used later is that any graph denoted by $`B_n^n(f,g)`$ necessarily includes at least one loop (an expression of the form $`_{i_m}\pi _{i_1j_1}_{i_1}\pi _{i_2j_2}\mathrm{}_{i_{m1}}\pi _{i_mj_m}`$ after possible reshuffling). As Kathotia \[Ka\] showed, such graphs vanish for nilpotent Lie algebras due to a theorem of Varadarajan \[V\]:
###### Background Theorem 2 (Varadarajan, Theorem 3.5.4.)
Let $`𝔤`$ be a d-dimensional nilpotent Lie algebra over a field of characteristic zero. Then there is a basis $`\{X_1,\mathrm{},X_d\}`$ of $`𝔤`$ such that the structure constants $`C_{ij}^k`$ defined by $`[X_i,X_j]=X_kC_{ij}^k`$ satisfy
$$C_{ij}^r=\mathrm{\hspace{0.33em}0},rmin(i,j).$$
An easy corollary to the above theorem assures that any cyclic product of the form $`C_{i_1j_1}^{j_m}C_{i_2j_2}^{j_1}\mathrm{}C_{i_mj_m}^{j_{m1}}(m>1)`$ vanishes and this, in turn, yields the vanishing of loop graphs.
Now since we have seen that in the nilpotent setting there are no loop graphs, the idea is to rewrite the integrals of $`B_k^n`$ as combinations of those of $`B_n^n`$. This is accomplished via integration by parts. Once again using the linearity of the Poisson structure, we have a vertex such that one outgoing edge acts on $`f`$, whereas the other acts on $`g`$ (otherwise $`k=n`$ and we are done). We use one of these edges to perform integration by parts. As a result we obtain a few integrals of admissible graphs of the form $`B_{k1}^n`$. To see that we first note that all nilpotent groups are unimodular. Indeed, if there is a $`C_{ij}^i0`$ with respect to any basis, we would have an $`X_i`$ such that $`[X_i,X_j]=X_iC_{ij}^i`$, and iterating this $`l`$ times we would violate the condition $`𝔤^l=0`$. So the admissibility of the resulting graphs as well as the decreasing of the number of edges acting on $`f`$ and $`g`$ follows from Proposition 1 and its proof. Inductively we get rid of all antisymmetric ($`n`$ odd) graphs. Q.E.D.
Remark. This result is in some sense the best attainable. Consider a unimodular Lie algebra $`𝔲`$ such that the Kontsevich quantization of $`𝔲^{}`$ does contain loop graphs. Then by choosing $`f,g`$ supported in a small neighborhood of the origin in $`𝔲^{}`$, we can assure that the contribution of the graphs with $`n`$ odd, $`k=n`$, is not eliminated by some freaky cancellation. Thus the statement of Theorem 1 is not true for arbitrary unimodular Lie algebras.
## 4 Morita Equivalence
The essential sources used here are \[W1\], \[X\], and \[GG\]. Following \[W1\] we recall that a full dual pair $`P_1\stackrel{\rho _1}{}W\stackrel{\rho _2}{}P_2`$ consists of two Poisson manifolds $`(P_1,\pi _1)`$ and $`(P_2,\pi _2)`$, a symplectic manifold $`W`$, and two submersions $`\rho _1:WP_1`$ and $`\rho _2:WP_2`$ such that $`\rho _1`$ is Poisson, $`\rho _2`$ is anti-Poisson, and the fibers of $`\rho _1`$ and $`\rho _2`$ are symplectic orthogonal to each other. A Poisson (or anti-Poisson) mapping is said to be complete if the pullback of a complete Hamiltonian flow under this mapping is complete. A full dual pair is called complete if both $`\rho _1`$ and $`\rho _2`$ are complete. The Poisson manifolds $`(P_1,\pi _1)`$ and $`(P_2,\pi _2)`$ are Morita equivalent if there exists a complete full dual pair $`P_1\stackrel{\rho _1}{}W\stackrel{\rho _2}{}P_2`$ such that $`\rho _1`$ and $`\rho _2`$ both have connected and simply connected fibers. The notion of Morita equivalence of Poisson manifolds was introduced and studied by Xu \[X\], as a classical analogue of the Morita equivalence of $`C^{}`$-algebras.
###### Theorem 2
The property of being equipped with a trace functional on the Kontsevich quantization algebra is an invariant of Morita equivalence.
Proof. By the hypothesis we have $`P_1`$, $`P_2`$ Morita equivalent, and $`_c(P_1)`$ possesses a trace functional. Whence
$$𝖳r_{P_1,\mu ^{}}(F)=\mathrm{}^{[d/2]}_{P_1}F\mu ^{},$$
where $`\mu ^{}`$ is a unimodular volume form. Now we recall a result from \[GG\]. Here mod$`(P_i)`$ denotes the modular class on respective manifolds.
###### Background Theorem 3 (Ginzburg, Golubev)
Let $`P_1`$ and $`P_2`$ be Morita equivalent and let, in addition, $`P_1`$ be locally unimodular. Then $`P_2`$ is also locally unimodular and $`\mathrm{m}\mathrm{o}\mathrm{d}(\mathrm{P}_1)`$ goes to $`\mathrm{m}\mathrm{o}\mathrm{d}(\mathrm{P}_2)`$ under the natural isomorphism of the first Poisson cohomology groups $`E:H_\pi ^1(P_1)\stackrel{}{}H_\pi ^1(P_2)`$, i. e. $`E(\mathrm{m}\mathrm{o}\mathrm{d}(P_1))=\mathrm{m}\mathrm{o}\mathrm{d}(P_2)`$.
Using the above theorem we conclude $`P_2`$ is unimodular too. Furthermore, starting off with $`\mu ^{}`$, utilizing the action of Brylinski’s symplectic star operator (see \[Br\] and also \[GG\] for more details), we arrive at a global unimodular volume form $`\mu ^{\prime \prime }`$ on $`P_2`$. At this point we invoke a theorem of Weinstein \[W1\] concerning transversal Poisson structures in full dual pairs:
###### Background Theorem 4 (Weinstein)
Let $`P_1\stackrel{\rho _1}{}W\stackrel{\rho _2}{}P_2`$ be a full dual pair. For each $`xW`$, the transverse Poisson structures on $`P_1`$ and $`P_2`$ at $`\rho _1(x)`$ and $`\rho _2(x)`$ are anti-isomorphic. Consequently, if $`\mathrm{d}\mathrm{i}\mathrm{m}(P_1)\mathrm{d}\mathrm{i}\mathrm{m}(P_2`$), $`P_2`$ is locally anti-isomorphic to the product of $`P_1`$ with a symplectic manifold.
Now we introduce a fixed open cover $`\{U_\alpha \}_{\alpha J}`$ of $`W`$ such that $`U_\alpha ^{}=\rho _1(U_\alpha )`$, $`U_\alpha ^{\prime \prime }=\rho _2(U_\alpha )`$ and the latter ones are open covers of $`P_1`$ and $`P_2`$ respectively. They are so fine that there is a local anti-isomorphism of Background Theorem 4 in each $`U_\alpha `$ (by refining the cover of $`W`$ we can always achieve that). From Background Theorem 4 we infer the relation between the local expressions for $`\mu ^{}`$ and $`\mu ^{\prime \prime }`$, which we denote by $`\mu _\alpha ^{}`$ and $`\mu _\alpha ^{\prime \prime }`$. Namely, either $`\mu _\alpha ^{}=\mu _\alpha ^{\prime \prime }\kappa `$ (this happens to be the case if $`\text{dim}(P_1)>\text{dim}(P_2)`$), or $`\mu _\alpha ^{\prime \prime }=\mu _\alpha ^{}\kappa `$ ($`\text{dim}(P_1)<\text{dim}(P_2)`$), or $`\mu ^{}=a\mu ^{\prime \prime }`$ ($`\text{dim}(P_1)=\text{dim}(P_2)`$), where $`\kappa `$ is a Liouville volume form and $`a`$ is a nonzero constant. By the standard manifold theory we may view $`U_\alpha `$’s as domains of $`^d`$. A crucial fact needed at this juncture is the following: the graphs in $`U_\alpha ^{}`$, $`U_\alpha ^{\prime \prime }`$ are unions of subgraphs (c. f. Definition 2) of symplectic and transversal components. Moreover, the symplectic components necessarily have equal number of edges acting on $`f`$ and $`g`$, so that up to a constant, the graphs can be thought of as those involving transversal components only. Now assembling the above volume forms and the graphs involving transversal components in the integral, we see that the vanishing of one on $`U_\alpha ^{}`$ entails the vanishing of its counterpart on $`U_\alpha ^{\prime \prime }`$. Finally, using a partition of unity subordinate to the cover, we globalize the trace functional. Q.E.D.
## 5 The Main Theorem via Symplectic Realizations
We recall the appropriate definitions. A symplectic realization of a Poisson manifold $`P`$ is a pair $`(W,\rho )`$, where $`W`$ is a symplectic manifold and $`\rho `$ is a Poisson morphism from $`W`$ to $`P`$. A symplectic realization $`\rho :WP`$ is called complete if $`\rho `$ is complete as a Poisson map and $`\rho `$ is said to be full if it is a submersion. There are many ways to construct symplectic realizations, but we are interested in just one particular way of doing it. Precisely, we want to realize Poisson manifolds as quotients of symplectic manifolds by group actions. In applications, as Weinstein \[W1\] pointed out, the symplectic manifold $`W`$ may represent a collection of states such that the points belonging to the same orbit of $`G`$-action are considered to be physically indistinguishable. Thus the set of ”true physical states” is none other than the Poisson manifold $`W/G`$, and the group $`G`$ is called a gauge group. This motivates the definition below.
###### Definition 4
A Poisson manifold $`P`$ is symplectically realizable with a gauge group $`G`$ if there exists a symplectic manifold $`W`$ and a Lie group $`G`$ freely acting on $`W`$ by symplectomorphisms, such that $`W/GP`$.
An upshot of this definition is that a) if $`P`$ is symplectically realizable with a gauge group $`G`$, $`\rho :WP`$ is a quotient map and necessarily a complete full surjective symplectic realization; b) $`P`$ is Morita equivalent to an open subset of $`𝔤^{}`$.
###### Theorem 3
The Kontsevich quantization algebra of a symplectically realizable Poisson manifold with a nilpotent gauge group possesses a trace functional.
Proof. Applying Theorem 1 we obtain a trace functional on $`𝔤^{}`$ in view of the hypothesis, and Theorem 2 now ensures that $`𝖳r_{P,\mu }`$ satisfies the trace property, where $`\mu `$ is the volume form obtained from the standard translation-invariant form on $`𝔤^{}`$. Q.E.D. |
warning/0001/cond-mat0001132.html | ar5iv | text | # Theory of magnetic order in the three-dimensional spatially anisotropic Heisenberg model
## I Introduction
The magnetic properties of spatially anisotropic antiferromagnetic (AFM) quantum spin systems, such as the quasi-two-dimensional (2D) parent compounds of high-$`T_c`$ superconductors (e. g. La<sub>2</sub>CuO<sub>4</sub>, Ca(Sr)CuO<sub>2</sub>) and the quasi-1D cuprates Sr<sub>2</sub>CuO<sub>3</sub>, Ca<sub>2</sub>CuO<sub>3</sub>, and SrCuO<sub>2</sub>, are of current interest. The main problem is the influence of spatial anisotropy on the staggered magnetization $`m`$ and the Neél temperature $`T_N`$ in the 3D spin-$`\frac{1}{2}`$ AFM Heisenberg model
$$H=J_x\left[\underset{ij_x}{}𝐒_i𝐒_j+R_y\underset{ij_y}{}𝐒_i𝐒_j+R_z\underset{ij_z}{}𝐒_i𝐒_j\right].$$
(1)
Here $`R_y=J_y/J_x`$, $`R_z=J_z/J_x`$ (throughout we set $`J_x=1`$) and $`ij_{x,y,z}`$ denote nearest-neighbor (NN) bonds along the $`x`$-, $`y`$\- or $`z`$-directions of a simple cubic lattice. For real systems, we consider $`0R_zR_y1`$.
In the paramagnetic phase, there exists a pronounced AFM short-range order (SRO) which is reflected by a maximum in the temperature dependence of the magnetic susceptibility at $`T_{max}`$, where $`0.64<T_{max}<1.2`$. However, the RPA spin-wave theories and the mean-field theories using auxiliary-field representations (Schwinger-boson, Holstein-Primakoff, Dyson-Maleev, and boson-fermion representations ) which were developed for the quasi-2D model with $`R_y=1`$, are valid only at sufficiently low temperatures. In those theories, the temperature dependent SRO is not adequately taken into account; in particular, the maximum in the magnetic susceptibility cannot be reproduced. In the chain mean-field approaches, recently improved by spin-fluctuation corrections which lower $`m`$ and $`T_N`$, an asymmetry between intrachain and interchain correlations is introduced. As was shown in Ref. on the basis of a detailed estimate of the exchange integrals for the quasi-1D cuprates using a first-principle calculation (Sr<sub>2</sub>CuO<sub>3</sub>: $`R_y0.004`$; Ca<sub>2</sub>CuO<sub>3</sub>: $`R_y0.02`$), all previous approaches overestimate both $`m`$ and $`T_N`$. This deficiency is calling for a theory that provides an improved description of SRO over the whole temperature region. In Ref. , a spin-rotation-invariant Green’s-function theory for the 2D isotropic Heisenberg and $`t`$-$`J`$ models was developed, which yields a good description of spin correlation functions of arbitrary range and at arbitrary temperatures. Moreover, the susceptibility maximum was obtained in good agreement with quantum Monte Carlo calculations. Applying this approach to the 2D anisotropic Heisenberg model, the short-ranged spin correlations at $`T=0`$ are well reproduced as compared with exact diagonalization (ED) data. Accordingly, we expect such a theory to describe the SRO properties quite well also in the 3D model (1).
In this paper we extend the Green’s-function approach of Refs. and and present a theory of AFM long-range order (LRO) and SRO for the 3D anisotropic Heisenberg model (1) (Sec. II). Thereby, the correlations along all spatial directions are described on the same footing. In Sec. III the ground state is investigated, where the magnetization and short-ranged spin correlation functions are calculated. In Sec. IV we present our finite-temperature results on the $`R_z`$-dependence of $`T_N`$, $`m(T)`$, and of the AFM correlation lengths. Moreover, for the first time, the effects of an arbitrary spatial anisotropy on the temperature dependence of the uniform static spin susceptibility, especially on $`T_{max}`$, are investigated. The results are compared with experiments on La<sub>2</sub>CuO<sub>4</sub> (magnetization, correlation length, magnetic susceptibility). The summary of our work can be found in Sec. V.
## II Dynamic spin susceptibility
To determine the dynamic spin susceptibility $`\chi ^+(𝐪,\omega )=S_𝐪^+;S_𝐪^{}_\omega `$ by the projection method outlined in Ref. , we choose the two-operator basis $`𝐀=(S_𝐪^+,i\dot{S}_𝐪^+)^T`$ and consider the two-time retarded matrix Green’s function in a generalized mean-field approximation, $`𝐀;𝐀^+_\omega =\left[\omega 𝐌^{}𝐌^1\right]^1𝐌`$ with $`𝐌=[𝐀,𝐀^+]`$ and $`𝐌^{}=[i\dot{𝐀},𝐀^+]`$, using Zubarev’s notation. We get
$$\chi ^+(𝐪,\omega )=\frac{M_𝐪^{(1)}}{\omega ^2\omega _𝐪^2}.$$
(2)
The spectral moment $`M_𝐪^{(1)}=[i\dot{S}_𝐪^+,S_𝐪^{}]`$ is given by
$$M_𝐪^{(1)}=4C_{1,0,0}(1\mathrm{cos}q_x)4R_yC_{0,1,0}(1\mathrm{cos}q_y)4R_zC_{0,0,1}(1\mathrm{cos}q_z).$$
(3)
The two-spin correlation functions $`C_𝐫=S_0^+S_𝐫^{}C_{n,m,l}`$ with $`𝐫=n𝐞_x+m𝐞_y+l𝐞_z`$ are calculated from
$$C_𝐫=\frac{1}{N}\underset{𝐪}{}C_𝐪\text{e}^{i\mathrm{𝐪𝐫}},C_𝐪=\frac{M_𝐪^{(1)}}{2\omega _𝐪}[1+2n(\omega _𝐪)],$$
(4)
where $`n(\omega _𝐪)=\left(\text{e}^{\omega _𝐪/T}1\right)^1`$. The NN correlation functions are directly related to the internal energy per site by $`ϵ=\frac{3}{2}\left(C_{1,0,0}+R_yC_{0,1,0}+R_zC_{0,0,1}\right)`$.
To obtain the spectrum in the approximation $`\ddot{S}_𝐪^+=\omega _𝐪^2S_𝐪^+`$, we take the site representation and decouple the products of three spin operators in $`\ddot{S}_i^+`$ along NN sequences introducing vertex parameters in the spirit of the scheme proposed by Shimahara and Takada and extending the decoupling given in Ref. ,
$$S_i^+S_j^+S_l^{}=\alpha _1^{x,y,z}S_j^+S_l^{}S_i^++\alpha _2S_i^+S_l^{}S_j^+.$$
(5)
Here $`\alpha _1^x`$, $`\alpha _1^y`$, and $`\alpha _1^z`$ are attached to NN correlation functions along the $`x`$-, $`y`$\- and $`z`$-directions, respectively, and $`\alpha _2`$ is associated with the longer ranged correlation functions. We obtain
$`\omega _𝐪^2`$ $`=`$ $`1+R_y^2+R_z^2\mathrm{cos}q_xR_y^2\mathrm{cos}q_yR_z^2\mathrm{cos}q_z+2\alpha _1^xC_{1,0,0}\mathrm{cos}(2q_x)2\alpha _2C_{2,0,0}\mathrm{cos}q_x`$ (13)
$`+2R_y^2(\alpha _1^yC_{0,1,0}\mathrm{cos}(2q_y)\alpha _2C_{0,2,0}\mathrm{cos}q_y)+2R_z^2(\alpha _1^zC_{0,0,1}\mathrm{cos}(2q_z)\alpha _2C_{0,0,2}\mathrm{cos}q_z)`$
$`2\alpha _1^xC_{1,0,0}\mathrm{cos}q_x+2\alpha _2C_{2,0,0}2R_y^2(\alpha _1^yC_{0,1,0}\mathrm{cos}q_y\alpha _2C_{0,2,0})2R_z^2(\alpha _1^zC_{0,0,1}\mathrm{cos}q_z\alpha _2C_{0,0,2})`$
$`+4R_y((\alpha _1^xC_{1,0,0}+\alpha _1^yC_{0,1,0})\mathrm{cos}q_x\mathrm{cos}q_y\alpha _2C_{1,1,0}(\mathrm{cos}q_x+\mathrm{cos}q_y))`$
$`+4R_z((\alpha _1^xC_{1,0,0}+\alpha _1^zC_{0,0,1})\mathrm{cos}q_x\mathrm{cos}q_z\alpha _2C_{1,0,1}(\mathrm{cos}q_x+\mathrm{cos}q_z))`$
$`+4R_yR_z((\alpha _1^yC_{0,1,0}+\alpha _1^zC_{0,0,1})\mathrm{cos}q_y\mathrm{cos}q_z\alpha _2C_{0,1,1}(\mathrm{cos}q_y+\mathrm{cos}q_z))`$
$`4R_y(\alpha _1^xC_{1,0,0}\mathrm{cos}q_y+\alpha _1^yC_{0,1,0}\mathrm{cos}q_x2\alpha _2C_{1,1,0})4R_z(\alpha _1^xC_{1,0,0}\mathrm{cos}q_z+\alpha _1^zC_{0,0,1}\mathrm{cos}q_x2\alpha _2C_{1,0,1})`$
$`4R_yR_z(\alpha _1^yC_{0,1,0}\mathrm{cos}q_z+\alpha _1^zC_{0,0,1}\mathrm{cos}q_y2\alpha _2C_{0,1,1}).`$
Note that our scheme preserves the rotational symmetry in spin space, i.e. $`\chi ^{zz}(𝐪,\omega )\chi (𝐪,\omega )=\frac{1}{2}\chi ^+(𝐪,\omega )`$. For $`|𝐪|1`$ we have
$$\omega _𝐪^2=c_x^2q_x^2+c_y^2q_y^2+c_z^2q_z^2,$$
(14)
with the squared spin-wave velocities
$`c_x^2`$ $`=`$ $`{\displaystyle \frac{1}{2}}3\alpha _1^xC_{1,0,0}+\alpha _2C_{2,0,0}2R_y(\alpha _1^xC_{1,0,0}\alpha _2C_{1,0,0})2R_z(\alpha _1^xC_{1,0,0}\alpha _2C_{1,0,1}),`$ (15)
$`c_y^2`$ $`=`$ $`R_y^2({\displaystyle \frac{1}{2}}3\alpha _1^yC_{0,1,0}+\alpha _2C_{0,2,0})2R_y(\alpha _1^yC_{0,1,0}\alpha _2C_{1,1,0})2R_yR_z(\alpha _1^yC_{0,1,0}\alpha _2C_{0,1,1}),`$ (16)
and (17)
$`c_z^2`$ $`=`$ $`R_z^2({\displaystyle \frac{1}{2}}3\alpha _1^zC_{0,0,1}+\alpha _2C_{0,0,2})2R_z(\alpha _1^zC_{0,0,1}\alpha _2C_{1,0,1})2R_yR_z(\alpha _1^zC_{0,0,1}\alpha _2C_{0,1,1}).`$ (18)
Considering the uniform static spin susceptibility $`\chi =lim_{𝐪0}M_𝐪^{(1)}/(2\omega _𝐪^2)`$, the ratio of the anisotropic functions $`M_𝐪^{(1)}`$ and $`\omega _𝐪^2`$ must be isotropic in the limit $`𝐪0`$. That is, the conditions
$$(c_y/c_x)^2=R_yC_{0,1,0}/C_{1,0,0}$$
(19)
and
$$(c_z/c_x)^2=R_zC_{0,0,1}/C_{1,0,0}$$
(20)
have to be fulfilled.
The critical behavior of the model (1) is reflected in our theory by the closure of the spectrum gap at $`𝐐=(\pi ,\pi ,\pi )`$ as $`T`$ approaches $`T_N`$ from above, so that $`lim_{TT_N}\chi ^1(𝐐)=0`$. At $`TT_N`$ we have $`\omega _𝐐=0`$ and, separating the condensation part $`C`$,
$$C_𝐫=\frac{1}{N}\underset{𝐪(𝐐)}{}C_𝐪\text{e}^{i\mathrm{𝐪𝐫}}+C\text{e}^{i\mathrm{𝐐𝐫}},$$
(21)
where $`C`$ results from (21) with $`𝐫=0`$ employing the sum rule $`C_{0,0,0}=\frac{1}{2}`$. Then the staggered magnetization $`m`$ is calculated as
$$m^2=\frac{1}{N}\underset{𝐫}{}𝐒_0𝐒_𝐫\text{e}^{i\mathrm{𝐐𝐫}}=\frac{3}{2}C.$$
(22)
The theory has 14 quantities to be determined self-consistently (9 correlation functions in $`\omega _𝐪^2`$, $`m`$, and 4 vertex parameters) and 13 self-consistency equations (10 Eqs. (21) including $`C_{0,0,0}=\frac{1}{2}`$, the LRO condition $`\omega _𝐐=0`$, and Eqs. (19) and (20)). If there is no LRO, we have $`\omega _𝐐>0`$, and the number of quantities and equations is reduced by one. As an additional condition for determining the free $`\alpha `$ parameter at $`T=0`$, we adjust the ground-state energy per site which we compose approximately as $`ϵ(R_y,R_z)=ϵ(R_y,0)+ϵ(0,R_z)ϵ(0,0)`$, where $`ϵ(R_y,0)`$ (and $`ϵ(0,R_z)`$) is taken from the Ising-expansion results by Affleck et al. for the 2D spatially anisotropic Heisenberg model, and $`ϵ(0,0)=0.4431`$ is the Bethe-ansatz value. This approximation is suggested to be good at least for $`R_zR_y`$ (or $`R_yR_z`$). To get an additional condition also at finite temperatures, where $`ϵ`$ data are not available and all vertex parameters are temperature dependent, we assume, following Refs. and , the ratio
$$r_\alpha (T)\frac{\alpha _2(T)1}{\alpha _1^x(T)1}=r_\alpha (0)$$
(23)
as temperature independent.
## III Ground-state properties
In Fig. 1 our results for the zero-temperature staggered magnetization $`m_0m(T=0)`$ as a function of $`R_y`$ and $`R_z`$ are shown. They indicate an order-disorder transition at the phase boundary $`R_{z,c}(R_y)`$ or $`R_{y,c}(R_z)`$ (cf. inset). For $`R_z=0`$ we get the critical ratio $`R_{y,c}(0)0.24`$ which was already found in Ref. . In that paper the suppression of LRO below the finite value of $`R_{y,c}`$ was interpreted, in combination with ED data, as indication of a rather sharp crossover in the spatial dependence of the spin correlation functions in the LRO phase at the coupling ratio $`R_{y,0}0.2`$. The finite value of $`R_{y,c}`$, however, seems to be due to the approximations in our theory, since there are strong indications for $`R_{y,c}=0`$ (see Ref. ). Accordingly, we cannot explain the tiny magnetic moments of Sr<sub>2</sub>CuO<sub>3</sub> and Ca<sub>2</sub>CuO<sub>3</sub>, since for $`R_y1`$, we have $`m=0`$. This result is just opposite to the overestimation of $`m`$ by all previous spin-wave theories. As seen in the phase diagram (inset of Fig. 1), the inclusion of the interplane coupling $`R_z`$ stabilizes the LRO, where this effect is quite considerable even at very small values of $`R_z`$.
Figure 2 exhibits some short-ranged spin correlation functions at $`T=0`$. For $`R_z=0`$, in Ref. the correlators $`C_{1,0,0},C_{0,1,0}`$, and $`C_{1,1,0}`$ as functions of $`R_y`$ were found to agree well with the ED data. For $`R_z=0.02`$ (cf. Fig. 2) our results deviate only slightly from those at $`R_z=0`$. The sign changes and magnitudes of $`C_𝐫`$ reflect the AFM SRO. In the limit $`R_y0`$ the correlations between the $`x`$-$`z`$-planes vanish. At $`R_z>R_{z,c}(0)0.24`$ the LRO enhances the inter-$`x`$-$`z`$ plane correlators and results in their sharp drop towards their limiting value $`C\text{e}^{i\mathrm{𝐐𝐫}}`$ as $`R_y0`$. This is visible in the data for $`R_z=0.35`$ in Fig. 2.
## IV Finite-temperature results
At nonzero temperatures we have solved the self-consistency equations (21) supplemented by the conditions (19), (20), and (23) to obtain the magnetization $`m(T)`$, the Néel temperature ($`m(T_N)=0`$), the static spin susceptibility, and the anisotropic correlation lengths.
In Fig. 3 the Néel temperature is plotted as a function of $`R_z`$. For $`R_z=0`$ we get $`T_N=0`$ (see Ref. ), in agreement with the Mermin-Wagner theorem. The increase of $`T_N`$ with $`R_z`$ is governed by the intra-$`x`$-$`y`$ plane anisotropy. At a fixed value of $`R_z`$, the decrease of $`T_N`$ with decreasing $`R_y`$ is in accordance with the reduced zero-temperature magnetization (cf. Fig. 1). Comparing our results for $`R_y=1`$ with previous RPA/mean-field approaches (see Table I), we ascribe the reduction of $`T_N`$ as compared with Refs. , and to the improved description of SRO. That is, the LRO is suppressed in favor of a paramagnetic phase with pronounced AFM SRO. If $`R_z`$ is fit to the Néel temperatures of real systems, the strong overestimation of $`T_N`$ by previous theories results in very small values of the interplane coupling. In our approach the resulting $`R_z`$ values turn out to be higher. Considering La<sub>2</sub>CuO<sub>4</sub> with $`T_N=325`$K and putting $`J=130`$meV ($`JJ_x=J_y`$) or $`J=117`$meV, we obtain $`R_z10^3`$ or $`R_z1.6\times 10^3`$, respectively, in contrast to $`R_z<10^4`$ according to Refs. and . For Ca<sub>0.85</sub>Sr<sub>0.15</sub>CuO<sub>2</sub> ($`T_N=540`$K, $`J=125`$meV) we get $`R_z1.2\times 10^2`$ as compared with $`R_z2.5\times 10^2`$ obtained from a fit of the low-temperature magnetization data.
Figure 4 shows the temperature dependence of the staggered magnetization at $`R_y=1`$ (for the zero-temperature values, compare with Fig. 1). The shape of the normalized curve $`m/m_0`$ versus $`T/T_N`$ (see inset) depends on the single parameter $`R_z`$ and is similar to that found in previous spin-wave theories. At low enough temperatures the system exhibits 3D behavior, so that the decrease of $`m`$ follows a $`T^2`$ law. This was also observed by NMR experiments on La<sub>2</sub>CuO<sub>4</sub> ($`T_N=312`$K) yielding $`m/m_0=1a(T/T_N)^2`$ with $`a=0.67`$ for $`T100`$K. The NMR data is indicated in the inset of Fig. 4 (marked by a bold curve) and agrees well with our theory for $`R_z=10^3`$ (as estimated above). For temperatures close to $`T_N`$ our numerical results for $`m(T)`$ are described by the law $`m(T)(1T/T_N)^{1/2}`$. The square-root temperature behavior agrees with the findings of Refs. , and with the neutron scattering data on La<sub>2</sub>CuO<sub>4</sub>, but contradicts the result of Ref. ($`m1T/T_N`$).
Considering the AFM correlation lengths above $`T_N`$ and for $`R_y=1`$, the expansion of $`\chi (𝐪)`$ around $`𝐐`$, $`\chi (𝐪)=\chi (𝐐)\left[1+\xi _{xy}^2(k_x^2+k_y^2)+\xi _z^2k_z^2\right]^1`$ with $`𝐤=𝐪𝐐`$, yields the intraplane correlation length
$`\xi _{xy}^2`$ $`=`$ $`\omega _𝐐^2[{\displaystyle \frac{1}{2}}+11\alpha _1^xC_{1,0,0}+\alpha _2(C_{2,0,0}+2C_{1,1,0})+`$ (25)
$`+2R_z(\alpha _1^xC_{1,0,0}+2\alpha _1^zC_{0,0,1}+\alpha _2C_{1,0,1})]{\displaystyle \frac{2C_{1,0,0}}{M_𝐐^{(1)}}}`$
and the interplane correlation length
$`\xi _z^2`$ $`=`$ $`R_z\omega _𝐐^2[4(2\alpha _1^xC_{1,0,0}+\alpha _1^zC_{0,0,1}+\alpha _2C_{1,0,1})+`$ (27)
$`+R_z({\displaystyle \frac{1}{2}}+5\alpha _1^zC_{0,0,1}+\alpha _2C_{0,0,2})]{\displaystyle \frac{2R_zC_{0,0,1}}{M_𝐐^{(1)}}}.`$
In Fig. 5 the influence of the interplane coupling on the temperature dependence of $`\xi _{xy}^1`$ and $`\xi _z^1`$ (inset) is shown. For comparison, the intraplane correlation length at $`R_z=0`$ (see also Ref. ) is plotted, where the low-temperature expansion $`\xi _{xy}=2(2\alpha _1^x|C_{1,0,0}(0)|)^{1/2}T^1\mathrm{exp}[2\pi \alpha _1^xm_0^2/(3T)]`$ holds up to $`T=0.2`$ within a deviation of about $`6\%`$ from the full temperature dependence calculated by Eq.(25). For $`R_z>0`$ the correlation lengths diverge at $`T_N`$, since the gap $`\omega _𝐐`$ closes as $`T`$ approaches $`T_N`$ from above. In the vicinity of $`T_N`$, $`\xi _{xy}^1`$ and $`\xi _z^1`$ behave as $`TT_N`$ also found by previous mean-field approaches.
Let us compare our results for the intraplane correlation length with the neutron-scattering data on La<sub>2</sub>CuO<sub>4</sub> in the range 340K$`T820`$K shown in Fig. 6. Taking $`J`$ as obtained previously from a least-squares fit of $`\xi _{xy}`$ in the 2D model ($`a=3.79`$Å), $`J=117`$meV, for $`T>500`$K and $`R_z3.5\times 10^3`$ we get a good quantitative agreement with experiments. In Ref. the deviation of the theory for $`R_z=0`$ and $`T<500`$K from the experimental data was ascribed to the appearance of the preexponential factor $`T^1`$ in the low-temperature expansion of $`\xi _{xy}`$ which is an artifact of our mean-field approach. However, this deviation may be reduced by the inclusion of the interplane coupling, since $`\xi _{xy}^1(T_N)=0`$. For $`T_N=325`$K we get $`R_z1.6\times 10^3`$ (see above, Fig. 3), and the theoretical $`\xi _{xy}^1`$ curve lies between the $`R_z=0`$ result and the experiments. The discrepancy between the theoretical and experimental low-temperature correlation lengths may be further reduced by the choice of higher $`R_z`$ values. Taking, for example, $`R_z=3.4\times 10^3`$, we get a very good quantitative agreement (cf. Fig. 6) down to 360K; however, the Néel temperature turns out to be somewhat higher ($`T_N=353`$K).
Finally, we consider the uniform static spin susceptibility $`\chi (T)=lim_{𝐪\mathrm{𝟎}}\chi (𝐪)`$. In Fig. 7 the anisotropy effects on the temperature dependence are demonstrated. For $`R_z=0`$ and a strong intraplane anisotropy ($`R_y<0.2`$) the minimum of $`\chi (T)`$ at a finite temperature, being an artifact of our approach, may signal the crossover in the spatial dependence of the spin correlation functions at $`R_{y,0}0.2`$ as was discussed in Sec. III. Note that such a minimum in the 1D model ($`R_y=0`$) was also found in Ref. . At $`R_y>0.2`$, the increase of $`\chi `$ with temperature, the maximum at $`T_{max}`$ near the exchange energy $`J_x=1`$ (see inset), and the crossover to the high-temperature Curie-Weiss behavior are due to the decrease of AFM SRO with increasing temperature (cf. Ref. ). Let us point out that the susceptibility maximum is totally missed in RPA theories. With increasing $`R_y`$, we obtain an increase of $`T_{max}`$ which agrees with a general tendency found in various spin–$`1/2`$ Heisenberg models and analyzed in Ref. . For comparison, the exact values at $`R_y=0`$ and $`R_y=1`$ are given by $`T_{max}=0.64`$ and $`T_{max}=0.94`$, respectively. Since our theory allows the calculation of $`T_{max}`$ at any spatial anisotropy, it may provide a reliable interpretation of experimental data on low-dimensional spin systems. Considering the maximum spin susceptibility $`\chi _{max}=\chi (T_{max})`$, again our results are in accordance with the general behavior: $`\chi _{max}`$ increases with decreasing $`T_{max}`$, i.e. with decreasing $`R_y`$. Concerning the influence of the interplane coupling, the enhancement of the low-temperature susceptibility by $`R_z`$ may be explained by the weakening of the SRO effect in higher dimensions. As seen from Figs. 7 and 3, the uniform susceptibility has no singularity at the Néel temperature, contrary to the RPA result of Ref. revealing a peak of $`\chi (T)`$ at $`T_N`$. Concerning the maximum in $`\chi (T)`$ of La<sub>2</sub>CuO<sub>4</sub>, we get $`T_{max}=1.19J=1615`$K (cf. Fig. 7, $`J=117`$meV). This value roughly agrees with the estimate given by Johnston, $`T_{max}=1460K`$, by means of a scaling analysis of the susceptibility data below 800K.
## V Summary
In this paper we have extended the spin-rotation-invariant Green’s-function theory of magnetic LRO and SRO in 2D Heisenberg models to the 3D Heisenberg antiferromagnet with arbitrary spatial anisotropy. Our theory provides a satisfactory interpolation between the low-temperature and high-temperature behavior, where the temperature dependent SRO, described in term of two-spin correlation functions, is adequately taken into account. The main results are summarized as follows.
(i) The incorporation of SRO results in a strong suppression of Néel order with increasing anisotropy and in a reduced Néel temperature as compared with previous spin-wave approaches.
(ii) The temperature dependence of the uniform static spin susceptibility reveals a maximum in the short-range ordered paramagnetic phase and a crossover to the Curie-Weiss law. The position of the maximum is influenced by the spatial anisotropy.
(iii) Comparing the theory with experiments on the magnetization and correlation length of La<sub>2</sub>CuO<sub>4</sub>, a good quantitative agreement is found.
From the results of our theory we conclude that the application of this approach to extended Heisenberg models (anisotropy in spin space, frustration) may be promising to describe the SRO effects on the unconventional magnetic properties of real low-dimensional spin systems.
Acknowledgments: The authors, especially L. Siurakshina, are very grateful to the DFG for financial support. Additional support by the Max-Planck society and the INTAS organisation (INTAS-97-1106) is acknowledged. The authors thank S.-L. Drechsler for many useful discussions.
TABLE I. Néel temperature $`T_N/J_x`$ at $`R_y=1`$ compared with other approaches
| $`\mathrm{log}_{10}(R_z)`$ | Fig. 3 | Ref. | Ref. | Ref. |
| --- | --- | --- | --- | --- |
| 4 | 0.17 | 0.48 | 0.29 | 0.25 |
| 3 | 0.22 | 0.65 | 0.38 | 0.34 |
| 2 | 0.36 | 0.80 | 0.54 | 0.47 |
| 1 | 0.56 | 1.15 | | 0.68 |
Figures
Fig. 1. Staggered magnetization at $`T=0`$ as a function of spatial anisotropy. The inset shows the stability region of Néel order.
Fig. 2. Spin correlation functions at $`T=0`$ for different spatial anisotropies.
Fig. 3. Néel temperature as a function of $`R_z=J_z/J_x`$.
Fig. 4. Staggered magnetization vs. temperature for $`R_y=1`$. The inset shows the $`R_z`$ dependence of the normalized curves compared with the NMR data on La<sub>2</sub>CuO<sub>4</sub> (bold curve).
Fig. 5. Inverse antiferromagnetic correlation lengths within ($`\xi _{xy}^1`$) and between the $`x`$-$`y`$ planes ($`\xi _z^1`$, see inset) for $`R_y=1`$.
Fig. 6. Inverse antiferromagnetic intraplane correlation length in La<sub>2</sub>CuO<sub>4</sub> obtained by the neutron-scattering experiments of Ref. and from the theory ($`R_y=1`$) for different $`R_z`$ values.
Fig. 7. Uniform static spin susceptibility vs. $`T`$. The inset exhibits the position $`T_{max}`$ of the maximum in $`\chi (T)`$ vs. $`R_y`$. |
warning/0001/hep-ph0001120.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In this note we study scattering of virtual and real photons
$$\gamma ^{}(q)+\gamma ^{}(p)X.$$
(1)
in the high-energy regime of large Regge parameter $`1/x`$ which depends on virtualities of photons as
$$\frac{1}{x}=\frac{W^2+Q^2+P^2}{Q^2+P^2+\mu ^2}1.$$
(2)
and has correct parton model limit if either $`Q^2P^2`$ or $`P^2Q^2`$. In eq.(2) $`W^2=(q+p)^2`$ is the center-of-mass energy squared of colliding space-like photons $`\gamma ^{}(q)`$ and $`\gamma ^{}(p)`$ with virtualities $`q^2=Q^2`$ and $`p^2=P^2`$, respectively.
The recent strong theoretical and experimental (see also a compilation in ) interest in high-energy $`\gamma ^{}\gamma ^{},\gamma ^{}\gamma ,\gamma \gamma `$ scattering stems from the fact that virtualities of photons give a handle on the size of color dipoles in the beam and target photons and, eventually, on short distance properties of the QCD pomeron exchange. For earlier development of the subject see the pioneering paper .
As noticed by Fadin, Kuraev and Lipatov in 1975 and discussed in more detail by Lipatov in the incorporation of asymptotic freedom into the BFKL equation makes the QCD pomeron a series of isolated poles in the angular momentum plane. The contribution of each isolated pole to the high-energy scattering amplitude satisfies the familiar Regge factorization . In we reformulated the consequences of the Regge factorization in our color dipole (CD) approach to the BFKL pomeron. In this communication we address several closely related issues in photon-photon scattering in the Regge regime (2) which can be tested at LEP200 and Next Linear Collider (NLC).
First, following our early work we discuss how the color dipole (CD) BFKL-Regge factorization leads to the parameter-free predictions for total cross sections of $`\gamma ^{}\gamma ^{}`$, $`\gamma ^{}\gamma `$, $`\gamma \gamma `$ scattering. We find good agreement with the recent experimental data from the L3 and OPAL experiments at LEP .
Second, we discuss the interplay of soft and hard dynamics of the vacuum exchange and comment on the onset of the soft plus rightmost hard BFKL-pole dominance in $`\gamma ^{}\gamma ^{}`$ diffractive scattering. The nodal properties of eigenfunctions of the CD BFKL equation suggest an interesting possibility of decoupling of sub-leading BFKL singularities when the virtuality of one or both of photons is in the broad vicinity of $`Q^220`$ GeV<sup>2</sup>. Hence, very efficient the leading hard plus soft approximation (LHSA) advocated by us previously in .
Third, we discuss the impact of running CD BFKL on contentious issue of testing the factorization properties of photon-photon scattering in the $`Q^2,P^2`$-plane which has earlier been discussed only in the approximation of $`\alpha _S=`$const to the BFKL equation (for the color dipole picture in the $`\alpha _S=`$const approximation see ). Our finding is that the non-perturbative corrections break down the Regge factorization. The experimental observation of this phenomenon would contribute to better understanding of the non-perturbative dynamics of high-energy processes.
## 2 Overview of color dipole BFKL-Regge factorization
In the color dipole basis the beam-target scattering is viewed as transition of $`\gamma ^{}`$ into quark-antiquark pair and interaction of the beam ($`A`$) and target ($`B`$) color dipoles of the flavor $`A,B=u,d,s,c`$. As a fundamental quantity we use the forward dipole scattering amplitude and/or the dipole-dipole cross section $`\sigma (x,𝐫,𝐫^{})`$. Once $`\sigma (x,𝐫,𝐫^{})`$ is known the total cross section of $`AB`$ scattering $`\sigma ^{AB}(x)`$ is calculated as
$$\sigma ^{bt}(x)=𝑑zd^2𝐫𝑑z^{}d^2𝐫^{}|\mathrm{\Psi }_A(z,𝐫)|^2|\mathrm{\Psi }_B(z^{},𝐫^{})|^2\sigma (x,𝐫,𝐫^{}),$$
(3)
where $`𝐫`$ and $`𝐫^{}`$ are the two-dimensional vectors in the impact parameter plane. In the color dipole factorization formula (3) the dipole-dipole cross section $`\sigma (x,𝐫,𝐫^{})`$ is beam-target symmetric and universal for all beams and targets, the beam and target dependence is concentrated in probabilities $`|\mathrm{\Psi }_A(z,𝐫)|^2`$ and $`|\mathrm{\Psi }_B(z^{},𝐫^{})|^2`$ to find a color dipole, $`𝐫`$ and $`𝐫^{}`$ in the beam and target, respectively. Hereafter we focus on cross sections averaged over polarizations of the beam and target photons, in this case only the term $`n=0`$ of the Fourier series
$$\sigma (x,𝐫,𝐫^{})=\underset{n=0}{\overset{\mathrm{}}{}}\sigma _n(x,r,r^{})\mathrm{exp}(in\phi ),$$
(4)
where $`\phi `$ is an azimuthal angle between $`𝐫`$ and $`𝐫^{}`$, contributes in (3).
Fadin, Kuraev and Lipatov noticed in 1975 , see also Lipatov’s extensive discussion that the incorporation of asymptotic freedom into the BFKL equation makes the QCD pomeron a series of isolated poles in the angular momentum plane. The contribution of the each pole to scattering amplitudes satisfies the standard Regge-factorization , which in the CD basis implies the CD BFKL-Regge expansion for the vacuum exchange dipole-dipole cross section
$$\sigma (x,r,r^{})=\underset{m}{}C_m\sigma _m(r)\sigma _m(r^{})\left(\frac{x_0}{x}\right)^{\mathrm{\Delta }_m}.$$
(5)
Here the dipole cross section $`\sigma _m(r)`$ is an eigen-function of the CD BFKL equation
$$\frac{\sigma _m(x,r)}{\mathrm{log}(1/x)}=𝒦\sigma _m(x,r)=\mathrm{\Delta }_m\sigma _m(x,r),$$
(6)
with eigen value (intercept) $`\mathrm{\Delta }_m`$. Arguably, for transition of $`\gamma ^{}`$ into heavy flavors, $`A=c,b,\mathrm{}`$, the hardness scale is set by $`Q^2+4m_A^2`$, for light flavors $`Q^2+m_\rho ^2`$ is a sensible choice which leads to the correct value of Regge parameter in the photoproduction regime, $`Q^20`$.
Hence, for the light-light transition we evaluate the Regge-parameter (2) with $`\mu ^2=m_\rho ^2`$, for the light-charm contribution we take $`\mu ^2=4m_c^2`$ and for the charm-charm contribution we take $`\mu ^2=8m_c^2`$
For the details on CD formulation of the BFKL equation, infrared regularization by finite propagation radius $`R_c`$ for perturbative gluons and freezing of strong coupling at large distances, the choice of the physically motivated boundary condition for the hard BFKL evolution and description of eigenfunctions we refer to our early works , the successful application of CD BFKL-Regge expansion to the proton and pion structure functions (SF) and evaluation of hard-pomeron contribution to the rise of hadronic and real photo-absorption cross sections is found in . We only recapitulate the salient features of the formalism essential for the present discussion.
There is a useful analogy between the intercept $`\mathrm{\Delta }=\alpha (0)1`$ and binding energy for the bound state problem for the Schrödinger equation. The eigenfunction $`\sigma _0(r)`$ for the rightmost hard BFKL pole (ground state) corresponding to the largest intercept $`\mathrm{\Delta }_0\mathrm{\Delta }_{𝐈𝐏}`$ is node free. The eigenfunctions $`\sigma _m(r)`$ for excited states with $`m`$ radial nodes have intercept $`\mathrm{\Delta }_m<\mathrm{\Delta }_{𝐈𝐏}`$. Our choice of $`R_c=0.27`$fm yields for the rightmost hard BFKL pole the intercept $`\mathrm{\Delta }_{𝐈𝐏}=0.4`$ , for sub-leading hard poles $`\mathrm{\Delta }_m\mathrm{\Delta }_0/(m+1)`$. The node of $`\sigma _1(r)`$ is located at $`r=r_10.050.06\mathrm{fm}`$, for larger $`m`$ the rightmost nodes move to a somewhat larger $`r`$ and accumulate at $`r0.1\mathrm{fm}`$, for the more detailed description of the nodal structure of $`\sigma _m(r)`$ see . Here we only emphasize that for solutions with $`m3`$ the third and higher nodes are located at a very small $`r`$ way beyond the resolution scale $`1/\sqrt{Q^2}`$ of foreseeable deep inelastic scattering (DIS) experiments. Notice that the Regge cut in the complex angular momentum plane found in the much discussed approximation $`\alpha _S=`$const resembles an infinite, and continuous, sequence of poles. In the counterpart of our CD BFKL-Regge expansion (5) for approximation $`\alpha _S=`$const the intercept $`\mathrm{\Delta }_m`$ would be a continuous parameter in contrast to a discrete spectrum for standard running $`\alpha _S`$.
Because the BFKL equation sums cross sections of production of multigluon final states, the perturbative two-gluon Born approximation is an arguably natural boundary condition. This leaves the starting point $`x_0`$ as the only free parameter which fixes completely the result of the hard BFKL evolution for dipole-dipole cross section. We follow the choice $`x_0=0.03`$ made in . The very ambitious program of description of $`F_{2p}(x,Q^2)`$ starting from this, perhaps excessively restrictive, perturbative two-gluon boundary condition has been launched by us in and met with remarkable phenomenological success .
Because in the attainable region of $`r`$ the sub-leading solutions $`m3`$ can not be resolved and they all have similar intercepts $`\mathrm{\Delta }_m1`$, in practical evaluation of $`\sigma ^{AB}`$ we can truncate expansion (5) at $`m=3`$ lumping in the term with $`m=3`$ contributions of all singularities with $`m3`$. Specifically, if we endow
$$\sigma _3(r)=\sigma _{Born}(r)\underset{m=0}{\overset{2}{}}\sigma _m(r)$$
(7)
with the effective intercept $`\mathrm{\Delta }_3=0.06`$ the truncated expansion reproduces the numerical solution $`\sigma (x,r)`$ of our CD BFKL equation in the wide range of dipole sizes $`10^3\text{ }<r\text{ }<10`$ fm with accuracy $`10\%`$ even at moderately small $`x`$. Such a truncation can be justified a posteriori if such a contribution from $`m3`$ turns out to be a small correction, which will indeed be the case at small $`x`$.
Whereas scattering of small dipoles $`r\text{ }<R_c`$ is dominated by the exchange of perturbative gluons, interaction of large dipoles with the proton target has been modeled in Ref. by the non-perturbative soft pomeron with intercept $`\alpha _{\mathrm{soft}}(0)1=\mathrm{\Delta }_{\mathrm{soft}}=0`$. Then the extra term $`\sigma _{\mathrm{soft}}(r,r^{})`$ must be added in the r.h.s. of expansion (5).
From the early phenomenology of DIS and diffractive vector meson production off the proton target we only know the parameterization of $`\sigma _{\mathrm{soft}}(r,r^{})`$ when one of the dipoles is definitely large, of the order of the proton size. Evaluation of the soft contribution to $`\gamma ^{}\gamma ^{}`$ scattering when both dipoles are small inevitably introduces model dependence. Modelling of soft contribution by the exchange by two nonperturbative gluons suggests
$$\sigma _{\mathrm{soft}}(r,r^{})\frac{r^2r^2}{(r^2+r^2)}$$
and the non-factorizable cross section of the form
$$\sigma _{\mathrm{soft}}^{\gamma ^{}\gamma ^{}}(Q^2,P^2)\frac{1}{Q^2+P^2}.$$
Similar non-perturbative $`\sigma _{\mathrm{soft}}`$ is found in the soft pomeron models . The explicit parameterization is found in Appendix.
Finally, at moderately small values of $`x`$ the above described $`t`$-channel gluon tower exchange must be complemented by the $`t`$-channel $`q\overline{q}`$ exchange often associated with DIS off vector mesons (hadronic component) and off the perturbative (point-like) $`q\overline{q}`$-component of the target photon wave function. We add corresponding corrections only to the real photon structure function $`F_{2\gamma }(x,Q^2)`$ to estimate the interplay of vacuum and non-vacuum exchanges in the currently accessible kinematical region of not very small $`x`$. In all other cases of interst we concentrate on the pure vacuum exchange at $`x\text{ }<x_0`$ where the non-vacuum corrections are negligible small.
In our evaluation of the box diagram contribution to $`F_{2\gamma }^{pl}(x,Q^2)`$ which is
$$F_{2\gamma }^{pl}(x,Q^2)=\frac{3\alpha _{em}}{\pi }\underset{q=udsc}{}e_q^4x\left\{\left[x^2+(1x)^2\right]\mathrm{log}\frac{Q^2(1x)}{xQ_q^2}+8x(1x)1\right\}$$
(8)
we take the $`\rho `$-meson mass as the lower cut-off for the light-flavor-loop integral, $`Q_q^2=m_\rho ^2`$ for $`q=u,d,s`$, and the charm quark mass for the $`c`$-loop, $`Q_c^2=m_c^2`$. In eq.(8) $`e_q`$ is a quark charge.
To describe the hadronic component of $`F_{2\gamma }`$ we take the coherent mixture of $`\rho ^0`$ and $`\omega `$ mesons . Being supplemented with quite standard assumptions on the vector meson valence quark density this gives
$$F_{2\gamma }^{had}(x)=\frac{\alpha _{em}}{12}[4(g_\rho +g_\omega )^2+(g_\rho g_\omega )^2)]\sqrt{x}(1x),$$
(9)
where the coupling constants $`g_V^2=4\pi /f_V^2`$, in the Fock state expansion
$$|\gamma ^{had}=\frac{e}{f_\rho }|\rho +\frac{e}{f_\omega }|\omega +\mathrm{}$$
are as follows $`g_\rho ^2=0.5`$ and $`g_\omega ^2=0.043`$ . We neglect the $`Q^2`$-evolution which, at reasonable values of lower scale, is a small correction on the interval $`1.9Q^25`$ GeV<sup>2</sup> where the small-$`x`$ data on $`F_{2\gamma }`$ were taken.
Then, combining (5) and (3) and adding in the soft and non-vacuum components, we obtain our principal result for virtual-virtual scattering ($`m=0,1,2,3`$, $`A,B=u,d,s,c`$)
$`\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}(x,Q^2,P^2)={\displaystyle \frac{(4\pi ^2\alpha _{em})^2}{Q^2P^2}}{\displaystyle \underset{m}{}}C_m{\displaystyle \underset{A,B}{}}f_m^A(Q^2)f_m^B(P^2)\left({\displaystyle \frac{3x_0}{2x_{AB}}}\right)^{\mathrm{\Delta }_m}`$
$`+\sigma _{\mathrm{soft}}^{\gamma ^{}\gamma ^{}}(x,Q^2,P^2).`$ (10)
To make explicite the scale dependence discussed in Introduction we provide $`\mu `$ and $`x`$ defined by the eq.(2) with two indices, $`A`$ and $`B`$, pointing out the flavor of both the beam and target dipoles: $`\mu _{AB}^2=m_\rho ^2`$ for $`A,B=u,d,s`$ while $`\mu _{AB}^2=4m_c^2`$ if either $`A=c`$ or $`B=c`$ and the second dipole is made of light quarks and $`\mu _{AB}^2=8m_c^2`$ if $`A=B=c`$.
For the DIS off real (quasireal) photons, $`P^20`$, we have $`(A=u,d,s,c)`$
$$F_{2\gamma }(x,Q^2)=\underset{m}{}A_m^\gamma \underset{A}{}f_m^A(Q^2)\left(\frac{3}{2}\frac{x_0}{x_A}\right)^{\mathrm{\Delta }_m}+F_{2\gamma }^{\mathrm{soft}}(x,Q^2)+F_{2\gamma }^{\mathrm{nvac}}(x,Q^2),$$
(11)
where
$$x_A=\frac{Q^2+\mu _A^2}{W^2+Q^2}$$
and $`\mu _A^2=m_\rho ^2`$ for $`A=u,d,s`$ while $`\mu _A^2=4m_c^2`$ for $`A=c`$. The $`c\overline{c}`$ component of the target photon wave function is strongly suppressed at $`P^20`$ and for all the practical purposes can be neglected as well as $`c\overline{c}`$ content of the target proton. This observation simplifies the factorization relation (11) for the real photon structure function. In eq.(11) the non-vacuum component denoted by $`F_{2\gamma }^{\mathrm{nvac}}`$ is
$$F_{2\gamma }^{\mathrm{nvac}}(x,Q^2)=F_{2\gamma }^{\mathrm{had}}(x,Q^2)+F_{2\gamma }^{\mathrm{pl}}(x,Q^2)$$
(12)
and the cross sections
$$\sigma _m^\gamma ^{}(Q^2)=\gamma _T^{}|\sigma _m(r)|\gamma _T^{}+\gamma _L^{}|\sigma _m(r)|\gamma _L^{}.$$
(13)
are calculated with the well known color dipole distributions in the transverse (T) and longitudinal (L) photon of virtuality $`Q^2`$ derived in , and the eigen SFs are defined as usual:
$$f_m(Q^2)=\frac{Q^2}{4\pi ^2\alpha _{\mathrm{em}}}\sigma _m^\gamma ^{}(Q^2).$$
(14)
The factor $`\frac{3}{2}`$ in the Regge parameter derives from the point that in a scattering of color dipole on the photon the effective dipole-dipole collision energy is $`\frac{3}{2}`$ of that in the reference scattering of color dipole on the three-quark nucleon at the same total c.m.s. energy $`W`$. The analytical formulas for the eigen-SFs $`f_m(Q^2)`$ and $`f_m^c(Q^2)`$ are found in the Appendix. Here as well as in all our previous calculations we put $`m_c=1.5`$ GeV. We do not need any new parameters compared to those used in the description of DIS and real photoabsorption on protons (for an alternative approach see ), the results for the expansion parameters $`A_m^\gamma =C_m\sigma _m^\gamma `$, $`C_m=1/\sigma _m^p`$ and $`\sigma _m^\gamma \sigma _m^\gamma ^{}(0)`$ are summarized in the Table 1.
We recall that because of the diffusion in color dipole space, exchange by perturbative gluons contributes also to interaction of large dipoles $`r>R_c`$ . However at moderately large Regge parameter this hard interaction driven effect is still small. For this reason in what follows we refer to terms $`m=0,1,2,3`$ as hard contribution as opposed to the genuine soft interaction.
Table 1. CD BFKL-Regge expansion parameters.
| $`m`$ | $`\mathrm{\Delta }_m`$ | $`\sigma _m^p,\mathrm{mb}`$ | $`C_m,\mathrm{mb}^1`$ | $`A_m^\gamma /\alpha _{\mathrm{em}}`$ | $`\sigma _m^\gamma ,\mu \mathrm{b}`$ | $`\sigma _m^{\gamma \gamma },\mathrm{nb}`$ |
| --- | --- | --- | --- | --- | --- | --- |
| 0 | 0.402 | 1.243 | 0.804 | 0.746 | 6.767 | 36.84 |
| 1 | 0.220 | 0.462 | 2.166 | 0.559 | 1.885 | 7.69 |
| 2 | 0.148 | 0.374 | 2.674 | 0.484 | 1.320 | 4.65 |
| 3 | 0.06 | 3.028 | 0.330 | 0.428 | 9.456 | 29.53 |
| soft | 0. | 31.19 | 0.0321 | 0.351 | 79.81 | 204.2 |
## 3 Isolating the soft plus rightmost hard BFKL pole in highly virtual-virtual $`\gamma ^{}\gamma ^{}`$ scattering
We start with the theoretically cleanest case of the highly virtual photons, $`P^2,Q^2`$1 GeV<sup>2</sup> and focus on the vacuum exchange component of the total cross section. The CD BFKL approach with asymptotic freedom predicts uniquely that subleading eigen SFs have a node at $`Q^220`$ GeV<sup>2</sup> in which region of $`Q^2`$ the rightmost hard pole contribution will dominate. This suppression of the subleading hard background is shown in Fig. 1, in which we plot the ratio ($`m=0,1,2,3,soft`$)
$$r_m(Q^2)=\frac{\sigma _m^{\gamma ^{}\gamma ^{}}(\frac{3}{2}x_0,Q^2,Q^2)}{\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}(\frac{3}{2}x_0,Q^2,Q^2)},$$
which defines the relative size of different contributions to $`\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}`$ at $`x=\frac{3}{2}x_0`$. At this value of $`x`$ the contribution of subleading hard BFKL poles remains marginal in a broad range of $`Q^2`$, although the contribution from the single-node component $`m=1`$ becomes substantial at $`Q^2\text{ }>10^3`$ GeV<sup>2</sup>.
The soft-pomeron exchange contributes substantially over all $`Q^2`$ and dominates at $`Q^2\text{ }<1`$ GeV<sup>2</sup>. However, at very large $`W100`$ GeV of the practical interest at LEP and LHC, such small values of $`Q^2`$ correspond to very small $`x`$, where the soft and subleading hard contributions are Regge suppressed by the factor $`\left(x/x_0\right)^{\mathrm{\Delta }_{𝐈𝐏}}`$ and $`\left(x/x_0\right)^{0.5\mathrm{\Delta }_{𝐈𝐏}}`$, respectively. The latter is clearly seen from Fig. 2 where the effective pomeron intercept
$$\mathrm{\Delta }_{eff}=\frac{\mathrm{log}\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}}{\mathrm{log}x}$$
(15)
is presented for the diagonal case $`Q^2=P^2`$ at three different values of $`W`$.
According to the results shown in Fig. 1 the dominance of the soft plus rightmost hard BFKL pomeron exchange in virtual-virtual $`\gamma ^{}\gamma ^{}`$ scattering holds in a very broad range of $`Q^2,P^2\text{ }<500`$ GeV<sup>2</sup> which nearly exhausts the interesting kinematical region at LEP200 and NLC. The quality of the leading hard pole plus soft approximation (LHSA) can be judged also from Fig. 3 for the diagonal case of $`Q^2=P^2`$, in which we show separately the soft component of the cross section (the dashed curve). The point that the contribution from subleading hard BFKL exchange is marginal is clear from the finding that approximation of soft-pomeron plus the rightmost hard BFKL exchange (LHSA) shown by long-dashed curve nearly exhausts the result from the complete CD BFKL-Regge expansion for vacuum exchange.
Recently the L3 collaboration reported the first experimental evaluation of the vacuum exchange in equal virtuality $`\gamma ^{}\gamma ^{}`$ scattering. Their procedure of subtraction of the non-vacuum reggeon and/or the Quark Parton Model contribution is described in , arguably the subtraction uncertainties are marginal within the present error bars. In Fig. 4 we compare our predictions to the L3 data. The experimental data and theoretical curves are shown vs. the variable $`Y=\mathrm{log}(W^2/\sqrt{Q^2P^2})`$. The virtuality of two photons varies in the range of $`1.2GeV^2<Q^2,P^2<9GeV^2`$ ($`Q^2,P^2=3.5`$GeV<sup>2</sup>) at $`\sqrt{s}91GeV`$ and $`2.5GeV^2<Q^2,P^2<35GeV^2`$ at $`\sqrt{s}183GeV`$ ($`Q^2,P^2=14`$GeV<sup>2</sup>). We applied to the theoretical cross sections the same averaging procedure as described in . The solid curve is a result of the complete BFKL-Regge expansion for the vacuum exchange, the long-dashed curve is a sum of the rightmost hard BFKL exchange and soft-pomeron exchange. Shown by the dashed line is the soft pomeron contribution. The agreement of our estimates with the experiment is good, the contribution of subleading hard BFKL exchange is negligible within the experimental error bars.
In Fig. 5 we compare our predictions for the vacuum exchange contribution to $`\sigma ^{\gamma ^{}\gamma ^{}}(Y)`$ with recent OPAL Collaboration measurements . In the applicability region of our approach corresponding to $`Y\text{ }>Y_0=\mathrm{log}(2/3x_0)3`$ the agreement with data is good. The discrepancy at smaller $`Y`$ may indicate significant non-vacuum contributions vanishing at large $`Y`$.
The early calculations of the perturbative vacuum component of $`\sigma ^{\gamma ^{}\gamma ^{}}`$ used the approximation $`\alpha _S=const`$ which predicts the $`P^2,Q^2`$-dependence different from our result for CD BFKL approach with running $`\alpha _S`$. Detailed comparison with numerical results by Brodsky, Hautmann and Soper (BHS) is reported by the L3 Collaboration , which founds that BHS formulas overpredict $`\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}`$ substantially. In the same perturbative fixed-$`\alpha _S`$ BFKL model with massive $`c`$-quark has been considered. At $`Q^2=14`$ GeV<sup>2</sup> and moderately small $`x`$, $`x\text{ }>3.10^2`$, the model is in agreement with the L3 data but at smaller $`x`$, already at $`x7.10^3`$, it substanially overpredicts $`\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}`$. At $`Q^2=3.5`$ GeV<sup>2</sup> the results are substantially above the L3 data over all $`x`$.
## 4 Virtual-real $`\gamma ^{}\gamma `$ scattering: the rightmost hard BFKL pole in the photon SF
The discussion of the photon SF follows closely that of the proton and pion SF’s in . Our normalization of eigen-functions is such that the vacuum (sea) contribution to the proton SF ($`m=\mathrm{soft},0,1,..,3`$)
$$F_{2p}(x,Q^2)=\underset{m}{}f_m(Q^2)\left(\frac{x_0}{x}\right)^{\mathrm{\Delta }_m}$$
(16)
has the CD BFKL-Regge expansion coefficients $`A_m^p=1`$. There is a fundamental point that the distribution of small-size color dipoles in the photon is enhanced compared to that in the proton which enhances the importance of the rightmost hard BFKL exchange. Indeed, closer inspection of expansion coefficients $`A_m^\gamma `$ shown in table 1 reveals that subleading hard BFKL exchanges are suppressed with respect to the leading one by the factor $``$ 1.5, whereas the soft-pomeron exchange contribution is suppressed by the factor $`2`$.
Our predictions for the photon SF are parameter-free and are presented in Fig. 6. At moderately small $`x0.1`$ there is a substantial non-vacuum reggeon exchange contribution from DIS off hadronic ($`q\overline{q}`$) component of the target photon wave function which can be regarded as well constrained by the large $`x`$ data. We use here the parameterizations presented above (eqs.(9,8, 12)). The solid curve shows the result from the complete BFKL-Regge expansion the soft-pomeron (the dashed curve) and quasi-valence (the dot-dashed curve) components included, the dotted curve shows the rightmost hard BFKL (LH) plus soft-pomeron (S) plus non-vacuum (NV) approximation (LHSNVA). A comparison of the solid and dotted curves shows clearly that subleading hard BFKL exchanges are numerically small in the experimentally interesting region of $`Q^2`$, the rightmost hard BFKL pole exhausts the hard vacuum contribution for $`2\text{ }<Q^2\text{ }<100`$ GeV<sup>2</sup>. The nodal properties of subleading hard BFKL SFs are clearly seen: LHSNVA underestimates $`F_{2\gamma }`$ slightly at $`Q^2\text{ }<10`$ GeV<sup>2</sup> and overestimates $`F_{2\gamma }`$ at $`Q^2\text{ }>50`$ GeV<sup>2</sup>. For still another illustration of the same nodal property of subleading hard components see Fig. 7 in which we show the vacuum component of virtual-real total cross section $`\sigma _{tot}^{\gamma ^{}\gamma }`$ as a function of $`Q^2`$ at fixed $`W`$. As seen from Fig. 1, the soft contribution rises towards small $`Q^2`$, but this rise is compensated to a large extent by the small-$`x`$ enhancement of the rightmost hard BFKL contribution by the large Regge factor $`\left(\frac{x_0}{x}\right)^{\mathrm{\Delta }_{𝐈𝐏}}`$. For this region the soft background (the dashed curve) remains marginal over the whole range of $`Q^2`$. Because of the node effect, the $`m=1`$ subleading component changes the sign and becomes quite substantial at very large $`Q^2`$ and moderately small $`x`$.
Recently the L3 and OPAL collaborations reported the first experimental data on the photon SF at sufficiently small-x . These data are shown in Fig. 6 and are in good agreement with the predictions from the CD BFKL-Regge expansion. A comparison with the long-dashed curve which is the sum of the rightmost hard BFKL and soft exchanges shows that the experimental data are in the region of $`x`$ and $`Q^2`$ still affected by non-vacuum reggeon (quasi-valence) exchange (dot-dashed line), going to smaller $`x`$ and larger $`Q^2`$ would improve the sensitivity to pure vacuum exchange greatly.
In order to give a crude idea on finite-energy effects at large $`x`$ and not so large values of the Regge parameter we stretch the theoretical curves a bit to $`x\text{ }>x_0`$ multiplying the BFKL-Regge expansion result by the purely phenomenological factor $`(1x)`$ motivated by the familiar behavior of the gluon SF of the photon $`(1x)^n`$ with the exponent $`n1`$.
## 5 The real-real $`\gamma \gamma `$ scattering
We recall that because of the well known BFKL diffusion in color dipole space, exchange by perturbative gluons contributes also to interaction of large dipoles $`r>R_c`$ . As discussed in this gives rise to a substantial rising component of hadronic and real photoabsorption cross sections and a scenario in which the observed rise of hadronic and real photon cross sections is entirely due to this intrusion of hard scattering. This is a motivation behind our choice of intercept $`\mathrm{\Delta }_{\mathrm{soft}}=0`$ for soft pomeron exchange. Furthermore, in order to make this picture quantitative one needs to invoke strong absorption/unitarization to tame too a rapid growth of large dipole component of hard BFKL the dipole cross section. The case of real-real $`\gamma \gamma `$ scattering is not an exception and the above discussed enhancement of small dipole configurations in photons compared to hadrons predicts uniquely that the hard BFKL exchange component of real-real $`\gamma \gamma `$ scattering will be enhanced compared to proton-proton and/or pion-proton scattering. This is clearly seen from table 1 in which we show the coefficients
$$\sigma _m^{\gamma \gamma }=\sigma _m^\gamma \sigma _m^\gamma C_m$$
(17)
of the expansion for the vacuum exchange component of the total $`\gamma \gamma `$ cross section ($`m=0,1,2,3,soft`$)
$$\sigma _{\mathrm{vac}}^{\gamma \gamma }=\underset{m}{}\sigma _m^{\gamma \gamma }\left(\frac{W^2x_0}{m_\rho ^2}\right)^{\mathrm{\Delta }_m}.$$
(18)
One has to look at the soft-hard hierarchy of $`\sigma _m^{\gamma \gamma }`$ and $`\sigma _m^\gamma ,\sigma _m^p`$ in the counterparts of (18) for $`\gamma p`$ and $`pp`$ scattering. This enhancement of hard BFKL exchange is confirmed by simplified vacuum pole plus non-vacuum reggeon exchange fits to real-real $`\gamma \gamma `$ total cross section: the found intercept of the effective vacuum pole $`ϵ^{\gamma \gamma }0.21`$ is much larger than $`ϵ0.095`$ from similar fits to the hadronic cross section data. In Fig. 8 we compare our predictions from the CD BFKL-Regge factorization for the single-vacuum exchange contribution to real-real $`\gamma \gamma `$ scattering with the recent experimental data from the OPAL collaboration and . The theoretical curves are in the right ballpark, but the truly quantitative discussion of total cross sections of soft processes requires better understanding of absorption/unitarization effects.
## 6 Regge factorization in $`\gamma ^{}\gamma ^{}`$ and $`\gamma \gamma `$ scattering
If the vacuum exchange were an isolated Regge pole, the well known Regge factorization would hold for asymptotic cross sections
$$\sigma _{tot}^{bb}\sigma _{tot}^{aa}=\sigma _{tot}^{ab}\sigma _{tot}^{ab}.$$
(19)
In the CD BFKL approach such a Regge factorization holds for each term in the BFKL-Regge expansion for vacuum exchange, but evidently the sum of factorized terms does not satisfy the factorization (19). One can hope for an approximate factorization only provided one single term to dominate in the BFKL-Regge expansion. Though corrections to the exact factorization still exist even for the single pole exchange because of the light $`q\overline{q}`$ and charm $`c\overline{c}`$ mass scale difference discussed above.
One such case is real-real $`\gamma \gamma `$ scattering dominated by soft-pomeron exchange (the factorization of the soft on-shell amplitudes though never proved gained strong support from the high-energy Regge-phenomenology). For this reason the CD BFKL-Regge expansion which reproduces well the vacuum exchange components of the $`pp`$ and $`\gamma p`$ scattering can not fail for the vacuum component in real-real $`\gamma \gamma `$ scattering. The rise of the contribution of hard-BFKL exchange breaks the Regge factorization relation
$$R_{\gamma \gamma }=\frac{\sigma _{\mathrm{vac}}^{\gamma \gamma }\sigma ^{pp}}{\sigma _{\mathrm{vac}}^{\gamma p}\sigma _{\mathrm{vac}}^{\gamma p}}=1,$$
(20)
which would restore at extremely high energies such that the rightmost hard BFKL exchange dominates. This property is illustrated in Fig. 9 where we show our evaluation of $`R`$ for single-vacuum component of total cross sections entering (19). At moderately high energies naive factorization breaks but the expected breaking is still weak, $`\text{ }<20\%`$. This curve must not be taken at face value for $`W\text{ }>`$0.1-1 TeV because of likely strong absorption effects, but the trend of $`R`$ being larger than unity and rising with energy should withstand unitarity effects. The second case is highly virtual-virtual $`\gamma ^{}\gamma ^{}`$ scattering. As we emphasized in section 3, here the CD BFKL approach predicts uniquely that because of the nodal property of subleading eigen SFs the the superposition of soft and rightmost hard BFKL poles dominate the vacuum exchange in a broad range of $`Q^2,P^2\text{ }<10^3`$ GeV<sup>2</sup>.
The above discussion suggests clearly that different cross sections must be taken at the same value of $`x^1=W^2/(Q^2+P^2)`$, in which case the vacuum components of $`\gamma ^{}\gamma ^{}`$ scattering at $`Q^2,P^24m_c^2`$ and $`xx_0`$ would satisfy
$$R_{\gamma ^{}\gamma ^{}}(x)=\frac{[\sigma ^{\gamma ^{}\gamma ^{}}(x,Q^2,P^2)]^2}{\sigma ^{\gamma ^{}\gamma ^{}}(x,Q^2,Q^2)\sigma ^{\gamma ^{}\gamma ^{}}(x,P^2,P^2)}=1.$$
(21)
In accordance to the results shown in Fig. 1, the soft exchanges break the factorization relation (21). The breaking is quite substantial at moderate $`x=0.01`$ (dotted line in Fig. 10), and breaking effects disappear rapidly, $`x^{\mathrm{\Delta }_0}`$, as $`x0`$. If the vacuum singularity were the Regge cut as is the case in approximation $`\alpha _S=`$const, then restoration of factorization is much slower, cf. our Fig. 10 and Fig. 9 in .
For an obvious reason that the soft-pomeron exchange is so predominant in real photon scattering, whereas the soft plus rightmost hard BFKL exchange is outstanding in virtual-virtual and real-virtual photon-photon scattering, it is ill advised to look at factorization ratio $`R_{\gamma ^{}\gamma ^{}}(W)`$ when one of the photons is quasireal, $`P^20`$. In this limit one would find strong departures of $`R_{\gamma ^{}\gamma ^{}}(W)`$ from unity. For precisely the same reason of predominance of soft-pomeron exchange in $`pp`$ scattering vs. nearly dominant rightmost hard BFKL pole exchange in DIS at small $`x`$ and 5-10$`\text{ }<Q^2\text{ }<`$ 100 GeV<sup>2</sup>, see , the naive factorization estimate
$$\sigma ^{\gamma ^{}\gamma ^{}}(W,Q^2,P^2)\frac{\sigma ^{\gamma ^{}p}(W,Q^2)\sigma ^{\gamma ^{}p}(W,P^2)}{\sigma ^{pp}(W)}$$
(22)
would not make much sense.
## 7 Conclusions
We explored the consequences for small-$`x`$ photon SFs $`F_{2\gamma }(x,Q^2)`$ and high-energy two-photon cross sections $`\sigma ^{\gamma ^{}\gamma ^{}}`$ and $`\sigma ^{\gamma \gamma }`$ from the color dipole BFKL-Regge factorization. Because of the nodal properties of eigen SFs of subleading hard BFKL exchanges the CD BFKL approach predicts uniquely that the vacuum exchange is strongly dominated by the combination of soft plus rightmost hard BFKL pole exchanges in a very broad range of photon virtualities $`Q^2,P^2`$ which includes much of the kinematical domain attainable at LEP200 and NLC. Starting with very restrictive perturbative two-gluon exchange as a boundary condition for BFKL evolution in the color dipole basis and having fixed the staring point of BFKL evolution in the early resulting CD BFKL-Regge phenomenology of the proton SF, we presented parameter-free predictions for the vacuum exchange contribution to the photon structure function which agree well with OPAL and L3 determinations. A good agreement is found between our predictions for the energy and photon virtuality dependence of the photon-photon cross section $`\sigma ^{\gamma ^{}\gamma ^{}}(W,Q^2,P^2)`$ and the recent data taken by the L3 Collaboration. We commented on the utility of Regge factorization tests of the CD BFKL-Regge expansion.
Acknowledgments: This work was partly supported by the grants INTAS-96-597 and INTAS-97-30494 and DFG 436RUS17/11/99.
## 8 Appendix
### 8.1 CD BFKL all flavor eigen-SF
In the early discussion of DIS off protons the results of numerical solutions of the CD BFKL equation for the all flavor ($`u+d+s+c`$) eigen-SF $`f_m(Q^2)`$ were parameterized as
$$f_0(Q^2)=a_0\frac{R_0^2Q^2}{1+R_0^2Q^2}\left[1+c_0\mathrm{log}(1+r_0^2Q^2)\right]^{\gamma _0},$$
(23)
$$f_m(Q^2)=a_mf_0(Q^2)\frac{1+R_0^2Q^2}{1+R_m^2Q^2}\underset{i=1}{\overset{m}{}}\left(1\frac{z}{z_m^{(i)}}\right),m1,$$
(24)
where $`\gamma _0=\frac{4}{3\mathrm{\Delta }_0}`$ and
$$z=\left[1+c_m\mathrm{log}(1+r_m^2Q^2)\right]^{\gamma _m}1,\gamma _m=\gamma _0\delta _m.$$
(25)
The parameters tuned to reproduce the numerical results for $`f_m(Q^2)`$ at $`Q^2\text{ }<10^5GeV^2`$ are listed in the Table 2.
Table 2. CD BFKL-Regge the all flavor SF parameters.
| $`m`$ | $`a_m`$ | $`c_m`$ | $`r_m^2,`$ $`\mathrm{GeV}^2`$ | $`R_m^2,`$ $`\mathrm{GeV}^2`$ | $`z_m^{(1)}`$ | $`z_m^{(2)}`$ | $`z_m^{(3)}`$ | $`\delta _m`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 0 | 0.0232 | 0.3261 | 1.1204 | 2.6018 | | | | 1. |
| 1 | 0.2788 | 0.1113 | 0.8755 | 3.4648 | 2.4773 | | | 1.0915 |
| 2 | 0.1953 | 0.0833 | 1.5682 | 3.4824 | 1.7706 | 12.991 | | 1.2450 |
| 3 | 1.4000 | 0.04119 | 3.9567 | 2.7706 | 0.23585 | 0.72853 | 1.13044 | 0.5007 |
| soft | 0.1077 | 0.0673 | 7.0332 | 6.6447 | | | | |
The soft component of the proton structure function as derived from $`\sigma _{\mathrm{soft}}(r)`$ taken from is parameterized as follows
$$f_{\mathrm{soft}}(Q^2)=\frac{a_{\mathrm{soft}}R_{\mathrm{soft}}^2Q^2}{1+R_{\mathrm{soft}}^2Q^2}\left[1+c_{\mathrm{soft}}\mathrm{log}(1+r_{\mathrm{soft}}^2Q^2)\right],$$
(26)
with parameters cited in the Table 2.
The cross section $`\sigma _{\mathrm{soft}}^{\gamma ^{}\gamma ^{}}(Q^2,P^2)`$ obtained by the continuation of the above
$$\sigma _{\mathrm{soft}}^{\gamma ^{}p}=\frac{4\pi ^2\alpha _{em}}{Q^2}f_{\mathrm{soft}}(Q^2)$$
into $`Q^2,P^2`$-plane reads
$$\sigma _{\mathrm{soft}}^{\gamma ^{}\gamma ^{}}(Q^2,P^2)=\frac{\sigma _{\mathrm{soft}}^{\gamma \gamma }}{1+R_{\mathrm{soft}}^2(Q^2+P^2)}\left[1+c_{\mathrm{soft}}\mathrm{log}\left(1+\frac{r_{\mathrm{soft}}^2Q^2}{1+r_{\mathrm{soft}}^2P^2}+\frac{r_{\mathrm{soft}}^2P^2}{1+r_{\mathrm{soft}}^2Q^2}\right)\right],$$
(27)
with parameters cited in the Table 2 and the on-shell cross section
$$\sigma _{\mathrm{soft}}^{\gamma \gamma }=\left[4\pi ^2\alpha _{em}a_{\mathrm{soft}}R_{\mathrm{soft}}^2\right]^2\frac{1}{\sigma _{\mathrm{soft}}^{pp}}.$$
(28)
### 8.2 CD BFKL charm eigen-SF
In practical evaluations one needs the charm eigen-SF, $`f_m^c(Q^2)`$. For the rightmost hard BFKL pole it is of the form
$$f_0^c(Q^2)=a_0\frac{R_0^2Q^2}{1+R_0^2Q^2}\left[1+c_0\mathrm{log}(1+r_0^2Q^2)\right]^{\gamma _0},$$
(29)
where $`\gamma _0=4/(3\mathrm{\Delta }_0)`$, while for the sub-leading hard BFKL poles
$$f_m^c(Q^2)=a_mf_0(Q^2)\frac{1+K_m^2Q^2}{1+R_m^2Q^2}\underset{i=1}{\overset{m_{max}}{}}\left(1\frac{z}{z_m^{(i)}}\right),m1,$$
(30)
where $`m_{max}=`$min$`\{m,2\}`$ and
$$z=\left[1+c_m\mathrm{log}(1+r_m^2Q^2)\right]^{\gamma _m}1,\gamma _m=\gamma _0\delta _m.$$
(31)
The parameters tuned to reproduce the numerical results for $`f_m^c(Q^2)`$ at $`Q^2\text{ }<10^5GeV^2`$ are listed in the Table 3.
The soft component of the charm SF is parameterized as
$$f_{\mathrm{soft}}^c(Q^2)=\frac{a_{\mathrm{soft}}R_{\mathrm{soft}}^2Q^2}{1+R_{\mathrm{soft}}^2Q^2}\left[1+c_{\mathrm{soft}}\mathrm{log}(1+r_{\mathrm{soft}}^2Q^2)\right],$$
(32)
with parameters cited in the Table 3.
Table 3. CD BFKL-Regge charm structure functions parameters. $`m`$ $`a_m`$ $`c_m`$ $`r_m^2,`$ $`R_m^2,`$ $`K_m^2,`$ $`z_m^{(1)}`$ $`z_m^{(2)}`$ $`\delta _m`$ $`\mathrm{GeV}^2`$ $`\mathrm{GeV}^2`$ $`\mathrm{GeV}^2`$ 0 0.02140 0.2619 0.3239 0.2846 1. 1 0.0782 0.03517 0.0793 0.2958 0.2846 0.2499 1.9249 2 0.00438 0.03625 0.0884 0.2896 0.2846 0.0175 3.447 1.7985 3 $`0.26313`$ 2.1431 $`3.742410^2`$ $`8.163910^2`$ 0.13087 158.52 559.50 0.62563 soft 0.01105 0.3044 0.09145 0.1303
Figure Captions
1. The normalized ratio of soft-to-rightmost-hard and subleading hard-to-rightmost hard expansion coefficients ($`m=1,2,3,\mathrm{soft}`$) $`r_m(Q^2)=\sigma _m^{\gamma ^{}\gamma ^{}}/\sigma _{\mathrm{vac}}^{\gamma ^{}\gamma ^{}}`$ of the BFKL-Regge expansion for $`\gamma ^{}\gamma ^{}`$ scattering at $`x=x_0`$.
2. Predictions from CD BFKL-Regge expansion for the effective intercept $`\mathrm{\Delta }_{eff}`$, eq.(15), for the diagonal case $`Q^2=P^2`$ and $`W=50,100,200`$ GeV.
3. Predictions from CD BFKL-Regge expansion for the vacuum exchange component of the virtual-virtial $`\gamma ^{}\gamma ^{}`$ cross section for the diagonal case of $`Q^2=P^2`$ and for cms collision energy W=50, 100 and 200 GeV (solid curves). The Leading Hard BFKL exchnage plus Soft-pomeron exchange Approximation (LHSA) is shown by the long dashed curve. The soft pomeron component of the cross section is shown separately by the dashed curve.
4. Predictions from CD BFKL-Regge expansion for the vacuum exchange component of the virtual-virtial $`\gamma ^{}\gamma ^{}`$ cross section for the diagonal case of $`Q^2=P^2`$ are confronted to the experimental data by the L3 Collaboration . The experimental data and theoretical curves are shown vs. the variable $`Y=\mathrm{log}(W^2/\sqrt{Q^2P^2})`$. The solid curve shows the result from the complete BFKL-Regge expansion, the soft-pomeron (the dashed curve) component included. the long dashed curve shows the rightmost hard BFKL (LH) plus soft-pomeron (S) approximation (LHSA).
5. Predictions from CD BFKL-Regge expansion for the vacuum exchange component of the virtual-virtial $`\gamma ^{}\gamma ^{}`$ cross section for the diagonal case of $`Q^2=P^2=17.9`$ GeV<sup>2</sup> are confronted to the experimental data by the OPAL Collaboration . The experimental data and theoretical curves are shown vs. the variable $`Y=\mathrm{log}(W^2/\sqrt{Q^2P^2})`$. The solid curve shows the result from the complete BFKL-Regge expansion, the soft-pomeron (the dashed curve) component included. The long dashed line corresponds to the rightmost hard BFKL (LH) plus soft-pomeron (S) approximation (LHSA).
6. Predictions from CD BFKL-Regge expansion for the photon SF. The solid curve shows the result from the complete BFKL-Regge expansion the soft-pomeron (the dashed curve) and valence (the dot-dashed curve) components included, the dotted curve shows the rightmost hard BFKL (LH) plus soft-pomeron (S) plus non-vacuum (NV) approximation (LHSNVA). The long dashed line corresponds to the LH plus S approximation (LHSA). Data points are from
7. Predictions from CD BFKL-Regge expansion for the vacuum exchange component of the the virtual-real $`\gamma ^{}\gamma `$ total cross section and for cms collision energy W=50, 100 and 200 GeV (solid curves). The result from the rightmost Hard BFKL (LH) plus Soft-pomeron (S) Approximation (LHSA) is shown by the long dashed curve. The soft-pomeron exchange component of the cross section is shown separately by the dashed curve.
8. Our predictions from the CD BFKL-Regge factorization for the single-vacuum exchange contribution to real-real $`\gamma \gamma `$ scattering are compared with the recent experimental data from the OPAL collaboration and .
9. Our evaluation of $`R_{\gamma \gamma }`$ for single-vacuum component of total cross sections.
10. The factorization cross section ratio $`R_{\gamma ^{}\gamma ^{}}(x)`$ at fixed $`x`$ and $`QP`$ as a function of $`Q/P`$ for $`x=10^2`$ (dotted line), $`x=10^3`$ (long-dashed) and $`x=10^4`$ (dashed). |
warning/0001/cond-mat0001156.html | ar5iv | text | # Thermodynamics and dielectric anomalies of DMAAS and DMAGaS crystals in the phase transitions region (Landau theory approach)
## 1 Introduction
Ferroelectric crystals (CH<sub>3</sub>)<sub>2</sub>NH<sub>2</sub>Al(SO<sub>4</sub>)<sub>2</sub> $``$ 6H<sub>2</sub>O (DMAAS) and (CH<sub>3</sub>)<sub>2</sub>NH<sub>2</sub>Ga(SO<sub>4</sub>)<sub>2</sub> $``$ 6H<sub>2</sub>O (DMAGaS) are intensively studied in recent years. Their interesting feature is possible existence of crystal in ferroelectric or antiferroelectric state depending on external conditions (e.g. temperature, hydrostatic pressure). There is a significant difference in thermodynamical behaviour of crystals despite on isomorphism of their structure. At ambient pressure DMAGaS crystal has three phases: paraelectric ($`T>T_c`$), ferroelectric ($`T_1<T<T_c`$) and antiferroelectric ($`T<T_1`$) with temperatures of phase transitions $`T_c=136`$ K (first order transition close to the tricritical point) and $`T_1=117`$ K (first order transition). There is only two phases in DMAAS crystal at ambient pressure: paraelectric ($`T>T_c`$) and ferroelectric ($`T<T_c`$) with $`T_c=155`$ K.
A set of structural , dilatometric, dielectric, pyroelectric and ultrasonic measurements is made for considered systems, what allows to establish their main dielectric, mechanical and dynamical characteristics (see below). At the same time these investigations are incomplete and of preliminary stage in many directions.
The nature of phase transitions in DMAAS and DMAGaS crystals was unclear up to recent time. During the last years conviction on important role of dimethyl ammonium (DMA) groups in phase transitions due to their orientational ordering-disordering is established (see, for example, ). In the microscopic approach based on the order-disorder model with account of different orientational states of DMA groups was proposed. In the framework of the model the phase transition to ferroelectric state has been described and conditions of realization of this transition as of the first or of the second order have been established. Order parameters of the system have been constructed. They are connected with differences of occupancies of four possible positions of nitrogen ions corresponding to different orientations of groups. As a result of symmetry analysis it has been established that components of the order parameters belonging to irreducible representation $`B_u`$ of point symmetry group $`2/m`$ of the high-temperature (paraelectric) phase describe ferroelectric ordering of DMA group along the ferroelectric axis OX (in crystallographic plane (ac)) and their antiferroelectric ordering along the OY axis (crystallographic axis b). The inverse ordering (antiferroelectric along OX and ferroelectric one along OY) corresponds to order parameter components belonging to irreducible representation $`A_u`$. Appearance of nonzero order parameters of $`B_u`$ type turns the system into ferroelectric state (point group $`m`$) while nonzero order parameters of $`A_u`$ type cause antiferroelectric state (point group 2).
Notwithstanding further perspectives of microscopic approach by means of the four-state order-disorder model, the more simple but more general thermodynamical description based on Landau expansion is of interest. One can construct corresponding Landau free energy and in standard way investigate possible phase transitions and obtain criteria of their realization with the use of data of the mentioned above symmetry analysis. This is a main goal of the present work. Results obtained in the framework of Landau expansion will be used for interpretation of the induced by the external pressure changes in the picture of phase transitions and for description of dielectric anomalies in the phase transition points of the investigated crystals.
## 2 Thermodynamics of phase transitions (Landau theory approach)
Let us make thermodynamical description of phase transition in DMAAS and DMAGaS crystals with the use of Landau expansion. We consider a simplified version when only one linear combination of the initial order parameters type is included for each of $`B_u`$ and $`A_u`$ irreducible representations. The combinations included are true order parameters: coefficients at their squared values tend to zero in the points of corresponding second order transitions.
Order parameters, which transform according to irreducible representations $`B_u`$ and $`A_u`$ of point symmetry group $`2/m`$ of high-symmetry phase, are denoted as $`\eta _b`$ and $`\eta _a`$ correspondingly. The first parameter $`\eta _b`$ describes polarization of ferroelectric type along the OX axis with simultaneous antiferroelectric type ordering along the OY axis; the second one corresponds to inverse orientation where antipolarization along OX is accompanied by polarization along OY.
We restrict ourself to the case of second order phase transition from the nonpolar high-temperature phase to ordered one. In this case Landau expansion of free energy can be limited by terms of the fourth order:
$$F=F_0+\frac{1}{2}a\eta _a^2+\frac{1}{2}b\eta _b^2+\frac{1}{4}c\eta _a^4+\frac{1}{4}d\eta _b^4+\frac{1}{2}f\eta _a^2\eta _b^2E_x\eta _bE_y\eta _a.$$
(1)
A linear dependence of coefficients $`a`$ and $`b`$ on temperature is assumed
$$a=a^{}(TT_c^{}),b=b^{}(TT_c),$$
(2)
where condition $`T_c>T_c^{}`$ is satisfied for normal state of the crystal what corresponds to the transition from the paraelectric phase (phase P) to the ferroelectric phase (phase F) as to the first one at lowering of the temperature.
Conditions of thermodynamical equilibrium correspond to the minimum of free energy and look like
$`{\displaystyle \frac{F}{\eta _a}}`$ $`=`$ $`\eta _a(a+c\eta _a^2+f\eta _b^2)E_y=0,`$
$`{\displaystyle \frac{F}{\eta _b}}`$ $`=`$ $`\eta _b(b+c\eta _b^2+f\eta _a^2)E_x=0.`$ (3)
At zero external fields there are following solutions
$$\eta _a=\eta _b=0$$
(4)
– paraphase (P-phase);
$`\eta _a=0,\eta _b0,`$
$`\eta _{b0}=\sqrt{b/d}=\sqrt{(b^{}/d)(T_cT)}`$ (5)
– ferroelectric phase (F-phase);
$`\eta _a0;\eta _b=0`$
$`\eta _{a0}=\sqrt{a/c}=\sqrt{a^{}/c(T_c^{}T)}`$ (6)
– antiferroelectric phase (AF-phase). <sup>1</sup><sup>1</sup>1We follow here to the terminology widely used in literature on the subject.
Corresponding expressions for free energy in these phases are as follows
$`F_{(P)}`$ $`=`$ $`F_0,`$
$`F_{(F)}`$ $`=`$ $`F_0{\displaystyle \frac{1}{4d}}b^2(TT_c)^2,`$ (7)
$`F_{(AF)}`$ $`=`$ $`F_0{\displaystyle \frac{1}{4c}}a^2(TT_c^{})^2.`$
The phase transition P$``$F which is of the second order in the used approximation takes place at temperature $`T_c`$. The phase transition F$``$AF which can take place at lower temperatures occurs at
$$F_{(F)}=F_{(AF)}.$$
(8)
The condition above determines the temperature of this first order phase transition:
$$T_1=\frac{\varkappa T_c^{}T_c}{\varkappa 1},$$
(9)
where
$$\varkappa =\frac{a^{}\sqrt{d}}{b^{}\sqrt{c}},\varkappa >1.$$
(10)
Nonequalities
$$0<T_1<T_c$$
(11)
define the region of temperature $`T_c^{}`$ values where the F-phase exists as an intermediate one:
$$\frac{1}{\varkappa }<\frac{T_c^{}}{T_c}<1.$$
(12)
These conditions are illustrated by the phase diagram in Fig. 1. In the case $`T_c^{}>T_c`$ a direct phase transition P$``$AF from the paraelectric phase to antiferroelectric one can take place.
Observed by experiment changes of temperatures of P$``$F and F$``$AF phase transitions and consecutive disappearance of the F-phase as the result of increasing of external hydrostatic pressure can be easily explained with the use of the obtained diagram. Under assumption that the influence of pressure leads mainly to shifts of temperatures $`T_c`$ and $`T_c^{}`$
$`T_c`$ $`=`$ $`T_{c0}+xp,`$
$`T_c^{}`$ $`=`$ $`T_{c0}^{}+x^{}p,`$ (13)
and the changes of other Landau expansion parameters are negligible, the following relation is obtained
$$T_1=T_1^0+\frac{\varkappa x^{}\varkappa }{\varkappa 1}p,$$
(14)
where
$$T_1^0=\frac{\varkappa T_{c0}^{}T_{c0}}{\varkappa 1}.$$
(15)
According to the data published in , $`dT_c/dpx=0.277`$ K/MPa; $`T_1/p=1.95`$ K/MPa and if one applies a linear approximation to the dependence of $`T_c^{}`$ on $`p`$ then $`dT_c^{}/dpx^{}=0.86`$ K/MPa.
The obtained relations are illustrated by the diagram shown in Fig. 2. This diagram qualitatively matches the experimental (T,p) diagram for DMAGaS crystal (at $`T_{c0}=136`$ K, $`T_{10}=116`$ K) . Obtained coordinates of triple point
$$T_3=\frac{x^{}T_{c0}xT_{c0}^{}}{x^{}x},P_s=\frac{T_{c0}T_{c0}^{}}{x^{}x},$$
(16)
where lines of phase transitions P$``$F, F$``$AF and P$``$AF come together are in good agreement with experimental ones ($`T_3^{exp}=140.3`$ C; $`P_3^{exp}=8.75`$ MPa). At $`pp_3`$ there take place a deviation of the theoretical prediction of temperature of the P$``$AF phase transition from experimental data. Unlike to relationship (13) experimental dependence is nonlinear at large pressures.
The pressure value
$$p^{}=\frac{\varkappa T_{c0}^{}T_{c0}}{x\varkappa x^{}}$$
(17)
(see Fig. 2) is an important characteristic of the model. At $`p^{}<0`$, what is realized at $`T_{c0}^{}/T_{c0}>1/\varkappa `$, AF-phase exists in the region of low temperatures at ambient pressure (this situation takes place for DMAGaS). At $`p^{}>0`$ (i.e. $`T_{c0}^{}/T_{c0}<1/\varkappa `$) and ambient pressure only P- and F-phases occur; this case can correspond to DMAAS crystal.
## 3 Dielectric susceptibility
The approach used in the previous section allows to derive expressions for components of dielectric susceptibility tensor in the vicinity of phase transition points and to describe their temperature dependencies in general. In the used approximation the components $`P_x`$ and $`P_y`$ of polarization vector are defined by parameters $`\eta _b`$ and $`\eta _a`$ correspondingly. Hence
$$\chi _{xx}=\frac{\eta _b}{E_x},\chi _{yy}=\frac{\eta _a}{E_y}$$
(18)
and proceeding from equations (3) one can obtain
$`\chi _{xx}`$ $`=`$ $`{\displaystyle \frac{1}{D}}(a+3c\eta _a^2+f\eta _b^2),`$
$`\chi _{yy}`$ $`=`$ $`{\displaystyle \frac{1}{D}}(b+3d\eta _b^2+f\eta _a^2),`$ (19)
where
$$D=(a+3c\eta _a^2+f\eta _b^2)(b+3d\eta _b^2+f\eta _a^2)4f^2\eta _a^2\eta _b^2.$$
(20)
The following particular cases follow from expression (19):
1. Paraphase (P):
$$\chi _{xx}=\frac{1}{b}=\frac{1}{b^{}(TT_c)},\chi _{yy}=\frac{1}{a}=\frac{1}{a^{}(TT_c^{})}.$$
(21)
2. Ferroelectric phase (F):
$$\chi _{xx}=\frac{1}{2b}=\frac{1}{2b^{}(T_cT)},\chi _{yy}=\frac{1}{(\xi 1)a^{}(T^{}T)},$$
(22)
here the notations are used:
$$T^{}=T_c+\frac{T_cT_c^{}}{\xi 1},\xi =\frac{fb^{}}{da^{}}(\xi >1).$$
(23)
In this case susceptibility $`\chi _{yy}`$ can be also expressed in the form
$$\chi _{yy}=[a+f\eta _{b0}^2]^1,$$
(24)
where $`\eta _{b0}`$ is a spontaneous value of order parameter (polarization $`P_s`$) in the ferroelectric phase.
3. Antiferroelectric phase (AF):
$$\chi _{xx}=\frac{1}{(\varkappa ^2\xi 1)b^{}(T^{}T)},\chi _{yy}=\frac{1}{2a}=\frac{1}{2a^{}(T_c^{}T)},$$
(25)
where the temperature
$$T^{}=T_c+\frac{T_cT_c^{}}{11/\varkappa ^2\xi },$$
(26)
is introduced such that $`T^{}>T^{}>T_c`$. A similar to the previous one expression
$$\chi _{xx}=[b+f\eta _{a0}^2]^1,$$
(27)
relating the temperature dependence of longitudinal susceptibility in AF phase with the equilibrium value of the order parameter (polarization in one of sublattices) takes place.
The temperature behaviour of dielectric susceptibility components and their anomalies in the phase transition points are illustrated in Fig. 4 and 4 as temperature dependencies of inverse susceptibilities $`\chi _{\alpha \alpha }^1`$.
The inverse susceptibility $`\chi _{xx}^1`$ is equal to zero at the temperature $`T_c`$. Its linear dependence on temperature in the vicinity of this point has an inclination $`b^{}`$ at $`T>T_c`$ and $`2b^{}`$ in the ferroelectric phase (Fig. 4). This typical behaviour for second order phase transition changes if the phase transition P$``$F is of the first order. Such a situation takes place in the DMAGaS crystal where the first order phase transition close to the tricritical point is observed. Then the susceptibility $`\chi _{xx}^1`$ remains nonzero at $`T_c`$ and has a small jump (according to data , $`T_cT_01.2`$ K, where $`T_0`$ is the temperature at which $`\chi _{xx}^10`$; $`\chi _{xx}^1(T=T_c)610^4`$). Mentioned changes are relevant only to a small vicinity of $`T_c`$; in a large temperature scale dependence $`\chi _{xx}^1(T)`$ in para- and ferroelectric phase is almost the same as for the second order transition. The phase transition F$``$AF is a well pronounced first order phase transition accompanied by jump of the $`\chi _{xx}^1`$ function. The continuation of the straight line describing the temperature dependence of $`\chi _{xx}^1`$ in the AF phase passes the point $`T^{}`$ (see Fig. 3). $`\chi _{xx}^1`$ has the following values at the ends of its jump
$`\chi _{xx}^1|_1`$ $`=`$ $`2b^{}{\displaystyle \frac{\varkappa }{\varkappa 1}}(T_cT_c^{}),`$
$`\chi _{xx}^1|_2`$ $`=`$ $`\left[{\displaystyle \frac{\varkappa }{\varkappa 1}}(\varkappa ^2\xi 1)+\varkappa ^2\xi \right]b^{}(T_cT_c^{}).`$ (28)
Value of susceptibility jump $`\mathrm{\Delta }\chi _{xx}^1=\chi _{xx}^1|_2\chi _{xx}^1|_1`$ can be positive or negative depending on values of theory parameters.
Temperature behaviour of the inverse susceptibility $`\chi _{yy}^1`$ is essentially different. In the point of the second order phase transition P$``$F it remains nonzero with value
$$\chi _{yy}^1(T_c)=a^{}(T_cT_c^{})$$
(29)
Its continuation to lower temperatures goes to zero at $`TT_c^{}`$. The continuation of the line of the inverse susceptibility in the antiferroelectric phase $`\chi _{yy}^1(T)=2a^{}(T_c^{}T)`$ also goes across this point. In the ferroelectric phase region the function $`\chi _{yy}^1(T)`$ is linear with the continuation passing the point $`T^{}`$. At the F$``$AF phase transition this function has a jump between points
$`\chi _{yy}^1|_3`$ $`=`$ $`{\displaystyle \frac{2a^{}}{\varkappa 1}}(T_cT_c^{}),`$
$`\chi _{yy}^1|_4`$ $`=`$ $`{\displaystyle \frac{\varkappa \xi 1}{\varkappa 1}}a^{}(T_cT_c^{}).`$ (30)
Similarly to the case of the function $`\chi _{xx}^1`$ the jump can have positive or negative value.
## 4 Discussion
Proceeding from obtained in the previous section formulae one can try to interpret available data on the temperature dependence of dielectric susceptibility components of DMAAS and DMAGaS crystals. The majority of performed measurements is devoted to the longitudinal dielectric permittivity $`\epsilon _x`$ (or its real part $`\epsilon _x^{}`$ for low frequency alternating current measurements) mainly in the region of the high-temperature phase transition for DMAGaS and the corresponding phase transition in DMAAS. Such data are reported in works (DMAGaS) and (DMAAS); only in paper the temperature behaviour of all permittivity components ($`\epsilon _a^{}`$, $`\epsilon _b^{}`$, $`\epsilon _c^{}`$) for DMAAS crystal in the wide range of temperatures (from $``$90 K to $``$280 K) was measured. In some papers dependence of spontaneous polarization on temperature in the ferroelectric phase was investigated and coercivity fields were measured (the value of $`P_s`$ in the state close to saturation is about 1.4–1.9 C/m<sup>2</sup> for DMAAS and 0.9–2.0 C/m<sup>2</sup> for DMAGaS). Particular investigation of the $`T_c`$ point vicinity in DMAGaS devoted to influence of the external electric field on the first order phase transition point and the difference $`T_cT_0`$ is made in . On the basis of available experimental data Curie-Weiss constant (from the paraphase side) is estimated as 2700–3060 K for DMAGaS crystal and 2700–3000 K for DMAAS crystal. The phase transition to the ferroelectric phase in DMAGaS crystal is of the first order and close to the tricritical point; this fact however does not affect the behaviour of $`\chi _{xx}`$ and $`\chi _{yy}`$ far from the $`T_c`$ point.
The mentioned experimental data are incomplete, hence only partial comparison with results of thermodynamical description is possible. For example one can obtain values of the temperature $`T_c^{}`$, parameters $`b^{}`$ and $`\varkappa `$ for the DMAGaS crystal $`T_c^{}=125`$ K, $`\varkappa =2.22`$, $`b^{}=0.3310^3`$ K<sup>-1</sup> with use of above mentioned data on the influence of external hydrostatic pressure on phase transitions in DMAGaS crystal and results of measurements of dielectric characteristics.
More comprehensive and selfconsistent evaluation of temperatures $`T_c^{}`$, $`T^{}`$ and $`T^{}`$ as well as Landau expansion parameters (or parameters $`a^{}`$, $`b^{}`$, $`\varkappa `$, $`\xi `$ and $`f`$) by means of presented in this section relationships become possible after goal-oriented investigations of temperature dependencies of $`\chi _{xx}^1`$ and $`\chi _{yy}^1`$ in a wide temperature interval including regions of existence of all phases for DMAGaS and DMAAS. Proceeding from obtained results will be possible to ascertain suitability of the simple thermodynamical description where Landau expansion is limited to only one order parameter for each of $`B_u`$ and $`A_u`$ representations. Such a description is obviously much simplified comparatively to results of the microscopic approach based on the four-state model of order-disorder type . Investigation of DMA group ordering in the configurational space of four orientational states needs two-component order parameters $`\eta _b^\alpha `$ ($`B_u`$) and $`\eta _a^\alpha `$ ($`A_u`$), $`\alpha =1,2`$. This fact could complicate temperature dependencies of dielectric characteristics of the model even for thermodynamical description in the framework of Landau expansion.
Furthermore, considered here expression for Landau expansion of free energy (1) includes terms up to the fourth order. A consistent description of the first order phase transition P$``$F and related dielectric anomalies demands the inclusion of the sixth order terms into expansion. Such a generalization is necessary for comprehensive description of experimental data and can be performed relatively easy. |
warning/0001/hep-ph0001105.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Recent polarized deep-inelastic experiments have yielded valuable information about the quark helicity distributions in the proton, but have left some crucial questions unanswered. One of the key unknowns is the size of the polarized gluon distribution. In light of the data, there have been numerous constituent models reflecting contributions of quarks and gluons to the proton spin. Most of these models reproduce existing data fairly well, but vary as to their physical assumptions and the overall functional form. There is a wide variety of results for the polarized strange sea and the polarized gluons in the proton, for example. Since the constituent distributions are used for predicting the spin observables, their measurement can be used to distinguish between the physical assumptions of these models. Hence, further experimental information is necessary in order to understand the nature of these constituents, in particular, the polarized glue.
The Relativistic Heavy Ion Collider (RHIC) at Brookhaven (BNL) is well suited for the polarized beam experiments which can help to reveal the nature of the proton constituents’ spin properties. The two major detectors, STAR and PHENIX, cover a similar kinematic region, but STAR has wider angular coverage, while PHENIX has finer granularity for more precise measurements.
This paper discusses the optimal observables for these detectors to determine $`\mathrm{\Delta }G`$. Section II will provide a brief overview of some theoretical models of the spin constituents of the proton. The physical assumptions and use of present polarized DIS data will be discussed. In section III, key experimental predictions for the determination of $`\mathrm{\Delta }G`$ will be compared. Finally, suggestions are made for the best possibilities to experimentally determine the nature of $`\mathrm{\Delta }G`$ in the accessible kinematic regions of the two detectors.
## 2 Models of the Polarized Distributions
In recent work we constructed three sets of polarized parton distributions for the valence and sea quarks, based upon three models for the polarized gluons. Each gluon model has a different physical basis, but the quark distributions comform to a set of reasonable theoretical assumptions. All of the distributions are generated at $`Q_0^2=1.0`$ GeV<sup>2</sup> and evolved in NLO, entirely in $`x`$-space, up to the $`Q^2`$ values necessary to predict the spin observables. All three sets are in good agreement with data.
The polarized valence distributions are generated from a modified SU(6) distribution with a spin dilution factor to control their small-$`x`$ behavior. The polarized valence distributions are written in terms of the unpolarized CTEQ distributions as:
$`\mathrm{\Delta }u_v`$ $`=`$ $`\mathrm{cos}\theta _D\left[u_v{\displaystyle \frac{2}{3}}d_v\right]`$
$`\mathrm{\Delta }d_v`$ $`=`$ $`\mathrm{cos}\theta _D\left[{\displaystyle \frac{1}{3}}d_v\right],`$ (1)
where the spin dilution factor is: $`\mathrm{cos}\theta _D=[1+R_0(1x)^2/\sqrt{(}x)]^1`$. The free parameter, $`R_0`$, is fixed by applying the Bjorken Sum Rule. In the $`Q^2`$ region of the present PDIS data, we find that $`R_0\frac{2}{3}\alpha _s`$. This parametrization gives the appropriate behavior at both small and large $`x`$.
For the polarized sea, we assume a broken SU(3) model, to account for mass effects in polarizing the strange sea. Our models separate all light flavors in the valence and sea. Charm is included via the evolution equations ($`N_f`$), at the appropriate $`Q^2`$ of charm production, to avoid any non-empirical assumptions about its size. The small-$`x`$ behavior is of the Regge type, and the large-$`x`$ behavior is compatible with the appropriate counting rules.
The polarized sea distributions are extracted from both the unpolarized CTEQ sea distributions and polarized deep-inelastic-scattering data. We assume a model of the sea which obtains its polarization from gluon Bremsstrahlung, so the net polarization of the quarks is dependent upon their density in hadrons in LO. We therefore assume that these densities are directly proportional and the flavor dependent sea distributions have the form
$$\mathrm{\Delta }q_i=\eta _i(x)xq_i(x).$$
(2)
The $`\eta _i`$ are determined from the integrated distributions and sum rules used to analyze polarized deep-inelastic-scattering data. We have chosen the functional form of $`\eta _i(x)`$ to have a reasonable Regge type behavior, to be consistent with positivity constraints and to yield the proper normalization (indicated here by the factor $`\eta `$) in reflecting the relative spin that each flavor contributes to that of the proton. Here, $`\eta _i(x)`$ may be interpreted as a modification of $`\mathrm{\Delta }q(x)`$ due to unknown soft effects at small-$`x`$. The normalization requires that $`_0^1\mathrm{\Delta }q_i(x)𝑑x=_0^1\eta _i(x)xq_i(x)𝑑x=\eta _0^1xq_i(x)𝑑x.`$ The overall parametrization for each of the polarized sea flavors, including the $`\eta (x)`$ functions, the anomaly terms and the up-down unpolarized asymmetry term can be written (with the CTEQ basis) in the form:
$$\mathrm{\Delta }q_i(x)=Ax^{0.143}(1x)^{8.041}(1B\sqrt{x})\left[1+6.112x+P(x)\right].$$
(3)
The values for the variables in this equation are given for each flavor and each gluon model in Table I. Note that $`\delta (x)0.278x^{0.644}`$ is due to the asymmetry of the unpolarized up and down anti-quarks.
We consider three distinct models for the polarized gluons, which have a moderate size. There exists no empirical evidence that the polarized gluon distribution is large at the relatively small $`Q^2`$ values of the data. Data from the Fermilab E704 experiment indicate that it is likely small at these $`Q^2`$ values. In addition, a theoretical model of $`\mathrm{\Delta }G`$ based on counting rules, implies that $`\mathrm{\Delta }G\frac{1}{2}`$. Our choice of models effectively includes two separate factorization schemes, Gauge Invariant (GI or $`\overline{MS}`$) and Chiral Invariant (CI or Adler-Bardeen), which can be used to represent the polarized sea distributions.
The first set of $`\mathrm{\Delta }q_i(x)`$ functions, quoted in Table I assumes a moderately polarized glue, normalized to $`\frac{1}{2}`$ using the CTEQ unpolarized gluons. The second polarized gluon model assumes $`\mathrm{\Delta }G=0`$. This is equivalent to writing the quark distributions in the gauge-invariant scheme, since the anomaly term vanishes. The third model is motivated by an instanton-induced polarized gluon distribution, which gives a negatively polarized glue at small-$`x`$. The three polarized gluon distributions are written as
$`\mathrm{\Delta }G(x)`$ $`=`$ $`xG(x)`$
$`\mathrm{\Delta }G(x)`$ $`=`$ $`0`$ (4)
$`\mathrm{\Delta }G(x)`$ $`=`$ $`7(1x)^7\left[1+0.474\mathrm{ln}(x)\right].`$
Our overall distributions agree very well with existing data.
The models of Gehrmann and Stirling are based upon an SU(3) symmetric sea with the small-$`x`$ behavior of the sea and glue assumed equal. All of the difference between the Ellis-Jaffe sum rule and the data are attributed to $`\mathrm{\Delta }G`$, with $`\mathrm{\Delta }s0`$. This differs considerably with our models. They also have three different polarized gluon models, which run the range of hard and soft gluons. The GSA $`\mathrm{\Delta }G`$ is much larger than that of our first model. The other two fall within the range of ours, and are therefore not discussed here. Our motive is to present a wide range of models for $`\mathrm{\Delta }G`$ so that the experiments can be used to distinguish the size of $`\mathrm{\Delta }G`$ as opposed to the overall accuracy of a particular model. The differences in the predictions of the GSA and the three GGR models (A, B, and C) are thus due to the $`x`$-behavior of the quark distributions and the relative sizes of $`\mathrm{\Delta }G`$.
## 3 A Comparison of Experimental Predictions
The polarization experiments planned for RHIC show great potential for extracting information on polarized distributions, especially $`\mathrm{\Delta }G`$. With polarized beams of $`70\%`$ polarization and luminosity of $`2\times 10^{32}`$ /cm<sup>-2</sup>/sec<sup>-1</sup>, both prompt-$`\gamma `$ production and jet production can be done in a kinematic region where determination of $`\mathrm{\Delta }G`$ is possible. If the planned integrated luminosity of 320 pb<sup>-1</sup> at $`\sqrt{s}=200`$ GeV is attained, the resulting data should be good enough to distinguish among many of the polarized gluon models which have been proposed.
The STAR detector will have a wide angular range to cover a large rapidity, especially for jet production. The PHENIX detector has a narrow rapidity, but finer granularity, and is well suited for measuring high $`p_T`$ prompt photons. Both are planned to have an accessible $`p_T`$ range of $`10p_T30`$ GeV at $`\sqrt{s}=200`$ GeV. The predictions shown here cover this kinematic range. Possible experiments which would provide a measure of $`\mathrm{\Delta }G`$ include:
* one and two jet production in $`ep`$ and $`pp`$ collisions,
* prompt photon production
* charm production
+ polarized heavy quark production: these asymmetries are between one and four percent for $`\sqrt{s}=200`$ GeV with $`p_T5`$ GeV
+ $`\chi _{0c}`$ and $`\chi _{1c}`$ production: for very large models of $`\mathrm{\Delta }G`$, these asymmetries range from eight down to one percent for $`\sqrt{s}=500`$ GeV and $`2p_T30`$ GeV.
+ $`\chi _{2c}`$ production: even for the larger models of $`\mathrm{\Delta }G`$, these asymmetries range from zero to four percent at $`\sqrt{s}=200`$ GeV and $`2p_T8`$ GeV.
* $`J/\psi `$ production these asymmetries are typically from two to six percent at $`\sqrt{s}=200`$ GeV and $`2p_T10`$ GeV.
* pion production .
All but the first two of these yield small asymmetries, even for moderately sizable gluon polarizations. Thus, we feel that jet production and prompt-$`\gamma `$ production are the best choices for extracting information about $`\mathrm{\Delta }G`$. The asymmetries are not large everywhere, but there are kinematic regions where models based upon the different sizes of $`\mathrm{\Delta }G`$ can be distinguished. We have calculated prompt-$`\gamma `$ production in LO and NLO, and jet production at LO. The figures compare the predictions of our three models with the Gehrmann and Stirling A model.
Figure 1 compares the predictions of the GSA (large $`\mathrm{\Delta }G`$) with the three GGR gluon models for the prompt photon asymmetry, $`A_{LL}^\gamma `$. The error bars shown are those expected for $`A_{LL}^\gamma `$ at the PHENIX detector in these kinematic regions. Figure 2 shows the same asymmetry predictions at fixed $`p_T`$ of $`15`$ GeV as a function of rapidity. The error bar shown is for PHENIX, which operates at essentially this point of rapidity. STAR has a much wider rapidity coverage, $`1y2`$. Figure 3 shows the comparison of the four models for jet production at $`\sqrt{S}=200`$ GeV for zero rapidity. The asymmetries for jet production at $`\sqrt{S}=500`$ GeV are much smaller than for the lower energy.
## 4 Extracting $`\mathrm{\Delta }G`$
The various polarized gluon models have different physical bases and provide a reasonable range of possibilities, which can be narrowed down by future experiments at RHIC. If the polarized gluon distribution is moderately positive at $`Q^2=1`$ GeV<sup>2</sup>, the asymmetry for prompt photon production is among the best candidates for determining the size of $`\mathrm{\Delta }G`$. Jet production will be a close contender for distinguishing among the various predictions for $`\mathrm{\Delta }G`$. Much depends upon the relative errors in the applicable kinematic regions of STAR and PHENIX.
According to the projected uncertainties for STAR and PHENIX, the most favorable region to study prompt photon production is for $`15p_T25`$ GeV at $`\sqrt{s}=200`$ GeV. Although the asymmetries are closer together here, the favorable small uncertainties should be able to separate the large and small models for $`\mathrm{\Delta }G`$ (Fig. 1). Also, at $`p_T=15`$ GeV, the large rapidity region $`1y2`$ is a favorable place for STAR to measure the asymmetry due to the separation of predictions in the models (Fig. 2). Since PHENIX is designed for the small rapidity region, the measurement of $`A_{LL}^\gamma `$ with the better uncertainties is also promising, and will provide a good cross check of the results obtained by STAR.
Jet production is also a good candidate for determination of whether $`\mathrm{\Delta }G`$ is large or small. Since the asymmetries are fairly close together in the $`p_T`$ region between $`15`$ and $`30`$ GeV, measurements of this asymmetry require the larger values of $`p_T`$ to distinguish the relative size of $`\mathrm{\Delta }G`$ extracted from these predictions (Fig. 3).
Prompt photon and jet production, measured by STAR and PHENIX in these kinematic regions, appear to be the best candidates for providing an indication of the nature of $`\mathrm{\Delta }G`$.
## Acknowledgements
These results are based upon work done with L.E. Gordon (Jefferson Lab) and M. Goshtasbpour (Shahid Beheshti Univ, Tehran). |
warning/0001/cond-mat0001428.html | ar5iv | text | # 𝑂(4)-Invariant Formulation of the Nodal Liquid
\[
## Abstract
We consider the $`O(4)`$ symmetric point in the phase diagram of an electron system in which there is a transition between $`d_{x^2y^2}`$ density-wave order and $`d_{x^2y^2}`$ superconductivity. If the pseudospin $`SU(2)O(4)`$ symmetry is disordered by quantum fluctuations, the Nodal Liquid can result. In this context, we (1) construct a pseudospin $`\sigma `$-model; (2) discuss its topological excitations; (3) point out the possibility of a pseudospin-Peierls state and (4) propose a phase diagram for the underdoped cuprate superconductors.
\]
Introduction. Competing interactions and fluctuations have led to a cornucopia of interesting phenomena in the cuprate superconductors. Unfortunately, these phenomena have not led to the unambiguous determination of the phase diagram of these materials, possibly because some of the phases realized in these materials are characterized by particularly subtle forms of order. This dilemma is rather acute on the underdoped side of the phase diagram, where it is still not clear if the pseudogap can be ascribed to a new phase of matter, a nearby critical point, or a crossover. Since a better understanding of proposed exotic phases and the transitions between them may mitigate this difficulty, we study the transition between the $`d_{x^2y^2}`$ superconducting state of the cuprates and a putative $`d_{x^2y^2}`$ density-wave state (also known as the staggered flux state ; see ). We ask if the pseudo-gap – which appears to have $`d_{x^2y^2}`$ symmetry – could be due to the proximity of the experimental system to this transition. The resulting phase diagram automatically includes the Nodal Liquid state , a state with spin-charge separation. We discuss the possible relevance of this theoretical cuprate phase diagram to the experimental one.
$`O(4)`$ Formulation of $`d_{x^2y^2}`$ Ordered States at Half-Filling. In , we adapted Yang’s pseudospin $`SU(2)`$ symmetry to a critical point between a $`d_{x^2y^2}`$ density-wave state and a $`d_{x^2y^2}`$ superconductor. The original pseudospin $`SU(2)`$ was germane to the transition between a CDW and an $`s`$-wave superconductor; Zhang’s closely related $`SO(5)`$, to the transition between an antiferromagnet and a $`d`$-wave superconductor.
We first consider a transition at half-filling between a singlet commensurate $`d_{x^2y^2}`$ density-wave and a $`d_{x^2y^2}`$ superconductor. We combine the order parameters into
$`\mathrm{\Phi }_{\underset{¯}{i}}(q)f(k)=\left(\begin{array}{c}\sqrt{2}\mathrm{Re}\left\{\psi _{}^{}(k+\frac{q}{2})\psi _{}^{}(k+\frac{q}{2})\right\}\\ \sqrt{2}\mathrm{Im}\left\{\psi _{}^{}(k+\frac{q}{2})\psi _{}^{}(k+\frac{q}{2})\right\}\\ i\psi ^\alpha (k+Q+\frac{q}{2})\psi _\alpha (k\frac{q}{2})\end{array}\right)`$ (4)
where $`f(k)=\mathrm{cos}k_xa\mathrm{cos}k_ya`$. Following Yang , we introduce the pseudospin $`SU(2)`$ generators $`O^3,O^+,O^{}=(O^+)^{}`$
$`O^3`$ $`=`$ $`{\displaystyle _{\mathrm{R}.\mathrm{B}.\mathrm{Z}.}}{\displaystyle \frac{d^2k}{(2\pi )^2}}\left(\psi ^\alpha (k)\psi _\alpha (k)+kk+Q\right)`$ (5)
$`O^+`$ $`=`$ $`{\displaystyle _{\mathrm{R}.\mathrm{B}.\mathrm{Z}.}}{\displaystyle \frac{d^2k}{(2\pi )^2}}i\psi _{}^{}(k)\psi _{}^{}(k+Q)`$ (6)
The order parameters form a triplet under this $`SU(2)`$<sup>*</sup><sup>*</sup>*We will use underlined lowercase Roman letters such as $`\underset{¯}{i}=\underset{¯}{1},\underset{¯}{2},\underset{¯}{3}`$ to denote pseudospin triplet indices and uppercase Roman letters to denote peudospin doublet indices $`A=1,2`$. Lowercase Roman indices $`a=1,2,3`$ will be vector indices (i.e. real spin triplet indices) and Greek letters $`\alpha =1,2`$ will be used for real spin $`SU(2)`$ spinor indices. Pauli matrices $`\tau ^{\underset{¯}{i}}`$ will be used for pseudospin, while $`\sigma ^a`$ will be reserved for spin.. The integrals are over the reduced Brillouin zone.
There is a small but important difference between our pseudospin $`SU(2)`$ and Yang’s : the factors of $`i`$ in $`O^\pm `$. They are necessary since a commensurate $`d_{x^2y^2}`$ density-wave breaks $`T`$, while a superconductor does not; hence, our pseudospin $`SU(2)`$ does not commute with $`T`$. Pseudospin $`SU(2)`$, spin $`SU(2)`$, and time-reversal combine to form the symmetry group $`O(4)`$.
The electron fields transform as a doublet under both $`SU(2)`$s. We will group them into $`\mathrm{\Psi }_{A\alpha }`$:
$`\left(\begin{array}{c}\mathrm{\Psi }_{1\alpha }\\ \mathrm{\Psi }_{2\alpha }\end{array}\right)=\left(\begin{array}{c}\psi _\alpha (k)\\ iϵ_{\alpha \beta }\psi ^\beta (k+Q)\end{array}\right)`$ (11)
Near the transition between a $`d_{x^2y^2}`$ density-wave and a $`d_{x^2y^2}`$ superconductor, we can focus on the low-energy degrees of freedom: the order parameters and the nodal fermionic excitations. We can write down an $`O(4)`$-invariant action for this transition:
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle 𝑑\tau \frac{d^2k}{(2\pi )^2}\mathrm{\Psi }^{A\alpha ^{}}\left(_\tau ϵ(k)\right)\mathrm{\Psi }_{A\alpha }}+`$ (16)
$`ig{\displaystyle }d\tau {\displaystyle \frac{d^2k}{(2\pi )^2}}{\displaystyle \frac{d^2q}{(2\pi )^2}}\mathrm{\Phi }_{\underset{¯}{i}}(q)f(k)\times `$
$`[ϵ^{\alpha \beta }\mathrm{\Psi }_{C\alpha }(k+{\displaystyle \frac{q}{2}})ϵ^{CA}\tau _A^{\underset{¯}{i}B}\mathrm{\Psi }_{B\beta }(k+{\displaystyle \frac{q}{2}})+`$
$`ϵ_{\alpha \beta }\mathrm{\Psi }^{A\alpha }(k+{\displaystyle \frac{q}{2}})\tau _A^{\underset{¯}{i}B}ϵ^{BC}\mathrm{\Psi }^{B\beta }(k+{\displaystyle \frac{q}{2}})]`$
$`+{\displaystyle 𝑑\tau d^2x\left(\left(_\mu \mathrm{\Phi }_{\underset{¯}{i}}\right)^2+\frac{1}{2}r\mathrm{\Phi }_{\underset{¯}{i}}\mathrm{\Phi }_{\underset{¯}{i}}+\frac{1}{4!}u\left(\mathrm{\Phi }_{\underset{¯}{i}}\mathrm{\Phi }_{\underset{¯}{i}}\right)^2\right)}`$
‘Microscopic’ models with this symmetry were constructed in . In this $`O(4)`$-symmetric action, we have, by a rescaling, set the $`\mathrm{\Phi }_{\underset{¯}{i}}`$ velocities, $`v_{\underset{¯}{i}}`$, and stiffnesses, $`\rho _{\underset{¯}{i}}`$, to $`1`$. This cannot be done in the asymmetric case, $`\rho _{\underset{¯}{1}}=\rho _{\underset{¯}{2}}\rho _s\rho _{DW}\rho _{\underset{¯}{3}}`$. In general, symmetry-breaking terms will be present, but they can scale to zero at a critical point, thereby dynamically restoring the symmetry, as we discuss later. Hence, we focus on the symmetric case.
When $`\mathrm{\Phi }_{\underset{¯}{i}}`$ is ordered, the fermionic spectrum is $`E(k)=\sqrt{ϵ^2(k)+g^2\mathrm{\Phi }_{\underset{¯}{i}}\mathrm{\Phi }_{\underset{¯}{i}}f^2(k)}`$. In the following, we ignore the fermionic excitations which are not associated with the nodes of the $`d_{x^2y^2}`$ order parameter. We linearize $`ϵ(k)`$ about the Fermi surface and $`f(k)`$ about the nodes. If we rotate our axes so that the $`k_x`$ axis is perpendicular to the Fermi surface at one antipodal pair of nodes, then we can write $`ϵ(k)v_Fk_x`$ and $`g\left|\mathrm{\Phi }_{\underset{¯}{i}}\right|f(k)v_\mathrm{\Delta }k_y`$. As in , we will have to introduce an additional index $`a=1,2`$ for the two sets of antipodal nodes which differ by the replacement $`k_xk_y`$. In order to avoid unnecessary clutter, this index will be suppressed.
It is convenient to adopt a non-linear $`\sigma `$-model approach and assume that the magnitude of $`\mathrm{\Phi }_{\underset{¯}{i}}`$ is fixed, $`\mathrm{\Phi }_{\underset{¯}{i}}^2=a^2`$. Following , we employ a $`CP^1`$ representation of the non-linear $`\sigma `$ model:
$$\mathrm{\Phi }^{\underset{¯}{i}}=z^A\tau _A^{\underset{¯}{i}B}z_B$$
(17)
with $`\left|z_1\right|^2+\left|z_2\right|^2=a^2`$ and rotate the pseudospins of the fermions to the local direction of the order parameter:
$$\mathrm{\Psi }_A=U_A^B\chi _B$$
(18)
where
$`U={\displaystyle \frac{1}{a}}\left(\begin{array}{cc}z_1& z_2^{}\\ z_2& z_1^{}\end{array}\right)`$ (21)
The latter change of variables is a direct $`SU(2)`$ analogue of the original $`U(1)`$ Nodal Liquid construction . As in that case, it is is double-valued, so we must introduce a Chern-Simons term as in which couples the $`\chi `$’s to the topological current $`j_\mu =ϵ_{\mu \nu \lambda }ϵ_{\underset{¯}{i}\underset{¯}{j}\underset{¯}{k}}\mathrm{\Phi }^{\underset{¯}{i}}^\nu \mathrm{\Phi }^{\underset{¯}{j}}^\lambda \mathrm{\Phi }^{\underset{¯}{k}}`$. This term is only important at the phase transitions since the topological current vanishes in the ordered phases since the pseudospins are aligned and in the disordered phases since it is odd under the $`Z_2`$ symmetry $`\mathrm{\Phi }^{\underset{¯}{3}}\mathrm{\Phi }^{\underset{¯}{3}}`$. We suppress this term below.
In terms of $`z_A`$, $`\chi _A`$, the effective action takes the form:
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle 𝑑\tau d^2x\chi ^{A\alpha }\left(_\tau +\alpha _\tau \tau ^{\underset{¯}{3}}v_Fi_xv_F\alpha _x\tau ^{\underset{¯}{3}}\right)\chi _{A\alpha }}`$ (27)
$`+i{\displaystyle }d\tau {\displaystyle \frac{d^2k}{(2\pi )^2}}[ϵ^{\alpha \beta }\chi _{C\alpha }ϵ^{CA}\tau _A^{\underset{¯}{3}B}v_\mathrm{\Delta }i_y\chi _{B\beta }+`$
$`ϵ_{\alpha \beta }\chi ^{A\alpha }\tau _A^{\underset{¯}{3}B}ϵ^{BC}v_\mathrm{\Delta }i_y\chi ^{B\beta }]`$
$`+{\displaystyle }d\tau {\displaystyle \frac{d^2k}{(2\pi )^2}}\chi ^\alpha (U^{}(_\tau A_\tau \tau ^{\underset{¯}{3}})U\alpha _\tau \tau ^{\underset{¯}{3}}`$
$`v_FU^{}(i_xA_x\tau ^{\underset{¯}{3}})U+v_F\alpha _x\tau ^{\underset{¯}{3}})\chi _\alpha `$
$`+{\displaystyle 𝑑\tau d^2x\left(\left|\left(i_\mu \alpha _\mu A_\mu \tau ^{\underset{¯}{3}}\right)z\right|^2+\lambda \left(z^{}za\right)\right)}`$
The $`U(1)`$ gauge field $`\alpha _\mu `$ is a Lagrange multiplier which removes the redundant phase variable in the parametrization of $`CP^1`$ by $`z_A`$. A coupling between $`\alpha _\mu `$ and $`\chi _A`$ has been added to the first term and subtracted from the $`U^{}U`$ terms so as to make the latter invariant under the gauge transformation $`z_Ae^{i\theta }z_A`$. $`\lambda `$ is a Lagrange multiplier which fixes $`\mathrm{\Phi }_{\underset{¯}{i}}^2=a^2`$. We have introduced the external electromagnetic field, $`A_\mu `$, in order to keep track of the charge quantum numbers of the fields. When $`a`$ is large, $`\mathrm{\Phi }^{\underset{¯}{i}}=z^A\tau _A^{\underset{¯}{i}B}z_B`$ condenses and the system is in one of the $`d_{x^2y^2}`$ ordered states. When $`a`$ is small, $`\mathrm{\Phi }^{\underset{¯}{i}}`$ is disordered. There is a critical point at $`a=a_c`$.
The Nodal Liquid Revisited. In the ordered phases, $`U`$ is a constant, so the $`U^{}U`$ terms in (27) can be dropped; the nodal quasiparticles are coupled to the external electromagnetic field. Note that the $`d_{x^2y^2}`$ density wave is an ordered state in this formalism, unlike in , where it is a disordered state. In the disordered phases, the $`z_A`$ sector of the theory develops a gap. Hence, the fourth and fifth lines of (27) can be dropped at low energies. To analyze these phases further, we introduce a dual representation for $`z_A`$, following . The effective action now takes the form:
$`S_{\mathrm{eff}}`$ $`=`$ $`S_F[\chi _A,\alpha _\mu ]+{\displaystyle \underset{A}{}}S_{GL}[\mathrm{\Phi }^A,{\displaystyle \frac{1}{2}}(a_\mu ^+\pm a_\mu ^{})]`$ (29)
$`+{\displaystyle 𝑑\tau d^2x\left(\alpha _\mu ϵ_{\mu \nu \lambda }_\nu a_\lambda ^++A_\mu ϵ_{\mu \nu \lambda }_\nu a_\lambda ^{}\right)}`$
where $`S_F[\chi _A,\alpha _\mu ]`$ is the first three lines of (27), $`\mathrm{\Phi }^A`$ annihilates a vortex in $`z_A`$, and
$$_{GL}(\mathrm{\Phi },a_\mu )=\frac{1}{2}|(i_\mu a_\mu )\mathrm{\Phi }|^2+V(\mathrm{\Phi })+\frac{1}{2}(f_{\mu \nu })^2$$
(30)
and $`J_\mu ^\pm =ϵ_{\mu \nu \lambda }_\nu a_\lambda ^\pm `$ are the $`z_A`$ number and pseudospin $`\underset{¯}{3}`$ currents. When the $`Z_2`$ symmetry $`\mathrm{\Phi }^A\mathrm{\Phi }^A`$ is unbroken, we can rewrite the effective action in terms of the fields $`\mathrm{\Phi }^+=\mathrm{\Phi }^1\mathrm{\Phi }^2`$, $`\mathrm{\Phi }^{}=\mathrm{\Phi }^1\mathrm{\Phi }^2`$. We now have:
$`S_{\mathrm{eff}}`$ $`=`$ $`S_F[\chi _A,\alpha _\mu ]+S_{GL}[\mathrm{\Phi }^+,a_\mu ^+]+S_{GL}[\mathrm{\Phi }^{},a_\mu ^{}]`$ (32)
$`+{\displaystyle 𝑑\tau d^2x\left(\alpha _\mu ϵ_{\mu \nu \lambda }_\nu a_\lambda ^++A_\mu ϵ_{\mu \nu \lambda }_\nu a_\lambda ^{}\right)}`$
Integrating out $`\alpha _\mu `$, we can solve the resulting constraint to express $`a_\mu ^+`$ in terms of $`\chi _A`$: $`J_0^+=\chi ^{}\tau ^{\underset{¯}{3}}\chi `$, $`J_x^+=v_F\chi ^{}\tau ^{\underset{¯}{3}}\chi `$.
Now suppose that the system becomes disordered as a result of the condensation of $`\mathrm{\Phi }^{}`$. By the Anderson-Higgs mechanism, $`a_\mu ^{}`$ aquires a gap. Integrating out $`a_\mu ^{}`$, we find no coupling of $`A_\mu `$ to the remaining degrees of freedom: $`\chi _A`$ is a neutral spin-$`1/2`$ fermion. The change of variables (18) has effectively ‘bleached’ the fermions by using the order parameter to screen their pseudospin (including their charge). This state is none other than the Nodal Liquid.
Pseudospin-Peierls Order. If the system is, instead, disordered by the condensation of $`\mathrm{\Phi }^{1,2}`$, then $`J_\mu ^\pm `$ must vanish at low energies. The only allowed excitations at low energies are those combinations of $`\chi _A`$s which are invariant under $`\tau _A^{\underset{¯}{3}B}`$ rotations, i.e. neutral excitations. At finite energy, there are also solitonic excitations which carry one quantum of $`(a_\mu ^+\pm a_\mu ^{})/2`$ flux, i.e. charge $`e`$ and spin-$`1/2`$. According to the analogy between the pseudospin $`SU(2)`$ physics of our system and the spin $`SU(2)`$ physics of a quantum antiferromagnet, we might, in this disordered phase, expect the pseudospin analog of spin-Peierls order, pseudospin Peierls order,
$`\stackrel{}{\mathrm{\Phi }}(k+K)\stackrel{}{\mathrm{\Phi }}(k)\stackrel{}{\mathrm{\Phi }}\times _\tau \stackrel{}{\mathrm{\Phi }}(k+K)\stackrel{}{\mathrm{\Phi }}\times _\tau \stackrel{}{\mathrm{\Phi }}(k)`$ (33)
$`=\mathrm{sin}k_xa`$ (34)
with $`K=(\pi /a,0)`$ or $`(0,\pi /a)`$, as a result of Berry phases which we have neglected in (27).
Phase Transitions at Half-Filling. The transition at half-filling between the $`d_{x^2y^2}`$ density-wave and the $`d_{x^2y^2}`$ superconductor is driven by a pseudospin-$`2`$ symmetry-breaking field,
$`S_u=u{\displaystyle 𝑑\tau d^2x\left(\mathrm{\Phi }_3^2\mathrm{\Phi }_1^2\mathrm{\Phi }_2^2\right)}`$ (35)
For $`u<0`$, the $`\underset{¯}{3}`$-axis is an easy axis and the $`d_{x^2y^2}`$ density-wave state is favored; for $`u>0`$, the $`\underset{¯}{1}\underset{¯}{2}`$-plane is an easy plane and the $`d_{x^2y^2}`$ superconducting state is favored. At $`u=0`$, a first-order pseudospin-flop transition occurs, provided $`a>a_c`$. At the bicritical point $`a=a_c`$, $`u=0`$, quantum fluctuations destroy order at the $`O(4)`$-symmetric point. This bicritical point and the quantum critical region are described by the physics of the critical fluctuations coupled to nodal fermionic excitations. For $`a<a_c`$, $`u=0`$ the system lies along the $`O(4)`$-symmetric line in the Nodal Liquid phase. A small increase or decrease of $`u`$ will not cause order, and the system will still be in the nodal liquid phase, albeit with lower symmetry, $`U(1)\times Z_2`$. Further increase or decrease of $`u`$ will lead to second-order phase transitions at $`u_{cr}^\pm (a)`$ into the $`d_{x^2y^2}`$ superconducting and $`d_{x^2y^2}`$ density-wave phases respectively.
At the second-order transition from the Nodal Liquid to the $`d_{x^2y^2}`$ density-wave, the $`Z_2`$ symmetry of translation by one lattice site is broken. At the second-order $`XY`$ transition from the Nodal Liquid to the $`d_{x^2y^2}`$ superconductor, electromagnetic $`U(1)`$ is broken. At the first-order pseudospin flop transition between the $`d_{x^2y^2}`$ superconductor and the $`d_{x^2y^2}`$ density-wave, $`U(1)`$ is restored and $`Z_2`$ is simultaneously broken. In the formulation discussed here, spin-charge confinement – which, in the language of (see also ) is due to the absence of vortex pairing – occurs simultaneously with translational symmetry breaking.
Topological Excitations. We can give a narrative for the destruction of superconductivity in the language of vortex condensation. In the superconducting phase, the pseudospin $`\mathrm{\Phi }^{\underset{¯}{i}}`$ lies in the $`\underset{¯}{1}\underset{¯}{2}`$ plane. In the core of a vortex – a meron in the $`\sigma `$-model – $`\mathrm{\Phi }^{\underset{¯}{i}}`$ must point out of the $`\underset{¯}{1}\underset{¯}{2}`$ plane. This can be done by pointing along the $`\pm \underset{¯}{3}`$ axis. When $`+\underset{¯}{3}`$ merons dominate (in the presence of an infinitesimal $`Z_2`$ symmetry-breaking field), the superconductor undergoes a transition to the $`d_{x^2y^2}`$ density-wave state. When there are equal numbers of $`\pm \underset{¯}{3}`$ merons, the superconductor instead undergoes a transition to the disordered state. This condition on the densities of $`\pm \underset{¯}{3}`$ merons is reminiscent of and cognate to the vortex-pairing scenario of , but is weaker since it allows for the two possibilities discussed earlier. The transition from the $`d_{x^2y^2}`$ density-wave state to the disordered state can be understood in terms of skyrmion condensation.
Discussion. Transitions of the type which we have discussed above do not in the cuprates occur at half-filling but – if at all – near $`x_c`$, the doping at which superconductivity first appears. We assume $`u<0`$ to suppress superconductivity at half-filling. In order to move away from half-filling, we vary the chemical potential, which can be done by adding the $`O(4)`$-breaking term:
$`S_\mu =\mu O^3=\mu {\displaystyle 𝑑\tau d^2x\left(ϵ_{\underset{¯}{3}\underset{¯}{i}\underset{¯}{j}}\mathrm{\Phi }_{\underset{¯}{i}}_\tau \mathrm{\Phi }_{\underset{¯}{j}}+\mathrm{\Psi }^{}\tau ^{\underset{¯}{3}}\mathrm{\Psi }\right)}`$ (36)
By increasing $`\mu `$, we can drive the system through a first-order pseudospin-flop transition into the superconducting state. As $`a`$ is decreased, a bicritical point will again be reached. The coupling between $`z_A`$ and $`\chi _A`$ only enters at two-loops; at one-loop, we can appeal to known results for the pure non-linear $`\sigma `$-model, which indicate that the $`O(4)`$ symmetry is dynamically restored at the bicritical point . As a result, the $`O(4)`$-symmetric critical theory discussed above will apply in the low-frequency, long-wavelength limit. A possible phase diagram for the cuprates, based on this scenario, is depicted in figure 1. An alternative, not depicted in figure 1, can occur if $`\rho _s<\rho _{DW}`$. In this case, there can be a phase with both $`d_{x^2y^2}`$ superconducting and $`d_{x^2y^2}`$ density-wave order, and a tetracritical point, $`T=T_{bc}`$, $`\mu =\mu _{bc}`$, at which both orders become critical. For $`\mu <\mu _{bc}`$, there will be a regime, $`T_c^{sc}<T<T_c^{dw}`$, above the superconducting transition temperature, which has density-wave order.
The dotted line in figure 1 is the pseudogap scale, which we interpret as the scale below which $`\mathrm{\Phi }^{\underset{¯}{i}}`$ has fixed magnitude and the non-linear $`\sigma `$ model description is available. Let us consider the physics below this scale. As $`\mu `$ is increased, Fermi pockets open at the nodes of the $`d_{x^2y^2}`$ density-wave state. Eventually, the system undergoes a transition from the $`d_{x^2y^2}`$ density-wave to the $`d_{x^2y^2}`$ superconductor. The nature of this transition depends on the value of $`a`$ which, ostensibly, varies among the materials in the cuprate family. It may, perhaps, be controlled by chemical substitution or applied pressure. For $`a`$ large, the transition will be first-order as depicted by the thick line. For $`a`$ small, it occurs via two second-order phase transitions; the Nodal Liquid is sandwiched between these two transitions. Neither the $`d_{x^2y^2}`$ superconductor nor the Nodal Liquid has Fermi pockets, the latter because the second term in (36) can be dropped in the disordered phase. Appealing to the phase diagram of the spin-flop transition in magnetically-ordered systems, we extend the first-order phase transition to finite-temperature, where it meets the second-order $`d_{x^2y^2}`$ density-wave and superconducting ordering transitions.
Whither the antiferromagnet? As Hsu and Gros pointed out, the $`d_{x^2y^2}`$ density-wave state has good short-ranged antiferromagnetic correlations, reflected in its excellent numerical variational energy. Hence, we will assume that the only additional physics needed to describe the antiferromagnetic state at half-filling is a moderate triplet quasiparticle-quasihole condensate . This will not affect our description of the critical regime. Our assumption appears to be supported by photemission experiments on the antiferromagnetic insulator, $`Ca_2CuO_2Cl_2`$ . Similar ideas may apply to the Nodal Liquid state, making it an equally good platform for the antiferromagnetic state at half-filling.
Our non-linear $`\sigma `$-model analysis mirrors that of , but is on firmer footing because the $`d_{x^2y^2}`$ density-wave – unlike the antiferromagnet – has a nodal fermionic spectrum similar to that of the $`d_{x^2y^2}`$ superconductor into which the pseudospin symmetry rotates it. Fluctuations between the $`d_{x^2y^2}`$ density-wave and superconducting states are also a key feature of the $`SU(2)`$ mean-field-theory of the $`tJ`$ model. In fact, a parallel approach to the Nodal Liquid state was taken in this framework in . However, the $`SU(2)`$ is local in that approach, which leads to complications arising from the concomitant gauge field. One virtue of the non-linear $`\sigma `$-model approach is that we can use the physics of quantum antiferromagnets as a guide. In this way, we identified pseudospin-Peierls order as a possible alternative to the Nodal Liquid phase. Another striking upshot of our analysis is the bicritical point at which the $`d_{x^2y^2}`$ density-wave, $`d_{x^2y^2}`$ superconducting, and Nodal Liquid phases touch. It is possible that it is responsible for recent experimental hints of quantum critical behavior in the cuprates .
I would like to thank S. Chakravarty for discussions, and S. Sachdev and T. Senthil for pointing out an error in an earlier version of this paper. |
warning/0001/hep-ph0001090.html | ar5iv | text | # Thermal Abundances of Heavy Particles
## Abstract
Matsumoto and Yoshimura \[hep-ph/9910393\] have argued that there are loop corrections to the number density of heavy particles (in thermal equilibrium with a gas of light particles) that are not Boltzmann suppressed by a factor of $`e^{M/T}`$ at temperatures $`T`$ well below the mass $`M`$ of the heavy particle. We argue, however, that their definition of the number density does not correspond to a quantity that could be measured in a realistic experiment. We consider a model where the heavy particles carry a conserved U(1) charge, and the light particles do not. The fluctuations of the net charge in a given volume then provide a measure of the total number of heavy particles in that same volume. We show that these charge fluctuations are Boltzmann suppressed (to all orders in perturbation theory). Therefore, we argue, the number density of heavy particles is also Boltzmann suppressed.
preprint: hep-ph/0001090 January 2000 Revised March 2000
Physical Review D, in press
In a series of papers, Matsumoto and Yoshimura (hereafter MY) have challenged the conventional wisdom concerning the number density of a gas of heavy particles in thermal equilibrium with a gas of light (or massless) particles. They consider a model of a heavy spin-zero boson (represented by a real scalar field $`\phi `$) interacting with a massless spin-zero boson (represented by a real scalar filed $`\chi `$). The scalar potential is
$$V(\phi ,\chi )=\frac{1}{2}M^2\phi ^2+\frac{1}{24}\lambda _\phi \phi ^4+\frac{1}{24}\lambda _\chi \chi ^4+\frac{1}{4}\lambda \phi ^2\chi ^2,$$
(1)
where the couplings are all real and positive, and we take $`\lambda _\phi \lambda ^2`$, $`\lambda \lambda _\chi `$, and $`\lambda _\chi <1`$. This hierarchy among the couplings allows the $`\chi `$ particles to function as an efficient heat bath for the $`\phi `$ particles.
At the level of free field theory, the equilibrium number density of $`\phi `$ particles is given by
$$n_{\phi 0}=g\frac{d^3p}{(2\pi )^3}f(E),$$
(2)
where $`g=1`$ counts the number of species of $`\phi `$ particles,
$$f(E)=\frac{1}{e^{\beta E}1}$$
(3)
is the Bose distribution function, $`\beta =1/T`$ is the inverse temperature, and $`E=(𝐩^2+M^2)^{1/2}`$ is the single-particle energy. For $`TM`$, we have
$$n_{\phi 0}=(2\pi )^{3/2}(MT)^{3/2}e^{\beta M};$$
(4)
the factor of $`e^{\beta M}`$ means that $`n_{\phi 0}`$ is Boltzmann suppressed. However, MY argue that there are loop corrections to $`n_{\phi 0}`$ that are not Boltzmann suppressed; specifically, they find
$$n_\phi =n_{\phi 0}+c\lambda ^2T^6/M^3+\mathrm{}$$
(5)
for $`TM`$, where $`c=1/69120`$, and the ellipses stand for all higher-order corrections.
They key issue that we wish to address (raised also in ) is the underlying definition of $`n_\phi `$. For $`TM`$, MY define $`n_\phi `$ via $`n_\phi =\rho _\phi /M`$, where $`\rho _\phi `$ is the energy density of the $`\phi `$ particles. This energy density is in turn defined (for all temperatures) via
$$\rho _\phi =_\phi =\frac{Tr_\phi e^{\beta H}}{Tre^{\beta H}}0|_\phi |0,$$
(6)
where $`H`$ is the total hamiltonian, and
$$_\phi =\frac{1}{2}\dot{\phi }^2+\frac{1}{2}(\phi )^2+\frac{1}{2}M^2\phi ^2+\text{counterterms}$$
(7)
is the free-field part of the $`\phi `$ hamiltonian, plus counterterms (some of which involve the $`\chi `$ field) that are necessary to remove infinities in this composite operator.
Eq. (6) is a highly plausible definition of $`\rho _\phi `$. However, it does not correspond in any obvious way to how the number density of $`\phi `$ particles would be determined experimentally. Standard methods all involve a search for individual, on-shell $`\phi `$ particles. Real-world examples of this include present-day dark matter searches, and measurements of the cosmic microwave background radiation.
What is needed theoretically is a measurable attribute that is carried by the $`\phi `$ particles only. To create one, we modify the model slightly by making the $`\phi `$ field complex, and requiring its interactions to conserve the corresponding U(1) charge. We leave the $`\chi `$ field real, and the $`\chi `$ particles neutral. The modified scalar potential is
$$V(\phi ,\chi )=M^2\phi ^{}\phi +\frac{1}{4}\lambda _\phi (\phi ^{}\phi )^2+\frac{1}{24}\lambda _\chi \chi ^4+\frac{1}{2}\lambda \phi ^{}\phi \chi ^2.$$
(8)
We can now study the net charge contained in a large but finite volume $`V`$. Of course, in thermal equilibrium, the average net charge $`Q`$ vanishes, but it has nonzero fluctuations $`Q^2`$. If we weakly gauge the U(1) symmetry with a small (and therefore dynamically irrelevant) gauge coupling $`e\lambda `$, we can in principle measure these charge fluctuations without tracking individual $`\phi `$ particles.
For $`TM`$, it is easy to compute $`Q^2_0`$, where the subscript 0 indicates that we are (for now) neglecting interactions. The number $`N_+`$ of positively charged particles in a volume $`V`$ is then controlled by a Poisson distribution; this implies $`N_+^2_0N_+_0^2=N_+_0`$. The number $`N_{}`$ of negatively charged particles is controlled by an independent Poisson distribution, with $`N_{}^2_0N_{}_0^2=N_{}_0`$. Overall charge neutrality implies $`N_{}_0=N_+_0=\frac{1}{2}n_{\phi 0}V`$, where $`n_{\phi 0}`$ is now given by Eq. (2) with $`g=2`$. The net charge is $`Q=N_+N_{}`$, and so
$`Q^2_0`$ $`=`$ $`(N_+N_{})^2_0`$ (9)
$`=`$ $`N_+^2_0+N_{}^2_02N_+_0N_{}_0`$ (10)
$`=`$ $`n_{\phi 0}V.`$ (11)
We see that the charge fluctuations give us a measurement of the total number of $`\phi `$ particles in a given volume. At higher temperatures (but still ignoring interactions), quantum effects modify this result to
$$Q^2_0=2V\frac{d^3p}{(2\pi )^3}f(E)[1+f(E)].$$
(12)
The ratio $`Q^2_0/n_{\phi 0}V`$ depends weakly on temperature; it rises slowly from one at $`TM`$ to $`\pi ^2/6\zeta (3)=1.37`$ at $`TM`$. (The result is similar for fermions, with a ratio of one at low temperatures and $`\pi ^2/9\zeta (3)=0.91`$ at high temperatures.) Furthermore, it seems highly unlikely that weak interactions could significantly modify Eq. (11). If we were to have either $`Q^2n_\phi V`$ or $`Q^2n_\phi V`$, we would be forced to conclude that the movements of positive and negative particles are highly correlated (in order to suppress or enhance the charge fluctuations in any particular volume). This is inconsistent with the usual notion of a gas of particles that move freely and independently between occasional scatterings, and would appear to require strong interactions.
We therefore propose to define the number density of $`\phi `$ particles, for $`TM`$, via
$$n_\phi =Q^2/V.$$
(13)
This definition has the advantage (not shared by the definition used by MY) of being directly connected to the experimentally measurable quantity $`Q^2`$. Adopting Eq. (13) as our definition of $`n_\phi `$, the question becomes whether or not the loop corrections to $`Q^2`$ are Boltzmann suppressed.
To compute these loop corrections, we introduce a chemical potential $`\mu `$ and the partition function
$$Z=Tre^{\beta (H\mu Q)}.$$
(14)
We then use
$$Q^2=\frac{1}{\beta ^2}\frac{^2}{\mu ^2}\mathrm{ln}Z|_{\mu =0}.$$
(15)
At the one-loop level (that is, for free $`\phi `$ particles, and ignoring the $`\mu `$-independent contributions of the $`\chi `$ particles), we have the textbook formula
$$\mathrm{ln}Z_0=V\frac{d^3p}{(2\pi )^3}\mathrm{ln}[1+f(E\mu )]+(\mu \mu ).$$
(16)
Using Eq. (16) in Eq. (15) yields Eq. (12).
As in , loop corrections to $`\mathrm{ln}Z`$ may be computed via finite temperature perturbation theory (see, e.g., ). At the two- and three-loop level, the contributing diagrams are shown in Fig. (1). Vertices joining four $`\phi `$ lines are neglected since we have assumed that $`\lambda _\phi \lambda ^2`$. Furthermore we neglect diagrams with only $`\chi `$ lines, since these are independent of $`\mu `$. The $`\phi `$ propagator is
$$\mathrm{\Delta }(n,𝐩)=\frac{1}{(2\pi n/\beta i\mu )^2+𝐩^2+M^2},$$
(17)
where $`n`$ is an integer specifying the discrete energy. (With $`M=\mu =0`$, this is also the $`\chi `$ propagator.) The direction of the flow of charge is the same as the flow of energy. A closed $`\phi `$ loop with a single vertex contributes a factor of
$$T\underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}\mathrm{\Delta }(n,𝐩)=\frac{1}{2E}\left[1+f(E+\mu )+f(E\mu )\right],$$
(18)
where $`E=(𝐩^2+M^2)^{1/2}`$. The “1” in square brackets yields a divergent, temperature independent term that must be removed by renormalization. The remaining two terms are Boltzmann suppressed, since at low temperatures
$$f(E\pm \mu )e^{\beta M}e^{\beta \mu }e^{\beta 𝐩^2/2M}.$$
(19)
This implies that the two-loop diagram in Fig. (1), and the first two three-loop diagrams, are all Boltzmann suppressed. The third three-loop diagram has a factor of $`_n\mathrm{\Delta }(n,𝐩)^2`$, and turns out to be suppressed as well. This leaves only the last (“basketball”) diagram.
According to the standard rules of finite temperature perturbation theory , the contribution of the basketball diagram is
$`{\displaystyle \frac{1}{\beta V}}\mathrm{ln}Z_{\mathrm{bb}}`$ $`=`$ $`C\lambda ^2T^4{\displaystyle \underset{n\text{}\mathrm{s}}{}}{\displaystyle \frac{d^3p_1}{(2\pi )^3}\frac{d^3p_2}{(2\pi )^3}\frac{d^3k_3}{(2\pi )^3}\frac{d^3k_4}{(2\pi )^3}(2\pi )^3\delta ^3(𝐩_1+𝐩_2+𝐤_3+𝐤_\mathrm{𝟒})}`$ (21)
$`\times \beta \delta _{n_1+n_2+n_3+n_4,0}\mathrm{\Delta }_1\mathrm{\Delta }_2\mathrm{\Delta }_3\mathrm{\Delta }_4.`$
Here $`𝐩_1`$ and $`𝐩_2`$ are the momenta on the $`\phi `$ lines, and $`𝐤_3`$ and $`𝐤_4`$ are the momenta on the $`\chi `$ lines. Each momentum and energy is taken to flow from the left vertex to the right vertex. Thus we have, in the individual propagators, masses $`M_1=M_2=M`$ and $`M_3=M_4=0`$, and chemical potentials $`\mu _1=\mu _2=\mu `$ and $`\mu _3=\mu _4=0`$. Note especially that $`\mu _2=\mu _1`$; this is because, with our convention on energy flow, the charge flow must be opposite to the energy flow on one of the $`\phi `$ lines, and this changes the sign of the chemical potential. Finally,
$$C=\frac{1}{2!}\left(\frac{1}{2}\right)^22=\frac{1}{4}$$
(22)
is a combinatoric factor. The $`1/2!`$ comes from the expansion of $`\mathrm{exp}(\beta H_{\mathrm{int}})`$, the $`(1/2)^2`$ comes from two vertex factors arising from the last term in the scalar potential, Eq. (8), and the $`2`$ comes from the number of ways to match up the $`\chi `$ lines.
Eq. (21) can be evaluated by the standard procedure of writing the Kroneker delta as an integral,
$$\beta \delta _{n_1+\mathrm{}+n_4,0}=_0^\beta 𝑑te^{2\pi i(n_1+\mathrm{}+n_4)t/\beta },$$
(23)
and then performing each sum via contour integration,
$`{\displaystyle \underset{n=\mathrm{}}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{e^{2\pi int/\beta }}{(2\pi n/\beta i\mu )^2+E^2}}`$ $`=`$ $`{\displaystyle \frac{\beta }{2\pi }}{\displaystyle 𝑑z\frac{e^{izt}}{e^{i\beta z}1}\frac{1}{(zi\mu )^2+E^2}}`$ (24)
$`=`$ $`{\displaystyle \frac{\beta }{2E}}\left[e^{(\beta t)(E+\mu )}f(E+\mu )+e^{t(E\mu )}f(E\mu )\right].`$ (25)
Here the $`z`$ contour originally encloses the real axis, and then is deformed into two small circles surrounding the poles at $`z=i(\mu \pm E)`$; this deformation is allowed for $`0t\beta `$. Putting all of this together we now have
$$\frac{1}{\beta V}\mathrm{ln}Z_{\mathrm{bb}}=C\lambda ^2\stackrel{~}{dp}_1\stackrel{~}{dp}_2\stackrel{~}{dk}_3\stackrel{~}{dk}_4(2\pi )^3\delta ^3(𝐩_1+𝐩_2+𝐤_3+𝐤_\mathrm{𝟒})_0^\beta 𝑑tD_1D_2D_3D_4,$$
(26)
where $`\stackrel{~}{dp}_i=d^3p_i/(2\pi )^3(2E_i)`$, and
$$D_i=\left[e^{(\beta t)(E_i+\mu _i)}f(E_i+\mu _i)+e^{t(E_i\mu _i)}f(E_i\mu _i)\right].$$
(27)
The integral over $`t`$ is tedious but straightforward to perform; the result is
$`{\displaystyle _0^\beta }𝑑tD_1D_2D_3D_4`$ $`=`$ $`{\displaystyle \underset{\delta _1=0}{\overset{1}{}}}{\displaystyle \underset{\delta _3=0}{\overset{1}{}}}{\displaystyle \underset{\delta _4=0}{\overset{1}{}}}[{\displaystyle \frac{(\delta _1+f_{1+})(\delta _1+f_2)(\delta _3+f_3)(\delta _4+f_4)}{\epsilon _1E_1+\epsilon _1E_2+\epsilon _3E_3+\epsilon _4E_4}}`$ (29)
$`+{\displaystyle \frac{(\delta _1+f_{1+})(1\delta _1+f_{2+})(\delta _3+f_3)(\delta _4+f_4)}{\epsilon _1E_1\epsilon _1E_2+\epsilon _3E_3+\epsilon _4E_4}}]+(\mu \mu ),`$
where $`\epsilon _i=(1)^{1\delta _i}`$, and $`f_{1+}=f(E_1+\mu )`$, $`f_2=f(E_2\mu )`$, $`f_3=f(E_3)`$, etc. We have made repeated use of the relation
$$e^{\beta E}f(E)=1+f(E)$$
(30)
in obtaining Eq. (29).
We are interested only in those terms (if any) in Eq. (29) that are not Boltzmann suppressed. As we see from Eq. (19), any term containing a factor of either $`f_{1\pm }`$ or $`f_{2\pm }`$ is Boltzmann suppressed. Thus, we can drop all such terms. At this point, we see immediately that the remaining terms are all independent of $`\mu `$, since the only $`\mu `$ dependence in Eq. (29) is that which is contained implicitly in the factors of $`f_{1\pm }`$ and $`f_{2\pm }`$. Therefore, all $`\mu `$-dependent contributions (at the two- and three-loop level) to $`\mathrm{ln}Z`$ are Boltzmann suppressed. Eqs. (15) and (13) then imply that $`Q^2`$ and $`n_\phi `$ are also Boltzmann suppressed. This is our main result.
For completeness, and for comparison with the results of MY , we continue the computation of the unsuppressed terms in $`\mathrm{ln}Z_{\mathrm{bb}}`$. We set $`f_{1\pm }=f_{2\pm }=0`$. All numerators in Eq. (29) are then either $`1`$, $`f_3`$, $`f_4`$, or $`f_3f_4`$. The first of these yields a divergent, temperature independent term that is canceled by a renormalization of the vacuum energy. The second and third yield divergent, temperature dependent terms that are canceled by renormalization of the $`\phi `$ and $`\chi `$ thermal self-energies (see the corresponding discussion for QED in ). Dropping these terms, we have
$$_0^\beta 𝑑tD_1D_2D_3D_4\underset{\epsilon _3=\pm 1}{}\underset{\epsilon _4=\pm 1}{}\frac{2f_3f_4}{E_1+E_2+\epsilon _3E_3+\epsilon _4E_4}.$$
(31)
Substituting this into Eq. (26), making a low temperature expansion, and dropping a final divergent contribution to the vacuum energy ultimately yields
$$\frac{1}{\beta V}\mathrm{ln}Z_{\mathrm{bb}}=\frac{\pi ^2\lambda ^2T^8}{648000M^4}+O(T^{10}).$$
(32)
The corresponding correction to the total energy density $`\rho `$ is then given by
$$\delta \rho =\frac{1}{V}\frac{}{\beta }\mathrm{ln}Z_{\mathrm{bb}}=\frac{7\pi ^2\lambda ^2T^8}{648000M^4}+O(T^{10}).$$
(33)
This is consistent with the results of MY , who also find that the leading correction to $`\rho `$ is of order $`T^8/M^4`$.
That there are corrections of this form to $`\rho `$ is not surprising. Without interactions, we have $`\rho =(\pi ^2/30)T^4`$ from the $`\chi `$ particles, plus the Boltzmann-suppressed contribution of the $`\phi `$ particles. At temperatures $`TM`$, we should be able to integrate the heavy $`\phi `$ field out of the functional integral, and be left with an effective lagrangian for the massless $`\chi `$ field alone. This lagrangian will contain nonrenormalizable interaction terms that are suppressed by powers of $`\lambda `$ and inverse powers of $`M`$. These will give rise to corrections like Eq. (33). From this point of view, it is clear that all such corrections should be thought of as modifications of the energy of the $`\chi `$ particles, and not as unsuppressed contributions to the energy of the $`\phi `$ particles.
We now argue that $`Q^2`$ (and hence $`n_\phi `$) is Boltzmann suppressed to all orders in perturbation theory. Consider an exact evaluation of the partition function
$$Z=\underset{\alpha }{}e^{\beta (E_\alpha \mu Q_\alpha )},$$
(34)
where the sum is over a basis of energy and charge eigenstates. States that yield a $`\mu `$-dependent contribution to $`Z`$ must have $`Q0`$. States consisting of a single stable $`\phi `$ particle are well defined exact energy eigenstates, with $`E=(𝐩^2+M^2)^{1/2}`$ and $`Q=\pm 1`$; their contribution to $`Z`$ is obviously Boltzmann suppressed. Furthermore, any other state with $`Q=\pm 1`$ must have energy $`E>M`$; otherwise, the $`\phi `$ particle would not be stable (since it could decay into this lighter charged state). At tree level, energy eigenstates with two or more $`\phi `$ particles all have $`E2M`$; interactions can modify this to $`E2M[1+O(\lambda ^2)]`$. The $`O(\lambda ^2)`$ term could in principle be negative (if there is a bound state of two or more like-charge particles), but cannot make $`E2M`$ in a weakly coupled theory. We therefore conclude that all $`\mu `$-dependent contributions to $`Z`$ are Boltzmann suppressed. We have already seen this explicitly at the level of two and three loops.
With this in mind, we can return to the original model of MY, Eq. (1). At the level of Feynman diagrams, the models are essentially identical, the only difference being the combinatoric factor associated with each diagram. We therefore expect that $`n_\phi `$ should be Boltzmann suppressed in this model as well. We are still in need, however, of a general definition of $`n_\phi `$ that does not rely on the trick of Eq. (13), but which corresponds to experimental measurements.
We conclude, in accord with , that the proper definition of the number density $`n_\phi `$ of heavy particles (in equilibrium with a thermal bath of light particles) is a delicate matter. We have argued that it is essential that $`n_\phi `$ be defined in a manner that renders it measurable in a realistic experiment. In this paper, we have considered a model in which the heavy particles carry a conserved U(1) charge, while the light particles are neutral. In this situation, the local charge fluctuations provide an experimentally accessible measure of the number density of heavy particles. We have shown explicitly that, at low temperatures, this number density is Boltzmann suppressed up through three-loop order in perturbation theory, and we have argued that this must in fact be true to all orders.
###### Acknowledgements.
I thank Anupam Singh and Scott Thomas for discussions. This work was supported in part by the National Science Foundation through grant PHY–97–22022, and by the Institute of Geophysics and Planetary Physics through grant 920. |
warning/0001/math0001107.html | ar5iv | text | # Some results on rational surfaces and Fano varieties
## 0. Introduction
The goal of this article is to study the equations and the syzygies of embeddings of rational surfaces and certain Fano varieties. Previously Butler, Homma, Kempf, and the authors had proved results regarding syzygies of (geometrically) ruled surfaces and surfaces of nonnegative Kodaira dimension. We will be interested in knowing under what conditions the resolution of the homogeneous coordinate ring $`S/I`$ of an embedded variety is “simple”. More precisely we want to know under what conditions the so-called property $`N_p`$ after M. Green is satisfied. We define this property next:
###### Definition 0.1
Let $`X`$ be a projective variety. A very ample line bundle $`L`$ is said to satisfy property $`N_0`$ if $`|L|`$ embeds $`X`$ as a projectively normal variety. A very ample line bundle $`L`$ satisfies property $`N_1`$ if $`L`$ satisfies property $`N_0`$ and the homogeneous ideal $`I`$ of the image of $`X`$ embedded by $`|L|`$ is generated by quadratic equations. Finally a very ample line bundle $`L`$ is said to satisfy property $`N_p`$, $`p1`$, if it satisfies property $`N_1`$ and the matrices in the minimal graded free resolution of $`S/I`$ have linear entries from the second to the $`p`$-th step.
Section 1 is devoted to rational surfaces. Given a rational surface $`X`$ and an ample and base-point-free line bundle $`L`$ on $`X`$, we observe there is an optimal criterion for $`L`$ to satisfy property $`N_p`$. This criterion (cf. Theorem 1.3 ) depends solely on the intersection number of $`L`$ with $`K_X`$, namely, $`L`$ satisfies property $`N_p`$ if $`K_XLp+3`$. This criterion turns out to be a characterization if $`X`$ is anticanonical, i.e., a rational surface with effective anticanonical class. Anticanonical surfaces have been extensively studied by several authors, among them B. Harbourne, and Theorem 1.3 improves and generalizes one of his (cf. \[Hb2\] ).
We also study the syzygies associated to adjunction bundles and prove a Reider type theorem for higher syzygies. More precisely, given any $`A_1,\mathrm{},A_n`$ ample line bundles on an anticanonical rational surface $`X`$ of fixed $`K_X^2`$, we find a sharp bound on $`n`$ so that $`K_X+A_1+\mathrm{}+A_n`$ satisfies property $`N_p`$. One of the easy consequences of this is Mukai’s conjecture, which is in fact optimal for anticanonical surfaces with $`K_X^2=1`$.
Properties such as base-point-freeness and very ampleness are governed numerically, as classical results on curves and Reider’s theorem on surfaces show. It is natural to ask whether the same philosophy holds for the property $`N_p`$, which are a natural generalization of base-point-freeness and very ampleness. For curves Green’s theorem (cf. \[G\] , Theorem 4.e.1), which says that a line bundle of degree greater than or equal to $`2g+p+1`$ satisfies property $`N_p`$, provides an affirmative answer. For rational surfaces results in this article, namely, the already mentioned Theorem 1.3 and Theorem 1.24 , which is a Reider type result for property $`N_p`$, show us as well that property $`N_p`$ only depends on numerical criteria. In fact, if $`X`$ is anticanonical, Theorem 1.3 tells that property $`N_p`$ for $`L`$ is exclusively governed by the intersection number of $`L`$ with a particular curve lying on $`X`$, namely, an anticanonical curve. This phenomenon can be also observed in other surfaces such as elliptic ruled surfaces. Indeed, Homma and the authors gave in \[Ho1\] , \[Ho2\] and \[GP1\] a characterization for properties $`N_0`$ and $`N_1`$ in terms of the intersection number of $`L`$ with a few curves lying on the elliptic surface. In this case the relevant curves are a minimal section and a fiber of the elliptic ruled surface in addition to the anticanonical curve. The authors also gave in \[GP2\] a criterion for property $`N_p`$ in terms of the intersection number of $`L`$ with the three above mentioned curves, and conjectured that those intersection numbers should also characterize property $`N_p`$ as they did characterize property $`N_0`$ and $`N_1`$. It is surprising that results of such similar spirit hold for both rational surfaces and elliptic ruled surfaces. Even though both are surfaces of Kodaira dimension $`\mathrm{}`$, they differ in the fact that the Picard group of a rational surface is discrete, the same does not happen for an elliptic ruled surface. The differences can be also be seen in the methods of proof used in \[GP1\] and \[GP2\] and in this article, which are indeed very distinct.
We also prove a result (cf. Theorem 1.29) connecting property $`N_p`$, which is an extrinsic property depending on the embedding of $`X`$, with the “termination” of ampleness of $`mK_X+L`$. In particular we show for what $`m`$ the line bundle $`mK_X+L`$ stops being ample for a line bundle $`L`$ satisfying property $`N_p`$. The formula we obtain for such an $`m`$ depends on $`K_X^2`$ and $`p`$.
In Section 1 we also construct several families of examples. These examples show that all the theorems and propositions proved are sharp, and that the bounds cannot be reduced.
In Section 2 we study $`n`$-dimensional Fano varieties of index greater than or equal to $`n1`$. Let $`H`$ be primitive line bundle such that $`K_X=iH`$. First we prove Theorem 2.1 , which tells exactly what property $`N_p`$ is satisfied by $`H`$. We prove results regarding very ampleness and on the higher syzygies of multiples of a primitive $`H`$ on $`X`$ such that $`K_X=iH`$. We derive these syzygy results from the vanishings of certain Koszul cohomology groups on the Fano variety $`X`$. We reduce these vanishings to the vanishings of similar Koszul cohomology groups on lower dimensional Fano varieties. These lower dimensional are subvarieties of $`X`$, but since we need them to be again Fano varieties, we do not obtain them by taking subsequent hyperplane sections. We take this more indirect approach because the techniques of Section 1 do not work for these higher dimensional varieties. The reason is that the information available on the resolution of the coordinate ring of the subsequent hyperplane sections is not good enough.
Finally in Section 3 we deal with $`n`$-dimensional Fano varieties of index $`n3`$. We first give a criterion as to when $`nH`$ satisfies very ampleness and property $`N_0`$ when $`n2`$. This criterion is actually a characterization if $`n3`$. Then we prove a result on the higher syzygies of multiples of $`H`$. The arguments and techniques used are similar to those used in Section 2. There is though a difference worth noting. When one is considering an $`n`$-dimensional Fano variety $`X`$ of index greater than or equal to $`n1`$, one deduces the vanishing of Koszul cohomology groups on $`X`$ from the vanishings of similar cohomology groups on a rational surface. However, if $`X`$ has index $`n3`$, one deduces the vanishing of Koszul cohomology groups on $`X`$ from the vanishings of similar cohomology groups on a Calabi–Yau threefold, and eventually, on a surface of general type. The reader might wonder about what the situation is for $`n`$-dimensional Fano varieties of index $`n2`$. Those surfaces were studied in \[GP5\] , where the authors obtained results which are similar to those in Sections 2 and 3 of this article.
## 1. Rational surfaces
In this section we study property $`N_p`$ for rational surfaces. The first result is a criterion for property $`N_p`$ in terms of a very precise numerical condition, namely the intersection number of the line bundle $`L`$ under consideration with the anticanonical class of the surface. In order to prove this theorem we will need two preliminary results we mention now:
###### \restobs(\GPfour\ns, Observation 2.3)
Let $`X`$ be a regular variety (i.e., a variety such that $`H^1(𝒪_X)=0).`$ Let $`E`$ be a vector bundle on $`X`$, let $`C`$ be a divisor such that $`L`$ $`=𝒪_X\left(C\right)`$ is globally generated line bundle and $`H^1(EL^1)=0.`$ If the multiplication map $`H^0(E𝒪_C)H^0(L𝒪_C)H^0(EL𝒪_C)`$ surjects, then the map $`H^0(E)H^0(L)H^0(EL)`$ also surjects.
The second result we need is a useful lemma on the vanishing of cohomology of big and base-point-free line bundles:
###### \vanlemma
Let $`X`$ be a smooth surface with $`p_g=q=0`$ and $`B`$ a big base-point-free line bundle on $`X`$ such that $`K_XB>0`$. Then $`h^1(B)=h^2(B)=0`$.
Proof. Since $`B`$ is big and base-point-free, by Bertini, there exists smooth, irreducible $`C|B|`$. Consider
$$0𝒪_XBB𝒪_C0.$$
Since $`p_g=0`$, $`h^2(B)=0`$. Since $`K_XB>0`$, by adjunction deg$`(B𝒪_C)>2g(C)2`$, and consequently, $`h^1(B)=0`$, since $`q=0`$. $`\mathrm{}`$
We are now ready to state and prove the numerical criterion for property $`N_p`$. This criterion turns out to be a characterization of property $`N_p`$ if the surface is anticanonical. We remark that the case of property $`N_0`$ was observed by B. Harbourne (cf. \[Hb2\] ).
###### \rsNp
Let $`X`$ be a rational surface and let $`L`$ be an ample line bundle on $`X`$. If $`L`$ is base-point-free and $`K_XLp+3`$, then $`L`$ satisfies property $`N_p`$. In addition, if $`X`$ is anticanonical and $`L`$ is ample, then $`L`$ satisfies property $`N_p`$ if and only if $`K_XLp+3`$.
Proof. First we prove the part of the result stated for general rational surfaces, and we start showing that if $`K_XL3`$, $`L`$ satisfies property $`N_0`$. We want to show that
$$H^0(rL)H^0(L)\mathrm{@}>\alpha >>H^0((r+1)L)\text{ for all }r1.$$
Since $`L`$ is ample and base-point-free, we can choose a smooth and irreducible curve $`C|L|`$. Then, by Lemma 1.2 and Observation 1.1 , the surjectivity of $`\alpha `$ follows from the surjectivity of
$$H^0(rL_C)H^0(L_C)\mathrm{@}>\beta >>H^0((r+1)L_C)\text{ for all }r1.$$
Since $`K_XL3`$, deg$`L_C2g(C)+1`$, hence by Castelnuovo’s Theorem $`\beta `$ surjects. This proves that $`L`$ satisfies property $`N_0`$.
Now we prove the result for general $`p`$. We have just proven that if $`K_XLp+3`$, $`L`$ satisfies property $`N_0`$, i.e., $`L`$ is very ample and embeds $`X`$ as a projectively normal variety. On the other hand, by Lemma 1.2 , Kodaira Vanishing Theorem and because $`X`$ is regular, $`H^1(rL)=0`$ for all $`r`$. Therefore the image $`Y`$ of $`X`$ by the embedding induced by $`|L|`$ is arithmetically Cohen-Macaulay. Let $`H`$ be a general hyperplane of $`^N`$ and let $`D`$ be the corresponding (smooth, irreducible) divisor on $`X`$. Since $`X`$ is regular, $`|L_D|`$ embeds $`D`$ in $`H`$ and the image of this embedding is $`YH`$, which is projectively normal. This is the same as saying that $`L_D`$ satisfies property $`N_0`$. Moreover, since $`Y`$ is arithmetically Cohen-Macaulay, the minimal resolution of $`YH`$ has the same graded Betti as the minimal resolution of $`Y`$ (see also \[G\] , Theorem 3.b.7). Since $`K_XLp+3`$ we have by adjunction that deg$`L_D2g+p+1`$, then by Green’s theorem (cf. \[G\] , Theorem 4.a.1) the homogeneous coordinate ring of $`YH`$ is generated by quadrics and the resolution of its homogeneous coordinate ring is linear until the $`p`$th step. Since we did already prove that $`L`$ satisfied property $`N_0`$, we see that $`L`$ satisfies property $`N_p`$.
Let now $`X`$ be anticanonical. On an anticanonical surface a nef line bundle $`L`$ such that $`K_XL2`$ is base-point-free (cf. \[Hb1\] , Theorem III.1). Therefore we have just proven one of the implications, namely that if $`K_XLp+3`$, then $`L`$ satisfies property $`N_p`$. We prove now the other implication. Let $`L`$ be a line bundle on $`X`$ satisfying property $`N_p`$. In particular $`L`$ is very ample and embeds $`X`$ as a projectively normal variety. Moreover, $`H^1(rL)=0`$ for all $`r`$. Indeed, $`H^1(𝒪_X)=0`$ because $`X`$ is rational and $`H^1(rL)=0`$ by Kodaira Vanishing Theorem if $`r`$ is negative. If $`r`$ is positive, $`K_XrL>0`$, because $`L`$ is ample and $`K_X`$ is effective. Hence $`H^1(rL)=0`$ by Lemma 1.2 . Therefore the image of the embedding of $`X`$ by $`|L|`$ is arithmetically Cohen-Macaulay. Let $`C`$ be a smooth curve in $`|L|`$. As observed before, $`L`$ satisfies property $`N_p`$ if and only if $`L_C`$ satisfies property $`N_p`$. Since $`K_X`$ is effective and $`C`$ moves (for $`L`$ is very ample), then $`K_X𝒪_C`$ is also effective. Hence $`L_C=K_C+N`$ with $`N`$ effective line bundle of degree $`K_XL`$. Now if $`K_XL`$ were less than or equal to $`p+2`$, $`L_C`$ would not satisfy property $`N_p`$. This follows from a result of Green and Lazarsfeld (cf. \[GL\] , Theorem 2) which says in particular that on a smooth, irreducible, genus $`g`$ curve $`C`$, a line bundle of the form $`K_C+N`$ with $`N`$ effective of degree $`p+2`$ does not satisfies property $`N_p`$.
###### \question
If $`X`$ is a rational surface which is not anticanonical and $`L`$ is a line bundle on $`X`$ which satisfies $`N_p`$, in general, $`K_XL`$ need not be greater than or equal to $`p+3`$. However, there are line bundles $`L`$ such that $`K_XL=p+3`$ and $`L`$ satisfies $`N_p`$ and not $`N_{p+1}`$. Hence $`Theorem1.3`$ is sharp for non-anticanonical surfaces.
Proof. Consider a non-anticanonical rational surface $`X`$ and let $`L`$ be an ample and base-point-free line bundle so that $`K_X+L`$ is also ample, and such that the general curve in $`|L|`$ is not hyperelliptic. Then $`K_XL=p+3`$ for some $`p`$. Let $`C`$ be general and therefore smooth and irreducible curve in $`|L|`$. We know by Theorem 1.3 that $`L`$ satisfies property $`N_p`$. Since $`K_X+L`$ is ample, by Kodaira vanishing and duality, $`H^1(K_XL)=0`$. Since $`X`$ is not anticanonical, $`K_X𝒪_C`$ is non effective. Then, arguing as in the proof of Theorem 1.3 we conclude, using \[GL\] , Theorem 2, that $`L`$ satisfies property $`N_{p+1}`$.
Now we show by means of an example the existence of $`(X,L)`$, where $`X`$ is a non-anticanonical rational surface and $`L`$ is an ample and free line bundle such that $`L`$ satisfies $`N_p`$ but not $`N_{p+1}`$ and $`K_XL=p+3`$. Let $`\pi :X𝔽_0`$ be the blowing up of $`𝔽_0`$ at $`9`$ points. We choose the $`9`$ points $`p_1,\mathrm{}p_9`$ so that $`X`$ is not anticanonical. Let $`E_1,\mathrm{},E_9`$ be the exceptional divisors and let $`L=\pi ^{}(2C_0+nf)E_1\mathrm{}E_9`$ with $`n4`$. Then $`L`$ is ample, because $`L=A+\pi ^{}((n3)f)`$, where $`A`$ is as in Example 1.16. $`L`$ is also base-point-free. This can be checked using Reider’s Theorem. Indeed, $`L`$ can be written as $`K_X+2A+\pi ^{}((n4)f)`$. Let $`C`$ be a smooth curve in $`|L|`$. We want to show that $`K_X𝒪_C`$ is effective. We look at the long exact sequence relating the cohomology of $`K_XL`$, $`K_X`$ and $`K_X𝒪_C`$. By our choice of $`p_1,\mathrm{},p_9`$, $`h^0(K_X)=0`$ and by Riemann–Roch $`h^1(K_X)=1`$. Also by Riemann–Roch, $`h^1(K_XL)=n3`$. Then $`h^0(K_X𝒪_C)=n31`$. Now $`K_XL=2n5`$. Then according to Theorem 1.3 $`L`$ satisfies property $`N_{2n8}`$ and by \[GL\] , Theorem 2, $`L`$ does not satisfy property $`N_{2n7}`$.
Remark 1.5 . Theorem 1.3 is an example of the following philosophy: On a variety $`X`$ the failure of a line bundle $`L`$ to satisfy property $`N_p`$ can be traced to the existence of an extremal curve $`C`$ on $`X`$. The curve $`C`$ is extremal with respect to $`L`$ and property $`N_p`$ in the following sense: $`L_C`$ satisfies property $`N_p`$ but not property $`N_{p+1}`$. In the proof of Theorem 1.3 the existence of such an extremal curve, namely a smooth curve $`C`$ in $`|L|`$, plays a key role. We would like to point out the existence of another extremal curve on $`X`$. If there is a smooth irreducible curve $`C^{}`$ in $`|K_X|`$ (for instance, if $`X`$ is a Del Pezzo surface), then $`C^{}`$ is a smooth elliptic curve. On an elliptic curve a line bundle satisfies property $`N_p`$ if and only if its degree is greater than or equal to $`p+3`$. This follows from Green’s theorem (cf. \[G\] , Theorem 4.a.1) and the self-duality of the resolution of any elliptic normal curve or alternatively by the theorem of Green and Lazarsfeld quoted above (cf. \[GL\] , Theorem 2). Now Theorem 1.3 says precisely that $`L`$ satisfies property $`N_p`$ but not property $`N_{p+1}`$ if and only if $`K_XL=p+3`$. This agrees with the philosophy just stated, since $`K_XL`$ is the degree of $`L_C^{}`$ and $`p+3`$ is the degree of those line bundles of $`C^{}`$ satisfying property $`N_p`$ but not property $`N_{p+1}`$.
In the remainder of this section we will focus on anticanonical surfaces. One of our main purposes is to find uniform and optimal bounds on $`n`$ so that a line bundle of the form $`K_X+A_1+\mathrm{}+A_n`$ satisfies property $`N_p`$, where each $`A_i`$ is an ample line bundle on $`X`$. The bounds we will obtain will be for anticanonical surfaces with a fixed value of $`K_X^2`$. They will therefore be finer than a uniform bound valid for any anticanonical surface. There are two ingredients we need to attack the problem. First we need to find a sharp, uniform lower bound on $`n`$ so that $`K_X+A_1+\mathrm{}+A_n`$ be very ample. This will be dealt with in Proposition 1.6. Second, because of the numerical characterization of property $`N_p`$ given by Theorem 1.3 , we need to find a sharp, uniform lower bound for $`K_XA`$ for an ample $`A`$. This is what we do in Proposition 1.9 and Proposition 1.10 . The sharpness of Propositions 1.6, 1.9 and 1.10 is shown by the Examples 1.11 to 1.20 presented later on.
###### Proposition 1.6
Let $`X`$ be an anticanonical surface, let $`A_1,\mathrm{},A_n`$ be ample line bundles on $`X`$, and let $`L=K_X+A_1+\mathrm{}+A_n`$. Then $`L`$ is very ample if
1) $`K_X^2=9`$ and $`n4`$.
2) $`K_X^2=8`$ and $`n3`$.
3) $`3K_X^27`$ and $`n2`$.
4) $`K_X^2=2`$ and $`n3`$ or $`n2`$ unless $`n=2,A_1=A_2=K_X`$.
5) $`K_X^2=1`$ and $`n4`$ or $`n2`$ unless
5a) $`n=2`$ and $`A_1=A_2=K_X`$ or $`A_1=K_X`$ and $`A_2=2K_X`$ (or vice versa); or
5b) $`n=3`$ and $`A_1=A_2=A_3=K_X`$.
6) $`K_X^2=0`$ and $`n3`$.
7) $`K_X^2<0`$ and $`n2`$.
In order to prove the result we need this
###### Lemma 1.7
Let $`X`$ be an anticanonical surface and let $`A`$ be an ample line bundle on $`X`$. Then $`A^22`$ unless $`(X,A)=(^2,𝒪_^2(1))`$ or $`K_X^2=1`$ and $`A=K_X`$.
Proof. If $`K_XA2`$, then by \[Hb1\] , Theorem III.1.a $`A`$ is base-point-free, hence if $`A^2=1`$, then $`(X,A)=(^2,𝒪_^2(1))`$. If $`K_XA=1`$ and $`A^2=1`$, then $`(K_X+A)A=0`$ and hence $`K_X+A`$ is effective by Riemann–Roch. Since $`A`$ is ample, $`A=K_X`$ and $`K_X^2=1`$. $`\mathrm{}`$
(1.8) Proof of Proposition 1.6. First, we note that if $`n4`$, then $`L`$ is very ample by Reider’s Theorem. Now, except if $`(X,A)=(^2,𝒪_^2(1))`$ or $`A=K_X`$ and $`K_X^2=1`$, an ample line bundle $`A`$ satisfies $`A^22`$ by Lemma 1.7. Therefore if $`n3`$, $`L`$ is very ample by Reider’s Theorem unless all $`A_i=A`$ and, either $`(X,A)=(^2,𝒪_^2(1))`$, or $`A=K_X`$ and $`K_X^2=1`$.
Now we prove that under the hypothesis stated in the proposition, $`L`$ is very ample if $`n2`$. First we show that the result is true if $`A_1=A_2=A`$ and $`A^23`$. By Reider’s Theorem the only obstacle to the very ampleness of $`L`$ would be the existence of an irreducible curve $`E`$ with $`AE=1`$ and $`E^2=0`$. Then we prove the following:
(1.8.1) Let $`X`$ be an anticanonical surface with $`K_X^20,8`$ and let $`A`$ be an ample line bundle on $`X`$. Then there does not exist an irreducible curve $`E`$ on $`X`$ with $`AE=1`$ and $`E^2=0`$.
We show that, under the above hypothesis, $`p_a(E)=0`$ or $`1`$. Since $`E`$ is irreducible and $`E^2=0`$, $`E`$ is nef, hence $`K_XE0`$ and $`(K_X+E)E0`$. We exclude both possibilities, starting with $`p_a(E)=1`$. In this case $`K_XA=1`$, otherwise $`A`$ would be base-point-free by \[Hb1\] , Theorem III.1 and this would contradict the fact that $`AE=1`$. Now $`K_X+E`$ is effective by Riemann–Roch. Indeed, $`(K_X+E)E=0`$ and $`E`$ is a curve, hence $`h^0(K_X+E)=1+h^1(K_X+E)1`$. On the other hand $`(K_X+E)A=0`$, hence $`E=K_X`$ and $`K_X^2=0`$, which is excluded by hypothesis. Now we rule out the second possibility, namely, $`p_a(E)=0`$. In this case $`E=^1`$ and it follows from Lemma 4.1.10 in \[BS\] that $`X`$ is a $`^1`$-bundle, which is also excluded by hypothesis.
By Lemma 1.7 the only case left when $`A_1=A_2=A`$ is $`A^2=2`$. Then $`K_XA2`$ and $`A`$ is base-point-free by \[Hb1\] , Theorem III.1.a. Then $`X`$ is either $`^1\times ^1`$ or a double cover of $`^2`$ ramified along a smooth quartic, for which $`A=K_X`$ and $`K_X^2=2`$. Both possibilities are excluded by hypothesis.
Now we deal with the case $`A_1A_2`$. If one $`A_i`$ is such that $`A_i^25`$, by (1.8.1) and Reider’s Theorem $`K_X+A_1+A_2`$ is very ample. If, say, $`A_1^2=4`$ and $`A_2^22`$, we prove the following:
(1.8.2) Let $`X`$ be an anticanonical surface and let $`A_1`$ and $`A_2`$ be ample line bundles on $`X`$ such that $`A_1^23`$. Then $`A_1A_22`$.
Then we are done by (1.8.1), (1.8.2) and Reider’s Theorem. Therefore we prove (1.8.2). Assume the contrary. Now for any ample bundle $`A`$ on $`X`$, and in particular for $`A_2`$, we have that $`h^2(A)=h^0(K_XA)=0`$, because any ample bundle intersects $`K_XA`$ strictly negatively. Then by Riemann–Roch $`A_2`$ is effective. On the other hand $`(K_X+2A_1)A_21`$. We have seen before that $`K_X+2A_1`$ is very ample, therefore $`(K_X+2A_1)A_2=1`$, $`K_XA_2=1`$ and each member of $`|A_2|`$ is a smooth $`^1`$. But if $`K_XA_2=1`$ then $`(K_X+A_2)A_20`$ and this contradicts the fact that the sectional genus of $`A_2`$ is $`0`$.
Now if, say, $`A_1^2=4`$ and $`A_2^2=1`$, again we know by (1.8.2) that $`A_1A_22`$, and if $`A_1A_23`$ we are again done by (1.8.1) and Reider’s Theorem. Thus we consider the case $`A_1^2=4`$, $`A_2^2=1`$ and $`A_1A_2=2`$. Since $`n=2`$ and $`A_2^2=1`$, by Lemma 1.7 $`A_2=K_X`$ and $`K_X^2=1`$. Then $`A_1(2K_X+A_1)=0`$. On the other hand $`h^2(2K_X+A_1)=h^0(K_XA_1)=0`$, because $`A_1(K_XA_1)=2`$. Then it follows by Riemann-Roch that $`2K_X+A_1`$ is effective, so $`A_1=2K_X`$. However the possibility $`K_X^2=1`$, $`A_2=K_X`$ and $`A_1=2K_X`$ is excluded by hypothesis.
Now the only remaining cases are $`A_1^2=3`$ and $`A_2^2=1,2,3`$, $`A_1^2=2`$ and $`A_2^2=2,1`$ and $`A_1^2=A_2^2=1`$. We analyze them case by case. If $`A_1^2=A_2^2=3`$, then by (1.8.2) $`A_1A_22`$, then by (1.8.1) and Reider’s Theorem we are done. If $`A_1^2=3`$ and $`A_2^2=2`$, then $`K_XA_22`$ and as already argued $`A_2=K_X`$ and $`K_X^2=2`$. If $`A_1^2=3`$ and $`A_2^2=1`$, then by Lemma 1.7, $`A_2=K_X`$ and $`K_X^2=1`$. But then in neither of the cases can it happen that $`A_1^2=3`$ and $`K_XA_1=2`$. If $`A_1^2=2`$ and $`A_2^2=2`$ or $`1`$, then either $`A_1=A_2=K_X`$ and $`K_X^2=2`$, which is excluded because we are assuming $`A_1A_2`$, or $`A_1=K_X`$ and $`K_X^2=2`$ and $`A_2=K_X`$ and $`K_X^2=1`$, which is absurd. Finally if $`A_1^2=A_2^2=1`$, then by Lemma 1.7, $`A_1=A_2=K_X`$ and $`K_X^2=1`$, but we are assuming $`A_1A_2`$. $`\mathrm{}`$
We now proceed to study lower bounds for $`K_XA`$. One lower bound is of course $`1`$. This bound is sharp if $`K_X^21`$, as proven by Examples 1.15, 1.18, 1.19 and 1.20. However, if $`K_X^22`$ a better bound holds: First, if $`K_X^2=9`$, then $`X=^2`$ and the bound is $`3`$. If $`K_X^2=8`$, $`X`$ is a Hirzebruch surface and we summarize the bounds on $`K_XA`$ in the next proposition. Finally if $`1K_X^27`$, $`K_XAK_X^2`$. This bound is sharp as shown in Examples 1.13 and 1.14. We end this study by classifying the boundary and near-boundary cases when $`1K_X^27`$ and by obtaining a better bound when $`(X,A)`$ does not fall in one of these boundary and near-boundary cases. As explained before, we need these technical results to prove Theorem 1.23 .
###### Proposition 1.9
Let $`X`$ be an anticanonical surface and let $`A`$ be an ample line bundle on $`X`$. If $`K_X^2=8`$, i.e., if $`X`$ is a (geometrically) ruled rational surface, then $`K_XA4`$. More precisely, if $`X=𝔽_e`$, then $`K_XAe+4`$.
Proof. If $`X=𝔽_e`$, $`K_X`$ is linearly equivalent to $`2C_0+(e+2)f`$, where $`C_0`$ is the minimal section of $`𝔽_e`$ and $`f`$ is a fiber. In any case there is a divisor in the anticanonical class having $`e+4`$ irreducible components counted with multiplicity. Then the intersection number of any ample line bundle $`A`$ with this divisor is greater than or equal to $`e+4`$. $`\mathrm{}`$
###### \KA
Let $`X`$ be an anticanonical surface and let $`A`$ be an ample line bundle such that $`1K_X^27`$.
Then $`K_XAK_X^2+3`$ unless one of the following happens:
a) $`A=K_X`$, in which case $`K_XA=K_X^2`$;
b) $`K_X^2=1`$ and $`A=2K_X`$, in which case $`K_XA=K_X^2+1`$;
c) $`K_X^2=1`$ and $`A=3K_X`$, in which case $`K_XA=K_X^2+2`$;
d) $`K_X^2=2`$ and $`A=2K_X`$, in which case $`K_XA=K_X^2+2`$;
e) $`K_X+A`$ is a base-point-free line bundle, $`A`$ is very ample and $`(X,A)`$ is a conic fibration under $`|K_X+A|`$, in which case $`K_XAK_X^2+2`$.
Proof. We assume throughout the proof that $`AK_X`$ and we want to see that, except for the other exceptions listed in the statement, $`K_X(K_X+A)3`$. We divide the proof in two cases:
Case 1: $`K_X^22`$. Assume first that $`A^25`$. Then by Hodge Index Theorem $`K_XA4`$; in particular, $`A`$ is very ample. We apply Reider’s Theorem to $`K_X+A`$ to see it is base-point-free. The only obstruction to $`K_X+A`$ being base-point-free is the existence of a reduced and irreducible curve $`E`$ such that $`AE=1`$ and $`E^2=0`$. Since $`A`$ is very ample, $`E=^1`$. Then it follows from \[BS\] , Lemma 4.1.10 that $`X`$ would be a $`^1`$-bundle, which is excluded by hypothesis. We apply now \[Hb1\] , Lemma II.6.a. If $`K_X+A`$ is composed with a pencil, since $`K_X+A`$ is base-point-free and $`A`$ is ample, $`(K_X+A)^2=0`$, $`h^1(K_X+A)=0`$ and $`K_X(K_X+A)=2r`$, for some $`r1`$. If $`r2`$, we are done. If $`r=1`$, then $`(K_X+A)^2=0`$ and $`K_X(K_X+A)=2`$. Hence $`(K_X+A)A=2`$, therefore the sectional genus of $`A`$ is $`2`$.
Then by \[BS\] , Theorem 10.2.7.2 $`(X,A)`$ is a conic fibration over $`^1`$ under $`|K_X+A|`$. Note that in this case $`K_X(K_X+A)`$ is an even number greater than or equal to $`2`$.
If $`K_X+A`$ is not compose with a pencil then $`(K_X+A)^2>0`$ by \[Hb1\] , Lemma II.6.b. We want to see under what conditions $`K_X(K_X+A)3`$. Since $`K_X^22`$ it follows from Hodge Index Theorem that $`K_X(K_X+A)2`$. Assume then that $`K_X(K_X+A)=2`$. Then $`(K_X+A)^22`$ and even. Now if $`(K_X+A)^24`$ or $`K_X^23`$, we are done by Hodge Index Theorem. Then the only case left is $`(K_X+A)^2=K_X^2=2`$.
On the other hand $`K_X+A`$ is big and nef, therefore by Kawamata-Viehweg $`h^1(2K_X+A)=h^2(2K_X+A)=0`$ and by Riemann–Roch $`2K_X+A`$ is effective. Since $`A`$ is ample and $`A(2K_X+A)=0`$, this implies that $`A=2K_X`$.
We deal now with the case $`A^24`$. By Lemma 1.7, $`A^22`$, and by Hodge Index Theorem $`K_XA2`$, hence by \[Hb1\] , Theorem III.1.a $`A`$ is base point free. If $`A^2=2`$, then either $`X=^1\times ^1`$ or $`K_X^2=2`$ and $`A=K_X`$. The two possibilities are excluded either by hypothesis or by assumption. Let now $`A^2=3`$ or $`4`$.
If $`A^2=3`$, by Hodge Index Theorem $`K_XA3`$, hence by Theorem 1.3 , $`A`$ is very ample. Then $`X`$ is a Del Pezzo cubic surface in $`^3`$ and $`A=K_X`$ or a rational normal scroll in $`^4`$. The two possibilities are excluded by hypothesis or by assumption.
Finally, if $`A^2=4`$, again by Hodge Index Theorem $`K_XA3`$ and $`A`$ is very ample. However the only linearly normal smooth surfaces of degree $`4`$ are K3 surfaces in $`^3`$, the Del Pezzo surface in $`^4`$ and the Veronese Surface and the rational normal scrolls in $`^5`$.
Case 2: $`K_X^2=1`$. Assume first that $`A^25`$. Then by Hodge Index Theorem $`K_XA3`$; in particular, $`A`$ is very ample. The same argument used in Case 1 proves that $`K_X+A`$ is base-point-free. We apply again \[Hb1\] , Lemma II.6.a. Then if $`K_X+A`$ is composed with a pencil, we get as before that either $`K_X(K_X+A)4`$ or, $`K_X(K_X+A)=2`$ and by \[BS\] , Theorem 10.2.7.2 $`(X,A)`$ is a conic fibration over $`^1`$ under $`|K_X+A|`$.
If $`K_X+A`$ is not compose with a pencil then $`(K_X+A)^2>0`$ by \[Hb1\] , Lemma II.6.b. We want to see under what conditions $`K_X(K_X+A)<3`$. If $`(K_X+A)^25`$, then $`K_X(K_X+A)3`$ by Hodge Index Theorem. Then we have to study the cases $`1(K_X+A)^24`$. Then by Hodge Index Theorem $`K_X(K_X+A)1`$.
Assume $`(K_X+A)^2=1`$. Then $`K_X(K_X+A)`$ is odd, so we only have to worry about $`K_X(K_X+A)=1`$. Recall that $`K_X+A`$ is free and big, therefore by Kawamata–Viehweg Theorem, $`h^1(2K_X+A)=h^2(2K_X+A)=0`$. Since $`(2K_X+A)(K_X+A)=0`$, by Riemann–Roch it follows that $`h^0(2K_X+A)=1`$. On the other hand, since $`K_X^2=(K_X+A)^2=1`$, then $`A(2K_X+A)=0`$, hence $`A=2K_X`$. This is excluded because at this point we are assuming $`A^25`$.
Assume $`(K_X+A)^2=2`$. Then $`K_X(K_X+A)`$ is even, so we only have to worry about $`K_X(K_X+A)=2`$. Note that $`A(K_X+A)=4`$. Then by Riemann-Roch $`|K_X+A|`$ maps $`X`$ as double cover of $`^2`$. Then either $`K_X^2=8`$, which is excluded by hypothesis or $`K_X^2=2`$, which is excluded by assumption.
Assume $`(K_X+A)^2=3`$. Then $`K_X(K_X+A)`$ is odd, so we only have to worry about $`K_X(K_X+A)=1`$. This cannot happen by Hodge Index Theorem.
Assume $`(K_X+A)^2=4`$. Then $`K_X(K_X+A)`$ is even, so we only have to worry about $`K_X(K_X+A)=2`$. First we see that $`(2K_X+A)`$ is effective. Indeed, $`(2K_X+A)(K_X+A)=2`$, then by Kawamata–Viehweg Theorem and Riemann–Roch, $`h^0(2K_X+A)=2`$. Then $`h^2(3K_X+A)=h^0(2K_XA)=0`$. We also have that $`(3K_X+A)(2K_X+A)=0`$. Then by Riemann–Roch $`3K_X+A`$ is effective. On the other hand $`A(3K_X+A)0`$, hence $`A(3K_X+A)=0`$ and $`A=3K_X`$.
We deal now with the case $`A^24`$. By Lemma 1.7, if $`A^2=1`$, then $`A=K_X`$, which is excluded by assumption. If $`A^22`$, by Hodge Index Theorem $`K_XA2`$, hence by \[Hb1\] , Theorem III.1.a $`A`$ is base point free. If $`A^2=2`$, then either $`X=^1\times ^1`$, which is excluded by hypothesis or $`K_X^2=2`$ and $`A=K_X`$, which is excluded by assumption.
If $`A^2=3`$, by Hodge Index Theorem and because $`K_XA`$ must be odd $`K_XA3`$, hence by Theorem 1.3 , $`A`$ is very ample. Then $`X`$ is a Del Pezzo cubic surface in $`^3`$ and $`A=K_X`$ or a rational normal scroll in $`^4`$. Both possibilities are excluded either by assumption or by hypothesis.
Finally let $`A^2=4`$. Again by Hodge Index Theorem $`K_XA2`$ and $`A`$ is base-point-free. If $`K_XA3`$, $`A`$ is actually very ample by Theorem 1.3 . However the only linearly normal smooth surfaces of degree $`4`$ are K3 surfaces in $`^3`$, the Del Pezzo surface in $`^4`$ and the Veronese Surface and the rational normal scrolls in $`^5`$. Then the only possibility left to study is $`A^2=4`$ and $`K_XA=2`$. Then $`A(2K_X+A)=0`$. On the other hand $`h^2(2K_X+A)=h^0(K_XA)=0`$, because $`A(K_XA)=2`$. Then it follows by Riemann-Roch that $`2K_XA_1`$ is effective, so $`A=2K_X`$. $`\mathrm{}`$
Before stating and proving Theorems 1.23 and 1.24, which are consequences of Propositions 1.6, 1.9 and 1.10 it is important to know that these propositions are sharp. We do this by means of the following series of examples. Then these examples will also imply the optimality of our theorems on rational surfaces.
###### Example 1.11
$`K_X^2=9`$, $`A`$ ample, $`K_XA=3`$, $`K_X+3A`$ not ample.
Of course the only rational surface with $`K_X^2=9`$ is $`X=^2`$, $`A=𝒪_^2(1)`$ attains the bound $`K_XA3`$ and $`K_X+3A=𝒪_^2`$, hence not very ample.
###### Example 1.12
$`K_X^2=8`$, $`A`$ ample, $`K_XA=e+4`$, $`K_X+2A`$ not ample.
The rational surfaces with $`K_X^2=8`$ are the Hirzebruch surfaces $`𝔽_e`$. With these examples we show that the bounds computed in Proposition 1.9 for $`K_XA`$ in terms of $`e`$ are sharp. If $`X=𝔽_e`$, let $`C_0`$ be the minimal section and $`f`$ a fiber. The divisor $`A=C_0+(e+1)f`$ is ample and $`K_XA`$ achieves the bound $`e+4`$. Moreover $`K_X+2A=ef`$ is free but not ample, hence $`A`$ provides an example of $`K_X+2A`$ not being very ample.
###### Example 1.13
$`3K_X^27`$, $`A`$ ample, $`K_XA=K_X^2`$, $`K_X+A`$ not ample.
We obtain $`X`$ with $`K_X^2=9i`$ by blowing up $`^2`$ at $`2i6`$ points, which we choose in sufficiently general position (not $`3`$ of them on a line, not $`6`$ of them on a conic). This is a Del Pezzo surface and $`K_X`$ is very ample (cf. \[Ht\] , Theorem V.4.6). Then $`(X,K_X)`$ achieves the bound $`A^2K_X^2`$ and also provides an example where $`K_X+A`$ is not very ample.
###### Example 1.14
$`K_X^2=2`$, $`A`$ ample, $`K_XA=K_X^2`$ and $`K_X+2A`$ not very ample.
Let $`p:X^2`$ be the double cover of $`^2`$ ramified along a smooth quartic. Such an $`X`$ is rational, $`K_X=p^{}(𝒪_^2(1))`$ and $`K_X=p^{}(𝒪_^2(1))`$. Therefore $`K_X^2=2`$ and, since $`p`$ is finite, $`K_X`$ is ample and base-point-free. In fact, $`p`$ is induced by the morphism of the complete anticanonical linear series. We remark that, actually, a rational surface $`X`$ with $`K_X^2=2`$ and with anticanonical divisor ample is a double cover of $`^2`$ as the one described above. This follows from Riemann–Roch. Thus the anticanonical divisor on such a surface $`X`$ attains the bound $`A^2K_X^2`$.
Such a surface $`X`$ can also be found embedded in $`^6`$. Indeed, $`2K_X`$ is very ample, and embeds $`X`$ as a degree $`8`$, sectional genus $`6`$, smooth surface in $`^6`$ (see \[BS\] , Example 10.2.4).
In addition $`A=K_X`$ provides an example of $`K_X+2A`$ not being very ample.
###### Example 1.15
$`K_X^2=1`$, $`A`$ ample, $`K_XA=1`$, $`K_X+3A`$ not very ample.
An example of a surface $`X`$ with $`K_X^2=1`$ and $`K_X`$ ample can be found in \[BS\] , Example 10.4.3. Indeed there exists a smooth rational surface with $`K_X^2=1`$ and $`K_X`$ ample. This surface is embedded by $`|3K_X|`$ as a degree $`9`$, sectional genus $`4`$ smooth surface in $`^6`$. Then $`(X,K_X)`$ achieves the bound $`A^2K_X^2`$. In addition $`A=K_X`$ provides an example of $`K_X+3A`$ not being very ample. Indeed, $`2K_X`$ is not very ample as its complete linear series induces a double cover of quadric cone in $`^3`$.
###### Example 1.16
$`1K_X^28`$, $`A`$ ample, $`K_XA=K_X^2+2`$, $`K_X+A`$ not ample.
The family of examples we consider now are conic bundles. Let us fix $`n=K_X^2`$, $`1n8`$. The examples are constructed from the pair $`(Y=𝔽_e,2C_0+mf)`$, where $`m=e+3`$, $`0e2`$, $`C_0`$ is the minimal section of $`Y`$ and $`f`$ is a fiber of $`Y`$. Let $`C`$ be a smooth anticanonical divisor on $`Y`$. Let $`l=8n`$. We choose $`\mathrm{\Sigma }=\{p_1,\mathrm{},p_l\}`$, $`l`$ distinct points on $`C`$ lying on different fibers of $`Y`$. Let $`\pi :XY`$ be the blowing-up of $`Y`$ along $`\mathrm{\Sigma }`$ and let $`E_1,\mathrm{},E_l`$ be the exceptional divisors. Let $`A=\pi ^{}(2C_0+mf)E_1\mathrm{}E_l`$. By the choice of $`m`$, $`K_X+A=\pi ^{}f`$. Hence $`K_X+A`$ is not ample and $`K_XA=K_X^2+2`$. We see now that $`A`$ is ample using Nakai–Moishezon’s criterion. Firstly, as $`n1`$, $`A^23`$. Secondly, it is clear that $`AE_i`$=1, for all $`1il`$. Finally we will check that the intersection of $`A`$ with any irreducible non-exceptional curve $`T`$ on $`X`$ is strictly positive. Let $`D=\pi (T)`$ and let $`DaC_0+bf`$. Let $`m_1,\mathrm{},m_l`$ be the multiplicities of $`D`$ at $`p_1,\mathrm{},p_l`$. Then $`AT=(2C_0+mf)Dm_1\mathrm{}m_l`$. Hence we want $`m_1+\mathrm{}+m_l<(2C_0+mf)D=3a+2bae`$. We first consider the case when $`C`$ and $`D`$ intersect properly. Since $`p_1,\mathrm{},p_lC`$, $`m_1+\mathrm{}+m_lCD`$. Therefore $`m_1+\mathrm{}+m_l2a+2bae`$. Then if $`a>0`$ we are done. If $`a=0`$, then $`D=f`$. Since we have chosen $`p_1,\mathrm{},p_l`$ in different fibers in this case $`m_1+\mathrm{}+m_l=1<2=3a+2bae`$, and we are also done. Now consider the case when $`C`$ and $`D`$ do not intersect properly. Since both $`C`$ and $`D`$ are irreducible, $`C=D`$. In this case $`m_1+\mathrm{}+m_l=l`$, for $`C`$ is smooth, and $`a=2`$ and $`b=e+2`$. Then $`3a+2bae=10`$ and we are done if $`l<10`$. The latter happens because $`n1`$.
###### Example 1.17
$`2K_X^28`$, $`A`$ ample, $`K_XA=K_X^2+3`$ and $`K_X+A`$ not ample.
Let $`Y=𝔽_1`$, let $`C_0`$ the minimal section, let $`f`$ be a fiber and let $`C`$ be a smooth irreducible anticanonical curve. Let $`\pi :XY`$ the blowing up of $`Y`$ at $`\mathrm{\Sigma }=\{p_1,\mathrm{},p_l\}`$, where $`p_1,\mathrm{},p_l`$ are distinct points of $`Y`$ on $`C`$ away from $`C_0`$ and $`0l10`$. Let $`E_1,\mathrm{},E_l`$ be the exceptional divisors lying over $`p_1,\mathrm{},p_l`$ respectively. Let $`A=\pi ^{}(3C_0+4f)E_1E_2\mathrm{}E_l`$. We claim that $`A`$ is ample, that $`K_X(K_X+A)=3`$ and that $`K_X+A`$ is not ample. The latter is clear, since $`K_X+A=\pi ^{}(C_0+f)`$. From this it also follows that $`K_X(K_X+A)=3`$. We check finally the ampleness of A using Nakai-Moishezon’s criterion. On the one hand $`A^2=15l5`$, since $`l10`$. Clearly $`AE_i=1`$. Let $`T`$ be now a nonexceptional irreducible curve on $`X`$, and let $`D=\pi (T)`$. Let $`DaC_0+bf`$ and let $`m_1,\mathrm{},m_l`$ be the multiplicities of $`D`$ at $`p_1,\mathrm{},p_l`$. Then $`AT=(3C_0+4f)Dm_1\mathrm{}m_l=a+3bm_1\mathrm{}m_l`$, so we want $`m_1+m_2+\mathrm{}+m_l<a+3b`$. We distinguish several cases: first we consider the case when $`C`$ and $`D`$ intersect properly. Then $`m_1+\mathrm{}+m_lCD=a+2b`$ since every $`p_i`$ lies on $`C`$. We consider two subcases: $`b>0`$ or $`b=0`$. If the former, $`a+2b<a+3b`$ and we are done. If the latter $`T=C_0`$ and $`m_1=\mathrm{}=m_l=0`$ by our choice of $`p_1,\mathrm{},p_l`$. Since $`0=m_1+\mathrm{}+m_l<1`$ we are also done. The only case left is when $`C`$ and $`D`$ do not intersect properly. Since both $`C`$ and $`D`$ are irreducible, $`C=D`$. Since $`C`$ is smooth, $`m_1+\mathrm{}+m_l=l`$. On the other hand $`a+3b=11`$ in this case, and, since $`l10`$, we are done.
###### Example 1.18
$`K_X^2=0`$, $`A`$ ample, $`K_XA=1`$, $`K_X+2A`$ not very ample.
We find a surface with $`K_X^2=0`$ with an ample line bundle $`A`$ such that $`K_XA=1`$ and $`K_X+2A`$ is not very ample. Let $`Y=^2`$. We consider a set $`\mathrm{\Sigma }Y`$ of $`9`$ distinct points being the complete intersection of $`2`$ cubics, and neither $`3`$ of them on a line nor $`6`$ of them on a conic. Let $`X`$ be the blowing up of $`Y`$ along $`\mathrm{\Sigma }`$. Then $`K_X`$ is base-point-free, $`h^0(K_X)=2`$, and all $`C|K_X|`$ are irreducible curves. Therefore $`|K_X|`$ turns $`Y`$ into an elliptic fibration $`\phi :X^1`$ with irreducible fibers and at least $`9`$ sections, namely, the $`9`$ exceptional divisors. Let $`E`$ be one of the exceptional divisors and let $`F`$ be a fiber of $`\phi `$. Then $`A=E+2F`$ is ample. Indeed, we use Nakai-Moishezon’s criterion. The self-intersection $`(E+2F)^2=E^2+4EF=3`$. If $`F^{}`$ is a fiber of $`\phi `$, $`(E+2F)F^{}=1`$. Let $`C`$ be an irreducible curve on $`X`$, and not a fiber. Then either $`C=E`$, in which case $`(E+2F)C=1`$ or $`(E+2F)C=EC+2FC2`$. Now we see that $`K_XA=1`$ and that $`K_X+2A`$ is not very ample. It is clear that $`K_XA=1`$. We see now that $`K_X+2A`$ is not very ample. Indeed, let $`C|K_X|`$. Then $`C`$ has arithmetic genus $`1`$; however $`K_X+2A=3K_X+2E`$, hence $`(K_X+2A)C=2`$. Since the arithmetic genus of $`C`$ is $`1`$, $`K_X+2A`$ is not very ample. Note that $`A^{}=E+nF`$ for some $`n2`$ satisfies as well that $`K_XA^{}=1`$ and $`K_X+2A^{}`$ is not very ample.
###### Example 1.19
$`K_X^2<0`$ odd, $`A`$ ample, $`K_XA=1`$ and $`K_X+A`$ not ample.
Given $`n`$ odd number strictly smaller than $`0`$, we find a surface $`X`$ with $`K_X^2=n`$ and an ample line bundle $`A`$ on $`X`$ such that $`K_XA=1`$. Let $`Y=𝔽_0`$, let $`l=8n`$. Let $`C`$ be a smooth anticanonical curve on $`Y`$ and $`f_1`$ and $`f_2`$ be two lines each belonging to one ruling of $`𝔽_0`$. Let $`\mathrm{\Sigma }=\{p_1,\mathrm{},p_l\}`$ be $`l`$ distinct points on $`C`$ chosen so that not two of them belong to a line of $`𝔽_0`$. Let $`\pi :XY`$ be the blowing up of $`Y`$ along $`\mathrm{\Sigma }`$ and let $`E_i`$ be the exceptional divisor over $`p_i`$. Let $`k=\frac{l3}{2}`$ and let $`A=\pi ^{}(𝒪_{𝔽_0}(2f_1+kf_2)E_1\mathrm{}E_l)`$. We see that $`K_XA=(2f_1+2f_2)(2f_1+kf_2)+E_1^2+\mathrm{}+E_l^2=1`$. Moreover $`K_X+A=\pi ^{}((k2)f_2)`$, and therefore it is not ample. Now we will see that $`A`$ is ample using Nakai–Moishezon’s criterion. First $`A^2=l63`$. Now we see the intersection of $`A`$ with irreducible curves on $`X`$. The intersection of $`A`$ with each $`E_i`$ is $`1`$. Let $`T`$ be an irreducible curve which is not an exceptional divisor and let $`D=\pi (T)`$. Let $`Daf_1+bf_2`$. Let $`m_1,\mathrm{},m_l`$ be the multiplicities of $`D`$ at $`p_1,\mathrm{},p_l`$. Then $`AT=(2f_1+kf_2)Dm_1\mathrm{}m_l`$. Hence we want to see that $`m_1+\mathrm{}+m_l<ka+2b`$. Since $`p_1,\mathrm{},p_lC`$, $`m_1+\mathrm{}+m_lCD`$ if $`C`$ and $`D`$ intersect properly. Then $`m_1+\mathrm{}+m_l2a+2b`$. Since $`n1`$, $`k3`$, hence if $`a1`$, $`2a+2b<ka+2b`$ and we are done. If $`a=0`$, then $`D=f_2`$ and by our choice of $`p_1,\mathrm{},p_l`$, $`m_1+\mathrm{}+m_l=1<2=ka+2b`$ and we are also done. Now suppose that $`C`$ and $`D`$ do not intersect properly. Since both $`C`$ and $`D`$ are irreducible, $`C=D`$ and $`a=b=2`$. In this case $`m_1+\mathrm{}+m_l=l`$, for $`C`$ is smooth. Since $`k=\frac{l3}{2}`$, $`l<2k+4`$.
###### Example 1.20
$`K_X^2<0`$ even, $`A`$ ample, $`K_XA=1`$ and $`K_X+A`$ not ample.
Given $`n`$ even number strictly smaller than $`0`$, we find a surface $`X`$ with $`K_X^2=n`$ and an ample line bundle $`A`$ on $`X`$ such that $`K_XA=1`$. Let $`Y=𝔽_0`$ and let $`l=8n`$. Let $`C`$ be a smooth anticanonical curve on $`Y`$ and $`f_1`$ and $`f_2`$ be two lines each belonging to one ruling of $`𝔽_0`$. Let $`\mathrm{\Sigma }=\{p_1,\mathrm{},p_l\}`$ be $`l`$ distinct points on $`C`$ chosen so that not two of them belong to a line of $`𝔽_0`$. Let $`\pi :XY`$ be the blowing up of $`Y`$ along $`\mathrm{\Sigma }`$ and let $`E_i`$ be the exceptional divisor over $`p_i`$. Let $`k=\frac{l4}{2}`$ and let $`A=\pi ^{}(𝒪_{𝔽_0}(3f_1+kf_2)2E_1E_2\mathrm{}E_l`$. Then $`K_XA=(2f_1+2f_2)(3f_1+kf_2)+2E_1^2+E_2^2+\mathrm{}+E_l^2=1`$. In addition $`K_X+A=\pi ^{}(f_1+(k2)f_2)E_1`$. Therefore $`K_X+A`$ is not ample, since its intersection with the strict transform of one of the lines passing through $`p_1`$ is $`0`$. We see now that $`A`$ is ample using Nakai–Moishezon’s criterion. First $`A^2=2l155`$. Now we see the intersection of $`A`$ with the irreducible curves on $`X`$. The intersection of $`A`$ with $`E_1`$ is $`2`$ and with the other exceptional divisors is $`1`$. Let $`T`$ be an irreducible curve which is not an exceptional divisor and let $`D=\pi (T)`$. Let $`Daf_1+bf_2`$. Let $`m_1,\mathrm{},m_l`$ the multiplicities of $`D`$ at $`p_1,\mathrm{},p_l`$. Then $`AT=(3f_1+kf_2)D2m_1m_2\mathrm{}m_l`$. Hence we want to see that $`2m_1+m_2+\mathrm{}+m_l<ka+3b`$. First we suppose that $`C`$ and $`D`$ intersect properly. Since $`p_1,\mathrm{},p_lC`$, then $`m_1+\mathrm{}+m_lCD=2a+2b`$. We distinguish now two cases. First, if $`m_1a`$, then $`2m_1+m_2+\mathrm{}+m_l3a+2b`$. Since $`n2`$, $`k3`$. Then if $`b1`$, $`3a+2b<ka+3b`$ and we are done. If $`b=0`$, then $`Df_1`$ and $`a=1`$. By the choice of $`p_1,\mathrm{},p_l`$, $`2m_1+m_2+\mathrm{}+m_l2<3ka+3b`$. Second, if $`m_1>a`$, then $`D`$ passes through $`p_1`$ and $`Df_2`$. Then, because of the choice of $`p_1,\mathrm{},p_l`$, $`2m_1+m_2+\mathrm{}+m_l2<3=ka+3b`$. Now we suppose that $`C`$ and $`D`$ do not intersect properly. Since $`C`$ and $`D`$ are irreducible, then $`C=D`$, and $`a=b=2`$. In this case $`2m_1+m_2+\mathrm{}+m_l=l+1`$, for $`C`$ is smooth. Then we are done if $`l+1<2k+6=ka+3b`$. This occurs because $`k=\frac{l4}{2}`$.
The lower bounds for $`K_XA`$ obtained in Proposition 1.9 and Proposition 1.10 combined with Theorem 1.3 yield several results. The following generalizes a well known result, namely, the equivalence of the notion of ampleness and very ampleness for $`^2`$ and for Hirzebruch surfaces (cf. \[Ht\] , Corollary V.2.18). The result we get is actually stronger than this classical result since we obtain the equivalence of ampleness and certain $`N_p`$ property for certain anticanonical surfaces:
###### \ava
Let $`X`$ be a rational surface such that $`d=K_X^23`$. Let $`A`$ be a line bundle on $`X`$. The following are equivalent:
1) $`A`$ is ample;
2) $`A`$ is very ample;
3) $`A`$ satisfies property $`N_0`$.
More precisely:
– If $`3K_X^27`$, then $`A`$ satisfies property $`N_{d3}`$ if and only if $`A`$ is ample and $`A`$ satisfies property $`N_{d1}`$ if and only if $`A`$ is ample different from $`K_X`$ ($`K_X`$ satisfies $`N_{d3}`$ but not $`N_{d2}`$).
– If $`X=𝔽_e`$, then $`A`$ satisfies property $`N_{e+1}`$ if and only if $`A`$ is ample.
– If $`X=^2`$, then $`A`$ satisfies property $`N_0`$ if and only if $`A`$ is ample.
In addition, if $`K_X^2=2`$, then $`A`$ satisfies property $`N_1`$ if and only if $`A`$ is an ample line bundle different form $`K_X`$.
In the next result we show that very ampleness and projective normality are equivalent for anticanonical surfaces:
###### \vapn
Let $`X`$ be an anticanonical surface. A line bundle on $`X`$ is very ample if and only if it satisfies property $`N_0`$.
Proof. Let $`L`$ be very ample and let $`C`$ be smooth curve in $`|L|`$. Since $`K_X𝒪_C`$ is effective, $`K_XL3`$, otherwise $`L𝒪_C`$ would not be very ample. Then $`L`$ satisfies property $`N_0`$ by Theorem 1.3 . $`\mathrm{}`$
We state and show now the result already announced dealing with line bundles of the form $`K_X+A_1+\mathrm{}+A_n`$, with $`A_1,\mathrm{},A_n`$ ample. It follows from Theorem 1.3 and Propositions 1.6, 1.9 and 1.10.
###### \Mukai
Let $`X`$ be an anticanonical surface. Let $`A_1,\mathrm{},A_n`$ be ample line bundles on $`X`$. If $`np+4`$, then $`L=K_X+A_1+\mathrm{}+A_n`$ satisfies property $`N_p`$. This bound is achieved for $`A_i=A`$ and $`(X,A)`$ as in Example 1.15.
More precisely:
1) If $`X=^2`$ and $`n\frac{p}{3}+4`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for $`A_i=𝒪_^2(1)`$ and $`p`$ multiple of $`3`$.
2) If $`K_X^2=8`$ and $`n\frac{p+3}{4}+2`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for $`X=𝔽_0`$, $`A_i=C_0+f`$ and $`p1(4)`$. More precisely, if $`X=𝔽_e`$ and $`n\text{ max}(3,\frac{p+11}{e+4})`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for instance for $`A_i=C_0+(e+1)f`$ and for $`p11(e+4)`$.
3) If $`1K_X^27`$ and $`n\frac{p+3}{K^2}+1`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for instance for $`(X,A_i)`$ as in Examples 1.13, 1.14 and 1.15, and for $`p3(K_X^2)`$.
4) If $`1K_X^27`$, $`A_iK_X`$, $`A_i2K_X`$ if $`K_X^2=1`$ and $`n\frac{p+K_X^2+3}{K_X^2+2}`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for $`(X,A_i)`$ as in Example 1.16, and for $`p+K_X^23(K_X^2+2)`$.
5) If $`2K_X^27`$, $`A_iK_X`$, $`A_i2K_X`$ when $`K_X^2=2`$, $`A_i2K_X,3K_X`$ when $`K_X^2=1`$, it does not happen that $`K_X+A_i`$ is base-point-free, $`A_i`$ is very ample and $`(X,A_i)`$ is a conic fibration under $`|K_X+A_i|`$, and $`n\text{ max}(2,\frac{p+K_X^2+3}{K_X^2+3})`$, then $`L`$ satisfies property $`N_p`$. The bound is achieved for $`(X,A_i)`$ as in Example 1.17, and for $`p+K_X^23(K_X^2+3)`$.
6) If $`1K_X^21`$ and $`np+3+K_X^2`$, then $`L`$ satisfies property $`N_p`$. The bound is sharp for $`(X,A_i)`$ as in Examples 1.15, 1.18 and 1.19.
7) If $`K_X^22`$ and $`n\text{ max}(2,p+3+K_X^2)`$, then $`L`$ satisfies property $`N_p`$. The bound is sharp for $`(X,A_i)`$ as in Examples 1.19 and 1.20.
The next result we will prove is Theorem 1.24 , which is an $`N_p`$ result in the same flavor of Reider’s theorem for base-point-freeness and very ampleness. Theorem 1.24 shows that when $`K_X^21`$, a high self-intersection number for $`L`$ implies by itself property $`N_p`$ for $`K_X+L`$ for a large value of $`p`$. This behavior is in contrast with the behavior that can be observed in surfaces of Kodaira dimension $`0`$.
###### \ReiderNp
Let $`X`$ be a rational surface. Let $`L`$ be a line bundle such that
1) $`LC3`$ for any curve $`C`$ on $`X`$ and $`L^210`$ or
1’) $`K_X+L`$ is very ample.
If $`K_X^21`$:
2a) Let $`L^2(p+3)^2+1`$.
If $`K_X^21`$ and $`L`$ is not a multiple of $`K_X`$ when $`K_X^2=1`$:
2b) Let $`L^2(p+3)^21`$.
If $`K_X^20`$:
2c) Let $`K_XLp+3`$.
Then $`K_X+L`$ satisfies property $`N_p`$.
To prove Theorem 1.24 we will use Theorem 1.3 , Proposition 1.22 and the technical lemma 1.25. The lemma is proven by contradiction and the argument connects property $`N_p`$ with termination of adjunction. In particular it shows that a failure of certain property $`N_p`$ to hold results in non termination of adjunction.
###### \term\ns
Let $`X`$ be a rational surface with $`K_X^2>0`$ and not isomorphic to $`^2`$. Let $`L`$ be a line bundle on $`X`$ such that $`K_X+L`$ is effective. Let $`p1`$ if $`K_X^27`$ and $`p2`$ if $`K_X^2=8`$. If $`L^2(p+3)^21`$ and $`L`$ is not a positive multiple of $`K_X`$ when $`K_X^2=1`$, then $`K_XLp+3+K_X^2`$.
Proof. Assume that, $`L`$ is not a positive multiple of $`K_X`$, or $`K_X^22`$. We will prove the result by way of contradiction. Assuming that $`K_X(K_X+L)p+2`$, we will prove by induction that $`mK_X+L`$ is effective for all $`m2`$. This contradicts the termination of adjunction on a surface of Kodaira dimension $`\mathrm{}`$.
If $`m=1`$, $`K_X+L`$ is effective by hypothesis. Let us see now that $`2K_X+L`$ is effective. Note that $`LK_X`$ because $`L^215`$ and $`K_X^28`$. Then, since $`K_X+L`$ is effective, $`K_XL`$ is not effective. Then by Riemann-Roch
$$h^0(2K_X+L)\frac{1}{2}(2K_X+L)(K_X+L)+1.$$
Therefore we need to see that $`(2K_X+L)(K_X+L)>2`$. This inequality is equivalent to
$$L^2+2K_X^2+13(K_XL).$$
Since by assumption $`K_XLp+2+K_X^2`$, it suffices to check that $`L^23(p+2)+K_X^21`$. Then it is enough to see that $`(p+3)^23p+6+K_X^2`$. This last inequality is equivalent to $`p^2+3p+3K_X^20`$, which holds for all $`p1`$ if $`K_X^27`$ and for all $`p2`$ if $`K_X^2=8`$.
Let $`m2`$. Now we assume $`mK_X+L`$ is effective, and we will show $`(m+1)K_X+L`$ is effective. First consider the case when $`K_X^2=1`$ and $`L`$ is not multiple of $`K_X`$. Since $`mK_X+L`$ is effective, $`h^0(mK_XL)=0`$. Then by Riemann–Roch
$$h^0((m+1)K_X+L)\frac{1}{2}((m+1)K_X+L)(mK_X+L)+1.$$
Thus we need to see that $`((m+1)K_X+L)(mK_X+L)1`$. This is equivalent to
$$L^2+m(m+1)K_X^2+1(2m+1)(K_XL).$$
Since by assumption $`K_XLp+2+K_X^2`$, it suffices to check that
$$L^2(2m+1)(p+2)+(m^2+m+1)K_X^21.$$
$`(\mathrm{1.25.1})`$
Recall that $`K_X^2=1`$. Then, in order to check (1.25.1) it is enough to show that
$$(p+3)^21(2m+1)(p+2)m^2+m.$$
This last inequality is equivalent to
$$(pm)^2+5(pm)+60.$$
This holds for all integers $`p`$ and $`m`$.
Now we assume that $`K_X^22`$. Since $`mK_X+L`$ is effective, $`h^0(mK_XL)1`$. Then by Riemann–Roch
$$h^0((m+1)K_X+L)\frac{1}{2}((m+1)K_X+L)(mK_X+L).$$
Thus we need to see that $`((m+1)K_X+L)(mK_X+L)>0`$. This is equivalent to
$$L^2+m(m+1)K_X^2>(2m+1)(K_XL).$$
Since by assumption $`K_XLp+2+K_X^2`$, it suffices to check that
$$L^2>(2m+1)(p+2)+(m^2+m+1)K_X^2.$$
$`(\mathrm{1.25.2})`$
Since $`m2`$, $`m^2+m+1<0`$. Then, in order to check (1.25.2) it is enough to show that
$$(p+3)^21>(2m+1)(p+2)2(m^2m1).$$
This last inequality is equivalent to
$$p^2+(52m)p+(2m^26m+4)>0.$$
This holds for any integer $`p1`$ and any $`m2`$.
Summing up, we have just shown that, under the hypothesis of the proposition and with the additional assumption that $`K_XLp+2+K_X^2`$, $`mK_X+L`$ is effective for all $`m2`$. As pointed out before, this is a contradiction. Therefore $`K_XLp+3+K_X^2`$, as wished. $`\mathrm{}`$
Remark 1.26. Lemma 1.25 follows from Hodge Index Theorem for some values of $`p`$ and $`K_X^2`$ but not for all.
(1.27) Proof of Theorem 1.24 . By hypothesis 1’) or by 1) and Reider’s Theorem it follows that $`K_X+L`$ is very ample. Therefore by Proposition 1.22 $`K_X+L`$ satisfies property $`N_0`$. If $`X=^2`$, it follows clearly from 2b) that $`K_X(K_X+L)p+3`$ if $`p2`$. Then by Theorem 1.3 $`K_X+L`$ satisfies property $`N_p`$ (if $`L`$ is such that $`L^2=16`$, then one could say that $`K_X+L`$ satisfies $`N_{\mathrm{}}`$). We consider now $`X=𝔽_e`$. If $`p2`$, by Lemma 1.25 , 2b) and Theorem 1.3 $`L`$ satisfies property $`N_p`$. If $`p=1`$ the result also follows because if $`K_X+L`$ is very ample it can be seen that $`L^218`$. Now if $`K_X^2=1`$ and $`L=m(K_X)`$, then 2a) implies $`mp+4`$. Therefore it follows by Theorem 1.3 that $`K_X+L`$ satisfies $`N_p`$. If $`L`$ is not a multiple of $`K_X`$ or if $`K_X^22`$ and $`X^2`$, then it follows from 2b) and Lemma 1.25 that $`K_X(K_X+L)p+3`$. Thus $`K_X+L`$ satisfies property $`N_p`$ by Theorem 1.3 . Finally if $`K_X^20`$, by 2c) $`K_X(K_X+L)p+3`$. Thus $`K_X+L`$ satisfies property $`N_p`$ by Theorem 1.3 . $`\mathrm{}`$
Remark 1.28 Theorem 1.24 is optimal. To see the optimality when 2a) is assumed, take $`X`$ as in Example 1.15 and $`L=(p+4)(K_X)`$. To see the optimality when 1’) and 2b) are assumed take $`A`$ as in Example 1.16, $`e=0`$, $`l=6`$ and $`L=K_X+A`$. Then $`K_X+L`$ satisfies property $`N_1`$ but not $`N_2`$ by Theorem 1.3 and $`L^2=16`$. On the other hand if $`K_X^20`$, assuming only hypothesis 1) and 2a) or 1’) and 2a) do not suffice. This can be seen taking $`(X,A)`$ as in Example 1.18 and $`L=3A`$. Indeed $`K_X+L`$ satisfies property $`N_0`$ but not $`N_1`$; however, $`L^2=27>(2+3)^2+1`$.
We end this section by showing the relation between the property $`N_p`$ satisfied by a line bundle $`L`$ and the termination of ampleness for $`mK_X+L`$.
###### Theorem 1.29
Let $`X`$ be an anticanonical surface and let $`L`$ be a line bundle satisfying property $`N_p`$ but not property $`N_{p+1}`$.
a) If $`X=^2`$ and $`m>\frac{p}{K_X^2}`$;
b) if $`X=𝔽_e`$ and $`m>\frac{pe1}{K_X^2}`$;
c) if $`1K_X^27`$, and $`m>\frac{p+3}{K_X^2}1`$;
d) if $`1K_X^27`$, $`L`$ is not a multiple of $`K_X`$ and $`m>\frac{p+1}{K_X^2}1`$;
e) if $`K_X^2<0`$, and $`m<\frac{p+2}{K_X^2}`$,
then $`mK_X+L`$ is not ample.
Proof. We outline the proof of c), d), e); a) and b) are similar. Let $`1K_X^27`$ and assume $`mK_X+L`$ is ample. Then $`K_X(mK_X+L)K_X^2`$ by Proposition 1.10. This implies $`K_XL(m+1)K_X^2`$. Now if $`L`$ satisfies $`N_p`$ but not $`N_{p+1}`$, $`K_XL=p+3`$. Then $`(m+1)K_X^2p+3`$, and $`m\frac{p+3}{K_X^2}1`$. If $`L`$ is not a multiple of $`K_X`$, then $`K_X(mK_X+L)K_X^2+2`$ and we argue similarly.
Let now $`K_X^2<0`$ and assume again that $`mK_X+L`$ is ample. Then $`K_X(mK_X+L)1`$. This implies $`K_XLmK_X^2+1`$. Again, if $`L`$ satisfies $`N_p`$ but not $`N_{p+1}`$, $`K_XL=p+3`$. Then $`mK_X^2+1p+3`$, and $`m\frac{p+2}{K_X^2}`$. $`\mathrm{}`$
Remark 1.30. The bounds for $`m`$ in the previous theorem are sharp. That is clear for the bound in a). For the bound in b) take $`L=C_0+(e+1)f`$. For the bound in c) take $`L=n(K_X)`$. For the bound in d) take $`L=A`$, where $`A`$ is as in Example 1.16.
## 2. Fano $`n`$-folds of index greater than or equal to $`n1`$
In this section and in the next we prove results on syzygies of certain Fano $`n`$-folds. The first attempt one would try to make to tackle this problem is to imitate the arguments we carried on in Section 1. Given a Fano $`n`$-fold $`X`$ and a very ample line bundle $`L`$, $`H^1(rL)=0`$ for all $`r0`$. Then one could try to obtain information about the free resolution of the image of $`X`$ from whatever information is available on the free resolution of its general hyperplane section $`X^{}`$, which is now of dimension $`n1`$. If no information is readily available on the resolution of $`X^{}`$, one would iterate the argument, and taking successive hyperplane sections one could read the Betti numbers of the resolution of $`X`$ from the resolution of a surface or even, of a curve. This worked very well in the case of rational surfaces because we ended with a curve $`C`$ and a line bundle $`L_C`$ on $`C`$ of relatively high degree, and because of this high degree, we knew relevant information about the resolution of the embedding of $`C`$ by $`|L_C|`$, thanks to the results of Green and Lazarsfeld. However when the dimension is higher we lose control of the hyperplane sections of $`X`$, or rather, there is much less information available on the syzygies of the hyperplane sections.
We will look at one example to explain what we mean. Let $`X`$ be $`^n`$ and let $`L=𝒪_^n(d)`$. If we consider the intersection of $`n2`$ general divisors of $`|L|`$ we end with a surface in $`^n`$ which is a $`(d,\mathrm{},d)`$ complete intersection. The only surfaces among those which are rational surfaces (in fact, anticanonical) other than linear $`^2`$ are the quadric and the cubic hypersurface in $`^3`$ and a $`(2,2)`$ complete intersection in $`^4`$. This means that, using essentially the same ideas as in Section 1, one is able to give a result as precise as Theorem 1.3 regarding the property $`N_p`$ for $`X=^3`$ and $`L=𝒪_^2(2),𝒪_^2(3)`$ and for $`X=^4`$ and $`L=𝒪_^2(2)`$. Precisely one has that $`𝒪_^3(2)`$ satisfies property $`N_5`$ but not $`N_6`$, $`𝒪_^3(3)`$ satisfies property $`N_6`$ but not $`N_7`$ and $`𝒪_^4(2)`$ satisfies property $`N_5`$ but not $`N_6`$. These particular cases were known (cf. \[JPW\] and \[OP\] ). However for other pairs $`(X,L)`$ the game turns out to be much more complicated. For instance, if $`(X,L)=(^3,𝒪_^3(4))`$ or $`(^5,𝒪_^5(2))`$ by the above process we arrive at a K3 surface. Then knowing the syzygies of a K3 surfaces is equivalent to knowing the syzygies of its hyperplane section, a canonical curve. On this much less information is known and moreover, this information would depend (at least conjecturally) on the Clifford index of the hyperplane section. Finally all the other complete intersection surfaces of type $`(d,\mathrm{},d)`$ are surfaces of general type. Then $`L_C`$ will be a line bundle on the hyperplane section $`C`$ of degree strictly less than $`2g(C)2`$, and for those line bundles our knowledge is even more incomplete than for the canonical line bundle.
Because of all the above we need to carry out different arguments in this and in the following section, but before that we will state Theorem 2.1, which is based upon the work done in Section 1. Theorem 2.1 gives a necessary and sufficient condition for the line bundle $`H`$ giving the index of the Fano $`n`$-fold of index $`n1`$ to satisfy property $`N_p`$. Since a Fano $`n`$-fold $`(X,H)`$ of index $`n+1`$ or $`n`$ is $`(^n,𝒪_^n(1))`$ or a quadric hypersurface, we do not consider Fano $`n`$-folds of these indices in the theorem:
###### \Fanopn
Let $`X`$ be a Fano $`n`$-fold. Assume there exists an ample and base-point-free line bundle $`H`$ such that $`K_X^{}=(n1)H`$ (e.g., if $`X`$ is a Fano $`n`$-fold of index $`n1`$). Assume furthermore that $`H^np+3`$. Then $`H^1(M_{rH}lH)=0`$ for all $`r,l1`$ and $`H`$ satisfies property $`N_p`$.
Proof. The result is proven by induction on the dimension $`n`$ of $`X`$, starting the induction in dimension $`2`$. If $`n=2`$, the results follows from Theorem 1.3 . Indeed, the only thing to be checked is that $`K_XHp+3`$, $`p0`$. Then $`K_X=H`$, so we have the required inequality by hypothesis.
Let us assume the result to be true for all dimensions from $`2`$ to $`n1`$, and we will prove it for $`X`$ of dimension $`n`$. We first see that $`H`$ satisfies property $`N_0`$. It suffices to prove that
$$H^0(rH)H^0(H)\mathrm{@}>\alpha >>H^0((r+1)H)$$
surjects for all $`r1`$. Let $`Y`$ be a smooth irreducible member of $`|H|`$. By Kodaira Vanishing Theorem, $`H^1(lH)=0`$ for all $`l0`$, hence by Observation 1.1 , it suffices to see that
$$0H^0(rH_Y)H^0(H_Y)\mathrm{@}>\beta >>H^0((r+1)H_Y)0$$
surjects for all $`r1`$.
The variety $`Y`$ is a Fano $`(n1)`$-fold, $`K_Y=(K_X+H)_Y=(n2)H_Y`$ and $`H_Y^{n1}=H^n3`$. Therefore $`\beta `$ surjects by induction.
Therefore we have just proven that $`H`$ satisfies property $`N_0`$. To see it does satisfies property $`N_p`$, we argue as in Theorem 1.3 . Since $`H^1(lH)=0`$ for all $`l0`$, $`H`$ satisfies the same property $`N_p`$ as $`H_Y`$. Then by induction, since $`H_Y^{n1}=H^np+3`$, $`H`$ satisfies property $`N_p`$. $`\mathrm{}`$
Remark 2.2. Theorem 2.1 is in fact a characterization of the property $`N_p`$ satisfied by $`H`$. This follows because the same is true for the successive hyperplane section, which is a rational surface as we explained in the proof of Theorem 2.1 .
Now we want to show a result about the syzygies associated to the multiples of the line bundle $`H`$. As pointed out at the beginning of this section we need to use different ideas than those used in Section 1. Any result on the graded Betti numbers of the resolution of a variety can be realized in terms of Koszul cohomology. This was shown by M. Green. For a base-point-free line bundle $`L`$ we define the vector bundle $`M_L`$ as
$$0M_LH^0(L)𝒪L0.$$
$`()`$
Then regarding property $`N_p`$ Green showed the following criterion:
###### \GLlemma(\EL\ns, Section 1.)
Let $`L`$ be an ample, globally generated line bundle on a variety $`X`$. If the group $`H^1(^{p^{}+1}M_LsL)`$ vanishes for all $`0p^{}p`$ and all $`s1`$, then $`L`$ satisfies the property $`N_p`$. If in addition $`H^1(rL)=0`$, for all $`r1`$, then the above is a necessary and sufficient condition for $`L`$ to satisfy property $`N_p`$.
According to Theorem 2.3 the results (see for instance Theorem 1.3 ) obtained in Section 1 for the syzygies of a rational surface $`X`$ embedded by $`L`$ are equivalent to the vanishing of $`H^1(^{p^{}+1}M_LsL)`$ for all $`0p^{}p`$ and $`s1`$. To carry out the arguments for Fano $`n`$-folds, we would need however to have the vanishing of $`H^1(M_L^{p+1}sL)`$, which does not follow in general from the vanishing of $`H^1(^{p+1}M_LsL)`$. That is the reason why we prove Theorem 2.6 . Before that, we need to state two auxiliary lemmas:
Observation 2.4 . Let $`E`$ and $`L_1,\mathrm{},L_r`$ be coherent sheaves on a variety $`X.`$ Consider the map $`H^0(E)H^0(L_1+\mathrm{}+L_r)\mathrm{@}>\psi >>H^0(EL_1+\mathrm{}+L_r)`$ and the maps
$$\begin{array}{c}\text{ }H^0(E)H^0(L_1)\mathrm{@}>\alpha _1>>H^0(EL_1),\text{ }\hfill \\ \text{ }H^0(EL_1)H^0(L_2)\mathrm{@}>\alpha _2>>H^0(EL_1+L_2),\text{ }\hfill \\ \text{ }\mathrm{},\text{ }\hfill \\ \text{ }H^0(EL_1+\mathrm{}+L_{r1})H^0(L_r)\mathrm{@}>\alpha _r>>H^0(EL_1+\mathrm{}+L_r).\text{ }\hfill \end{array}$$
If $`\alpha _1,\mathrm{},\alpha _r`$ are surjective then $`\psi `$ is also surjective.
###### \Splemma(\GPfour\ns, Lemma 2.9)
Let $`X`$ be a projective variety, let $`q`$ be a nonnegative integer and let $`F_i`$ be a base-point-free line bundle on $`X`$ for all $`1iq`$. Let $`Q`$ be an effective line bundle on $`X`$ and let $`𝔮`$ be a reduced and irreducible member of $`|Q|`$. Let $`R`$ be a line bundle and $`G`$ a sheaf on $`X`$ such that
1. $`\text{H}^1(F_iQ^{})=0`$
2. $`\text{H}^0(M_{(F_{i_1}𝒪_𝔮)}\mathrm{}M_{(F_{i_q^{}}𝒪_𝔮)}R𝒪_𝔮)\text{H}^0(G)`$
$`\text{H}^0(M_{(F_{i_1}𝒪_𝔮)}\mathrm{}M_{(F_{i_q^{}}𝒪_𝔮)}RG𝒪_𝔮)\text{surjects for all}0q^{}q`$.
Then, for all $`0q^{\prime \prime }q`$ and any subset $`\{j_k\}\{i\}`$ with $`\mathrm{\#}\{j_k\}=q^{\prime \prime }`$ and for all $`0k^{}q^{\prime \prime }`$,
$$\begin{array}{c}\text{ }\text{H}^0(M_{F_{j_1}}\mathrm{}M_{F_{j_k^{}}}M_{(F_{j_{k^{}+1}}𝒪_𝔮)}\mathrm{}M_{(F_{j_{q^{\prime \prime }}}𝒪_𝔮)}R𝒪_𝔮)\text{H}^0(G)\text{ }\hfill \\ \text{ }\text{H}^0(M_{F_{j_1}}\mathrm{}M_{F_{j_k^{}}}M_{(F_{j_{k^{}+1}}𝒪_𝔮)}\mathrm{}M_{(F_{j_{q^{\prime \prime }}}𝒪_𝔮)}GR𝒪_𝔮)\text{ }\hfill \end{array}$$
surjects.
Now we are ready to prove
###### \tensNp
Let $`X`$ be a rational surface. Let $`B`$ be an ample and base-point-free line bundle such that $`K_XB4`$ or $`(X,B)=(^2,𝒪_^2(1))`$. Then $`H^1(M_{rB}^{p+1}lB)=0`$ for all $`r1`$, $`lp`$ and $`p1`$. In particular $`lB`$ satisfies property $`N_p`$ for all $`lp`$.
Proof. We first do the case $`K_XB4`$. The proof goes by induction on $`p`$. We start proving the case $`p=1`$. We want to show that $`H^1(M_{rB}^2lB)=0`$ for all $`r1`$ and all $`l1`$. We tensor (\*) associated to $`rB`$ by $`M_{rB}lB`$ and take global sections. As a piece of the long exact sequence of cohomology we obtain:
$$\begin{array}{c}\text{ }H^0(M_{rB}lB)H^0(rB)\mathrm{@}>\alpha >>H^0(M_{rB}(r+l)B)\text{ }\hfill \\ \text{ }H^1(M_{rB}^2lB)H^1(M_{rB}lB)H^0(rB).\text{ }\hfill \end{array}$$
By Theorem 1.3 , the last term of the above sequence is $`0`$. Then the vanishing of $`H^1(M_{rB}^2lB)`$ is equivalent to the surjectivity of $`\alpha `$. To have the surjectivity of $`\alpha `$, by Observation 2.4 it suffices to see the surjectivity of
$$H^0(M_{rB}lB)H^0(B)\mathrm{@}>\beta >>H^0(M_{rB}(l+1)B).$$
Since $`B`$ is base-point-free and ample, we can choose $`C`$ smooth and irreducible in $`|B|`$. Then to see the surjectivity of $`\beta `$, by Observation 1.1 and Lemma 2.5 , it suffices to see the surjectivity of
$$H^0(M_{rB_C}lB_C)H^0(B_C)\mathrm{@}>\gamma >>H^0(M_{rB_C}(l+1)B_C).$$
To obtain the surjectivity of $`\gamma `$ we show that the cokernel of $`\gamma `$ vanishes. The cokernel of $`\gamma `$ is $`H^1(M_{rB_C}^2lB_C)`$. Since $`K_XB4`$, deg$`B_C2g(C)+2`$, then by \[B\] , Theorem 1.12, $`M_{rB_C}`$ is semistable, and by \[Mi\] , Corollary 3.7 so is $`M_{rB_C}^2lB_C`$. A simple calculation shows that $`\mu (M_{rB_C}^2lB_C)>2g(C)2`$ and we get the desired vanishing.
Now we assume the result to be true for $`1,\mathrm{},p1`$ and we will prove it for $`p`$, i.e., we want to see that $`H^1(M_{rB}^{p+1}lB)=0`$. We go over the steps given to prove the case $`p=1`$. Since by induction $`H^1(M_{rB}^plB)=0`$, the vanishing we seek is equivalent to the surjectivity of
$$H^0(M_{rB}^plB)H^0(rB)\mathrm{@}>\alpha >>H^0(M_{rB}^p(r+l)B).$$
By Observation 2.4 it suffices to see that
$$H^0(M_{rB}^plB)H^0(B)\mathrm{@}>\beta >>H^0(M_{rB}^p(l+1)B)$$
surjects. Now choosing an irreducible and smooth curve $`C`$ in $`|B|`$ and using Observation 1.1 and Lemma 2.5 ,
we see that it is enough to show that
$$H^0(M_{rB_C}^plB_C)H^0(B_C)\mathrm{@}>\gamma >>H^0(M_{rB_C}^p(l+1)B_C)$$
surjects.
Finally $`\gamma `$ surjects if $`H^1(M_{rB_C}^{p+1}lB_C)=0`$. By \[B\] , Theorem 1.12 and \[Mi\] , Corollary 3.7, the bundle $`M_{rB_C}^{p+1}lB_C`$ is semistable. On the other hand its slope $`\mu `$ is
$$\frac{(p+1)rB^2}{rB^2g(C)}+lB^2,$$
so we will conclude our argument if we see that $`\mu >2g(C)2`$. This follows from $`lp2,K_XB4`$.
As in the case $`p=1`$, one can alternatively deduce the surjectivity of $`\gamma `$ from \[B\] , Proposition 2.2.
To finish the proof we take care of the case $`(X,B)=(^2,𝒪_^2(1)`$. The proof goes along the same lines as before. We want the vanishing of $`H^1(M_{rB}^{p+1}lB)`$ for all $`r1`$ and all $`lp`$. This follows from the vanishing of $`H^1(M_{𝒪_{^1(r)}}^{p+1}𝒪_^1(l))`$, which can be easily checked as $`M_{𝒪_{^1(r)}}`$ is direct sum of copies of $`𝒪_^1(1)`$. $`\mathrm{}`$
We use Theorem 2.6 to obtain the following results on Koszul cohomology and syzygies for multiples of $`H`$.
###### \FanoNp
Let $`X`$ be a Fano $`n`$-fold. Assume there exists an ample and base-point-free line bundle $`H`$ such that $`K_X=mH`$, with $`mn1`$ (e.g., if $`X`$ is a Fano $`n`$-fold of index $`mn1`$). Assume furthermore that $`H^n4`$ if $`m=n1`$. Then $`H^1(M_{rH}^{p+1}lH)=0`$ for all $`r1`$, $`lp`$ and $`p1`$.
Proof. To prove $`H^1(M_{rH}^{p+1}lH)=0`$ we argue by induction on the dimension of $`X`$, the cornerstone being now Theorem 2.6 , and by induction on $`p`$.
Let first $`p=1`$. We want to obtain the vanishing of $`H^1(M_{rH}^2lH)`$ and, as announced, we do it by induction on the dimension $`n`$. If $`n=2`$ the result is a particular case of Theorem 2.6 . In fact, $`K_X=mH`$ with $`m1`$. If $`m=1`$, $`K_XH4`$ follows directly by hypothesis. If $`m2`$, $`K_XH2`$, since $`H`$ is ample, and $`K_XH4`$ unless $`(X,H)=(^2,𝒪_^2(1))`$.
Now we assume the result to be true for dimensions $`2`$ to $`n1`$ and we will prove it for dimension $`n`$. From (\*) we obtain
$$\begin{array}{c}\text{ }H^0(M_{rH}lH)H^0(rH)\mathrm{@}>\alpha >>H^0(M_{rH}(r+l)H)\text{ }\hfill \\ \text{ }H^1(M_{rH}^2lH)H^1(M_{rH}lH)H^0(rH).\text{ }\hfill \end{array}$$
By the vanishing of the first cohomology of the multiples of $`H`$, Theorem 2.1 implies the vanishing of the last term of the above sequence. Then the vanishing we seek is equivalent to the surjectivity of $`\alpha `$. By Observation 2.4 we see that it will suffice to prove that
$$H^0(M_{rH}lH)H^0(H)\mathrm{@}>\beta >>H^0(M_{rH}(l+1)H)$$
surjects for all $`r,l1`$. The conditions needed to apply Observation 1.1 and Lemma 2.5 follow from Theorem 2.1 and from the fact that $`H^1(M_{rH})=0`$, for $`H^1(𝒪_X)=0`$. Then it suffices to see
$$H^0(M_{rH_Y}lH_Y)H^0(H_Y)\mathrm{@}>\gamma >>H^0(M_{rH_Y}(l+1)H_Y)$$
surjects, which follows because the result is true for $`p=1`$ and $`Y`$ of dimension $`n1`$, by induction hypothesis.
Now we assume the result to be true for $`1,\mathrm{},p1`$ and we will prove it for $`p`$. We argue again by induction on the dimension. In dimension $`2`$ the result follows from Theorem 2.6 as in the case $`p=1`$. We therefore assume the result to be true for $`p`$ and dimensions $`2`$ to $`n1`$ and we will prove it for $`p`$ and dimension $`n`$, i.e., we will prove $`H^1(M_{rH}^{p+1}lH)=0`$ for all $`r1`$, $`lp`$. From (\*) we obtain
$$\begin{array}{c}\text{ }H^0(M_{rH}^plH)H^0(rH)\mathrm{@}>\alpha >>H^0(M_{rH}^p(r+l)H)\text{ }\hfill \\ \text{ }H^1(M_{rH}^{p+1}lH)H^1(M_{rH}^plH)H^0(rH).\text{ }\hfill \end{array}$$
By induction hypothesis on $`p`$, we have the vanishing of the last term of the above sequence. Therefore the vanishing we seek is equivalent to the surjectivity of $`\alpha `$. By Observation 2.4 we see that it will suffice to prove that
$$H^0(M_{rH}^plH)H^0(H)\mathrm{@}>\beta >>H^0(M_{rH}^p(l+1)H)$$
surjects for all $`r1,lp`$. The conditions needed to apply Observation 1.1 and Lemma 2.5
follow from induction hypothesis on $`p`$. Then it suffices to see
$$H^0(M_{rH_Y}^plH_Y)H^0(H_Y)\mathrm{@}>\gamma >>H^0(M_{rH_Y}^p(l+1)H_Y)$$
surjects, which follows because the result is true for $`p`$ and $`Y`$ of dimension $`n1`$, by induction hypothesis on $`n`$. $`\mathrm{}`$
###### Corollary 2.8
Let $`X`$ be a Fano $`n`$-fold. Assume there exists an ample and base-point-free line bundle $`H`$ such that $`K_X=mH`$, with $`mn1`$ (e.g., if $`X`$ is a Fano $`n`$-fold of index $`m1`$). Assume furthermore that $`H^n4`$ if $`m=n1`$. Then $`lH`$ satisfies property $`N_p`$ for all $`lp`$.
Proof.
Since we work in characteristic $`0`$, then
$$H^1(\stackrel{i}{}M_{lH}slH)=0,\text{ for all }lp,s1,1ip+1.$$
Therefore, according to \[GL\] , $`lH`$ satisfies property $`N_p`$.
## 3. Fano $`n`$-folds of index $`n3`$.
In this last section we deal with Fano $`n`$-folds of index $`n3`$, $`n4`$. We use the same ideas as in the second part of Section 2. We start by studying under what conditions $`mH`$ is very ample and satisfies property $`N_0`$:
###### \CYpn
Let $`X`$ be a Fano $`n`$-fold such that $`K_X=mH`$, $`H`$ is ample and base-point-free, and $`m=n31`$ (for instance, if $`X`$ is a Fano $`n`$-fold of index $`m=n3`$). Let $`L=kH`$.
(1) If $`k4`$, then $`L`$ satisfies property $`N_0`$.
(2) Let $`k=3`$; $`L`$ satisfies property $`N_0`$ if and only if the morphism induced by $`|H|`$ does not map $`X2:1`$ onto $`^n`$.
(3) Let $`k=2`$; if $`|H|`$ does not map $`X`$ onto a variety of minimal degree other than $`^n`$ nor maps $`X2:1`$ onto $`^n`$, then $`L`$ satisfies property $`N_0`$.
Proof. We will prove the result by induction on $`n`$. We start at $`n=4`$. We have to deal with several cases.
Case 1: $`k4`$. We want to prove that $`L=kH`$ satisfies property $`N_0`$, or equivalently, that
$$H^0(skH)H^0(kH)H^0((s+1)kH)$$
surjects for all $`s1`$. This follows from a more general result, namely,
$$H^0(kH)H^0(H)\mathrm{@}>\alpha >>H^0((k+1)H)$$
surjects for all $`k4`$. Indeed, the surjectivity of $`\alpha `$ follows from \[Mu\] , p. 41, Theorem 2, since by Kodaira Vanishing Theorem, $`H`$ is $`3`$-regular.
Case 2: $`k=3`$. We distinguish two subcases.
Case 2.1: First let us assume that $`|H|`$ does not map $`X`$ onto $`^4`$. We will show
$$H^0(kH)H^0(H)\mathrm{@}>\alpha >>H^0((k+1)H)$$
for all $`k3`$. If $`k4`$, we have already seen that $`\alpha `$ surjects. To show the surjectivity for $`k=3`$, let $`Y`$ be a smooth irreducible member of $`|H|`$. By Observation 1.1 , since $`H^1(lH)=0`$, for all $`l0`$ by Kodaira Vanishing Theorem, it will suffice to check that
$$H^0(3H_Y)H^0(H_Y)\mathrm{@}>\beta >>H^0(4H_Y)$$
surjects. By adjunction, $`(Y,H_Y)`$ is a polarized Calabi-Yau threefold, and, since $`H^1(𝒪_X)=0`$, $`|H_Y|`$ does not map $`Y`$ onto $`^3`$. Then $`\beta `$ surjects (cf. \[GP3\] , proof of Theorem 1.4, case 1).
Case 2.2: Now assume $`|H|`$ does map $`X`$ onto $`^4`$. We assume first that the map induced by $`|H|`$ is not $`2:1`$ and we will prove that
$$H^0(3lH)H^0(3H)H^0(3(l+1)H)$$
surjects for all $`l1`$. For that, by Observation 2.4 it is enough that
$$\begin{array}{c}\text{ }H^0(rH)H^0(H)\mathrm{@}>\alpha >>H^0((r+1)H)\text{ for all }r5\text{ }\hfill \\ \text{ }H^0(3H)H^0(2H)\mathrm{@}>\gamma >>H^0(5H)\text{ }\hfill \end{array}$$
surject. The map $`\alpha `$ was seen to be surjective in Case 1. For the surjectivity of $`\gamma `$ we choose smooth irreducible $`Y`$ in $`|H|`$ and we write the following commutative diagram,
$$\begin{array}{ccccc}H^0(2H)H^0(2H)& & H^0(2H)H^0(3H)& & H^0(2H)H^0(3H_Y)\\ \mathrm{@}VV\delta V\mathrm{@}VV\gamma V\mathrm{@}VVϵV\\ H^0(4H)& & H^0(5H)& & H^0(5H_Y),\end{array}$$
obtained from the sequence
$$0H^{}𝒪_X𝒪_Y0$$
$`(\mathrm{3.1.1}),`$
having in account that $`H^1(rH)=0`$ for all $`r0`$. To see the surjectivity of $`\delta `$ we construct another diagram arising from $`(\mathrm{3.1.1})`$:
$$\begin{array}{ccccc}H^0(2H)H^0(H)& & H^0(2H)H^0(2H)& & H^0(2H)H^0(2H_Y)\\ \mathrm{@}VV\eta V\mathrm{@}VV\delta V\mathrm{@}VV\nu V\\ H^0(3H)& & H^0(4H)& & H^0(4H_Y).\end{array}$$
Thus we need to see that $`\eta `$, $`ϵ`$ and $`\nu `$ are surjective. By Observation 1.1 , for the surjectivity of $`\eta `$ it suffices to see the surjectivity of
$$H^0(2H_Y)H^0(H_Y)\mathrm{@}>\mu >>H^0(3H_Y)$$
and for the surjectivity of $`ϵ`$ and $`\nu `$, since $`H^1(H)=0`$, it suffices to see that
$$\begin{array}{c}\text{ }H^0(3H_Y)H^0(2H_Y)\mathrm{@}>\pi >>H^0(5H_Y)\text{ }\hfill \\ \text{ }H^0(2H_Y)H^0(2H_Y)\mathrm{@}>\rho >>H^0(4H_Y)\text{ }\hfill \end{array}$$
both surject. $`(Y,H_Y)`$ is a polarized Calabi-Yau threefold such that the morphism induced by $`|H_Y|`$ maps $`Y`$ onto to $`^3`$ but is not $`2:1`$. Then the surjectivity of $`\pi `$ and $`\rho `$ is explicitly claimed (and proved) in the proof of Theorem 1.4, \[GP3\] and the surjectivity of $`\mu `$ is also proved there, although not explicitly said.
To end Case 2.2 assume $`|H|`$ maps $`X`$ $`2:1`$ onto $`^4`$. If $`Y`$ is a smooth and irreducible member of $`|H|`$, then $`|H_Y|`$ maps $`Y`$ $`2:1`$ onto $`^3`$, so according to \[GP3\] , Theorem 1.4, $`H_Y`$ is not very ample, nor is $`H`$.
Case 3: $`k=2`$. We want to see that if $`H`$ is ample, base-point-free and $`|H|`$ does not map $`X`$ onto a variety of minimal degree (different from $`^4`$), nor does it map $`X`$ $`2:1`$ onto $`^4`$, then
$$H^0(2lH)H^0(2H)H^0((2l+2)H)$$
surjects for all $`l1`$.
Case 3.1: Assume $`h^0(H)6`$. By Observation 2.4 , it suffices to prove that
$$H^0(rH)H^0(H)\mathrm{@}>\alpha >>H^0((r+1)H)$$
surjects for all $`r2`$. If $`r3`$, we have seen that $`\alpha `$ surjects while proving Case 1 and Case 2.1. To see that
$$H^0(2H)H^0(H)\mathrm{@}>\alpha >>H^0(3H)$$
surjects let $`Y`$ be as before a smooth and irreducible member of $`|H|`$. By Observation 1.1 it suffices to see that
$$H^0(2H_Y)H^0(H_Y)\mathrm{@}>\beta >>H^0(3H_Y)$$
surjects. Now $`(Y,H_Y)`$ is a polarized Calabi-Yau satisfying the hypothesis of Theorem 1.7.1, \[GP3\] . Then the map $`\beta `$ surjects, as seen in the proof of Theorem 1.7, \[GP3\] .
Case 3.2: Assume now that $`h^0(H)=5`$ and that the map induced by $`|H|`$ from $`X`$ onto $`^4`$ is not $`2:1`$. By Observation 2.4 it suffices to see that
$$\begin{array}{c}\text{ }H^0(rH)H^0(H)H^0((r+1)H)\text{ for all }r4\text{ and }\text{ }\hfill \\ \text{ }H^0(2H)H^0(2H)H^0(4H)\text{ }\hfill \end{array}$$
surject, and this has been proved in Case 1 and Case 2.2.
To prove the result for a Fano variety of arbitrary dimension $`n`$ we will argue by induction. Let $`Y`$ be smooth and irreducible member of $`|H|`$. If $`K_X=(n3)H`$, then $`K_Y=(n2)H_Y`$ and $`h^0(H_Y)=h^0(H)1`$. Moreover, if the image of $`X`$ by the morphism induced by $`|H|`$ is a variety of minimal degree, so is the image of $`Y`$ by the morphism induced by $`|H_Y|`$, and the degree of both morphisms is the same. Going over the arguments in the case $`n=4`$, we see that the key point was to show the surjectivity of the maps
$$\begin{array}{c}\text{ }H^0(rH)H^0(H)\mathrm{@}>\alpha >>H^0((r+1)H)\text{ for all }k4,\text{ }\hfill \\ \text{ }H^0(3H)H^0(H)\mathrm{@}>\beta >>H^0(4H)\text{ if }h^0(H)n+2\text{ }\hfill \\ \text{ }H^0(2H)H^0(H)\mathrm{@}>\gamma >>H^0(3H)\text{ if }h^0(H)n+2\text{ and the image of }X\text{ }\text{ }\hfill \\ \text{ }\text{ by the map induced by }|H|\text{ is not a variety of minimal degree }\text{ }\hfill \\ \text{ }H^0(3H)H^0(2H)\mathrm{@}>\delta >>H^0(5H)\text{ if }h^0(H)=n+1\text{ and the map induced by }|H|\text{ is not }2:1\text{ }\hfill \\ \text{ }H^0(2H)H^0(2H)\mathrm{@}>ϵ>>H^0(4H)\text{ if }h^0(H)=n+1\text{ and the map induced by }|H|\text{ is not }2:1\text{ }\hfill \\ \text{ }H^0(2H)H^0(H)\mathrm{@}>\eta >>H^0(3H)\text{ if }h^0(H)=n+1\text{ and the map induced by }|H|\text{ is not }2:1.\text{ }\hfill \end{array}$$
Let us assume therefore the surjectivity of all the above maps for Fano varieties $`(X,H)`$ of dimension $`4,\mathrm{},n1`$. Arguing as in the case $`n=4`$, by Observation 1.1 and because $`H^1(rH)=0`$ for all $`r0`$, we have that all the above maps are also surjective if the dimension of $`X`$ is $`n`$.
Then by Observation 2.4 the surjectivity of $`\alpha `$ implies (1), the surjectivity of $`\alpha `$, $`\beta `$ and $`\delta `$ implies the “if” part of (2), and the surjectivity of $`\alpha `$, $`\beta `$, $`\gamma `$ and $`ϵ`$ implies (3). Finally, if the morphism induced by $`|H|`$ is a double cover of $`^n`$, so is the morphism induced on $`Y`$ by $`|H_Y|`$, when $`Y`$ is an irreducible and smooth member of $`|H|`$. Then $`3H_Y`$ is not very ample (by induction hypothesis), nor is $`3H`$. $`\mathrm{}`$
###### \CYNp
Let $`X`$ be a Fano $`n`$-fold of index $`m=n3`$. Assume that $`K_X=mH`$, and $`H`$ is ample and base-point-free. Let $`L=kH`$. Assume furthermore that $`h^0(H)n+2`$, i.e., that $`|H|`$ do not map $`X`$ onto $`^n`$. If $`kp+2`$ and $`p1`$, then $`L`$ satisfies property $`N_p`$.
Proof. The proof is again by induction on the dimension. The first step is $`n=4`$. Given $`p1`$, we want to prove that
$$H^1(M_L^{p+1}sL)=0\text{ if }L=kH\text{}kp+2\text{}l1\text{}$$
$`(\mathrm{3.2.1})`$
We will prove this more general fact, namely that
$$H^1(M_{kH}^{p+1}sH)=0\text{ for all }kp+2\text{}sp+2\text{.}$$
$`(\mathrm{3.2.2})`$
We argue now by induction on $`p`$. If $`p=1`$, we want to prove that
$$H^1(M_{kH}^2sH)=0\text{ for all }k3\text{}s3\text{.}$$
$`(\mathrm{3.2.3})`$
Since $`H^1(M_{kH}sH)=0`$ by Theorem 3.1 , then (3.2.3) is equivalent to the surjectivity of
$$H^0(M_{kH}sH)H^0(kH)H^0(M_{kH}(k+s)H)\text{ for all }k3\text{}s3\text{,}$$
and by Observation 2.4 it suffices to see that
$$H^0(M_{kH}H)H^0(H)H^0(M_{kH}(k+1)H)\text{ surjects for all }k3\text{.}$$
Choose smooth and irreducible $`3`$-fold $`Y`$ in $`|H|`$. Because of Theorem 3.1 , Observation 1.1 and Lemma 2.5 it suffices to have the surjectivity of
$$H^0(M_{kH_Y}H_Y)H^0(H_Y)H^0(M_{kH_Y}(k+1)H_Y)$$
for all $`k3`$. By adjunction $`Y`$ is a Calabi-Yau threefold and since $`H^1(𝒪_X)=0`$, $`h^0(H_Y)5`$, hence from the proof of \[GP3\] , Theorem 1.4, the above map is surjective.
We complete now the proof of the result for $`n=4`$. We may assume the result proved until $`p1`$. Now we want to see that
$$H^1(M_{kH}^{p+1}sH)=0\text{ for all }kp+2\text{}sp+2\text{}$$
$`(\mathrm{3.2.4})`$
We argue similarly to the case $`p=1`$. By induction on $`p`$ we may conclude that the sought vanishing is equivalent to the surjectivity of
$$H^0(M_{kH}^psH)H^0(kH)H^0(M_{kH}^p(k+s)H)\text{ for all }kp+2\text{}sp+2\text{,}$$
and by Observation 2.4 it suffices to see that
$$H^0(M_{kH}^pH)H^0(H)H^0(M_{kH}^p(k+1)H)\text{ surjects for all }kp+2\text{.}$$
Finally we choose a smooth and irreducible $`3`$-fold $`Y`$ in $`|H|`$. Because of Theorem 3.1 , Observation 1.1 and Lemma 2.5 it suffices to have the surjectivity of
$$H^0(M_{kH_Y}H_Y)H^0(H_Y)H^0(M_{kH_Y}(k+1)H_Y)$$
for all $`kp+2`$. By adjunction $`Y`$ is a Calabi-Yau threefold and since $`H^1(𝒪_X)=0`$, $`h^0(H_Y)5`$, hence according to the proof of \[GP3\] , Theorem 1.4, the above map is surjective.
Now assume $`n>4`$. Recall that we want to show that
$$H^1(M_L^{p+1}sL)=0\text{ if }L=kH\text{}kp+2\text{}l1\text{}$$
$`(\mathrm{3.2.5})`$
and as before we will prove this more general fact, namely that
$$H^1(M_{kH}^{p+1}sH)=0\text{ for all }kp+2\text{}sp+2\text{.}$$
$`(\mathrm{3.2.6})`$
We argue now by induction on $`p`$ and $`n`$. If $`p=1`$, we want to prove that
$$H^1(M_{kH}^2sH)=0\text{ for all }k3\text{}s3\text{.}$$
$`(\mathrm{3.2.7})`$
Since $`H^1(M_{kH}sH)=0`$ by Theorem 3.1 , then (3.2.7) is equivalent to the surjectivity of
$$H^0(M_{kH}sH)H^0(kH)H^0(M_{kH}(k+s)H)\text{ for all }k3\text{}s3\text{,}$$
and by Observation 2.4 it suffices to see that
$$H^0(M_{kH}H)H^0(H)\mathrm{@}>\alpha >>H^0(M_{kH}(k+1)H)\text{ surjects for all }k3\text{.}$$
The surjectivity of this map has been proven under the hypothesis of the theorem, when the dimension $`n`$ of $`X`$ is $`4`$, and we will assume it proved also if dimension of $`X`$ is $`n1`$.
Choose smooth and irreducible $`(n1)`$-fold $`Y`$ in $`|H|`$. Because of Theorem 3.1 , Observation 1.1 and Lemma 2.5 it suffices to have the surjectivity of
$$H^0(M_{kH_Y}H_Y)H^0(H_Y)H^0(M_{kH_Y}(k+1)H_Y)$$
for all $`k3`$. By adjunction $`K_Y=(n4)H_Y`$ and since $`H^1(𝒪_X)=0`$, $`h^0(H_Y)n+1`$, hence by induction hypothesis, $`\alpha `$ surjects.
The proof of the general case follows the same steps as the proof for $`n=4`$. Recall that we want to proof that
$$H^1(M_{kH}^{p+1}sH)=0\text{ for all }kp+2\text{}sp+2\text{.}$$
$`(\mathrm{3.2.8})`$
Using now the induction hypothesis for $`n1`$ one conclude the result, exactly in the same fashion as we have just done when $`p=1`$. $`\mathrm{}`$
## References |
warning/0001/hep-ph0001326.html | ar5iv | text | # LIGHT MESON SPECTROSCOPY: RECENT DEVELOPMENTS and DAFNE
## 1 Introduction
The last few years have seen rapid and exciting developments in light meson spectroscopy, largely as a result of the analysis of high-statistics experiments using hadron beams. The most notable discoveries have come from studies of $`P\overline{P}`$ annihilation at LEAR and $`\pi ^{}P`$ at the AGS (BNL) and VES (Serpukhov). In both processes we have seen that detailed amplitude analyses of high-statistics events samples (ca.1M events) have made possible the identification of very interesting parent resonances in otherwise relatively mundane final states such as $`3\pi `$. This has led for example to the discovery of a glueball candidate in $`3\pi ^o`$ and an exotic hybrid candidate in $`(3\pi )^{}`$. Concurrently we have seen impressive progress in the study of conventional $`q\overline{q}`$ mesons (which must be identified as a background to more unusual resonances), and at this meeting we have heard important new results from VEPP which appear to confirm the predictions of Close, Isgur and Kumano for a $`K\overline{K}`$-molecule assignment for the scalars $`f_0(980)`$ and $`a_0(980)`$. In this case at least, progress has come from an $`e^+e^{}`$ facility rather than a hadronic one. In this introduction I will give a brief summary of the status of the various sectors of meson spectroscopy, and then discuss two areas in which DAFNE can make very important contributions, excited vectors and C=(+) mesons.
## 2 Recent developments in light meson spectroscopy.
### 2.1 Glueballs
The gluonic degree of freedom in QCD leads to more physical resonances than are predicted by the naive $`q\overline{q}`$ quark model. Pure-glue “glueball” states have been studied using many theoretical approaches, the most recent and (presumably) the most accurate of which is lattice gauge theory (LGT). In recent years LGT has largely displaced other theoretical methods for treating these most unfamiliar of hadrons. A recent high-statistics LGT study of the glueball spectrum to ca.4 GeV has been reported by Morningstar and Peardon(see Fig.1); for other recent discussions of glueballs and LGT see Teper and Michael. The lattice predicts that the lightest (assumed unmixed with $`q\overline{q}`$) glueball is a scalar, with a mass of about 1.7 GeV. Additional glueballs lie well above 2 GeV, with a $`0^+`$ and a $`2^{++}`$ appearing at masses of $`2.42.6`$ GeV. Spin-parity exotic glueballs are expected at rather higher masses; in the Morningstar and Peardon study the lightest exotic glueball was found to be a $`2^+`$ at just above 4 GeV. For experimental studies of meson spectroscopy below ca.2.2 GeV, the subject of glueballs thus reduces to the search for an extra scalar.
Scalars unfortunately comprise the most obscure part of the spectrum, and there are at least three states that might a priori be identified with a scalar glueball, the $`f_0(1370)`$, the LEAR state $`f_0(1500)`$ and the $`\psi `$ radiative candidate $`f_0(1710)`$.
There are outstanding problems with each of these assignments. In view of LGT mass predictions the $`f_0(1500)`$ and $`f_0(1710)`$ appear most plausible, but neither of these states shows the flavor-blind pattern of decay couplings naively expected for a flavor-singlet glueball. The $`f_0(1500)`$ as seen by Crystal Barrel in $`\pi ^o\pi ^o`$ is shown in Fig.2. The results of some analyses, taken from the 1998 PDG, are shown in Table 1. Although essentially all these numbers are controversial, it is clear that the $`\pi \pi /K\overline{K}`$ branching ratios of the $`f_0(1500)`$ and $`f_0(1710)`$ are both far from the approximate equality expected for a flavor-singlet. We also note that the two lighter states have large $`4\pi `$ modes, which have not been considered in glueball decay models.
The $`K\overline{K}`$ mode of the $`f_0(1500)`$ is difficult to isolate, but appears to be weaker than one would expect for flavor-singlet couplings to $`\pi \pi `$, $`K\overline{K}`$ and $`\eta \eta `$. Conversely, the $`f_0(1710)`$ has a strong $`K\overline{K}`$ mode but a weak $`\pi \pi `$ coupling. The determination of the $`K\overline{K}`$ branching fraction of the $`f_0(1500)`$ has recently been reanalysed by Ableev et al, who find a much larger branching fraction than quoted in Table 1, but still rather smaller than expected for a flavor singlet. Several models, for example that of Amsler and Close, invoke important $`n\overline{n}Gs\overline{s}`$ mixing to explain the observed branching fractions these scalar states, so the scalar glueball basis state may actually be distributed over several physical resonances. In the final section we will discuss how this possibility could be tested at an $`e^+e^{}`$ facility.
For completeness we note that BES has reported evidence for a possible narrow state in several channels, including $`P\overline{P}`$, $`\pi \pi `$, $`K\overline{K}`$ and $`\eta \eta `$, at about 2.2 GeV. Although one does expect a tensor glueball not far above this mass, and the narrow glueball candidate $`f_0(1500)`$ suggests that the tensor glueball might have a narrow width, the statistical significance of the reported signals near 2.2 GeV is rather low. Another problem is that the Crystal Barrel has shown that the $`P\overline{P}`$ and $`\eta \eta `$ modes cannot both be as large as claimed by BES, since the state does not appear with the corresponding strength in $`P\overline{P}\eta \eta `$. This state clearly “needs confirmation”.
### 2.2 Hybrid Mesons
In addition to glueballs, we also expect the glue degree of freedom to lead to “hybrid mesons” in which the $`q\overline{q}`$ pair is combined with glue in an excited state. Hybrids are especially attractive experimentally, because they span flavor nonets (so they can be searched for in many flavor channels), and have “exotic” $`J^{PC}`$ combinations such as $`1^+`$ that are forbidden to $`q\overline{q}`$ states. (Hybrids span all $`J^{PC}`$ quantum numbers, both exotic and non-exotic.) The $`J^{PC}`$ content of the lowest-lying hybrid multiplet is model dependent: The lowest-lying exotics in this first hybrid multiplet according to the flux-tube model are
$$J^{PC}(\mathrm{lightest}\mathrm{flux}\mathrm{tube}\mathrm{hybrid}\mathrm{exotics})=0^+,1^+,2^+$$
(1)
and are expected to be approximately degenerate. In contrast, in the bag model the lightest hybrid multiplet only has the single exotic
$$J^{PC}(\mathrm{lightest}\mathrm{bag}\mathrm{model}\mathrm{hybrid}\mathrm{exotic})=1^+.$$
(2)
The difference is due to assumptions about confinement; the bag model has a confining boundary condition that discriminates between color electric and magnetic fields, which gives a TM $`(1^{})`$ gluon more energy than TE $`(1^+)`$. The flux-tube model in contrast simply has a spatially excited interquark string and makes no reference to color field vectors. (Preliminary LGT results found the $`1^+`$ hybrid at a significantly lower mass than the $`0^+`$, as expected in the bag model but not the flux-tube model; more recent results by the same collaboration now find the $`1^+`$ and $`0^+`$ exotic hybrids rather closer in mass.) The mass of the lightest hybrid meson multiplet is expected by theorists to be near 1.9 GeV. The bag model typically finds a somewhat lower scale of ca.1.5 GeV, which is now deprecated because it disagrees with LGT. This 1.9 GeV estimate was originally due to the flux tube model, and has been (approximately) confirmed by recent LGT studies, which find a mass of about 2.0 GeV for the lightest hybrid. For a recent review of LGT predictions for these states see Michael.
We now have strong evidence for a true $`J^{PC}=1^+`$ exotic at 1.6 GeV in $`\rho \pi `$ at BNL and VES (see Fig.3 for the $`\rho \pi `$ mode), and $`\eta ^{}\pi `$ and $`b_1\pi `$ at VES. In addition a rather lighter state at 1.4 GeV in $`\eta \pi `$ has been reported by BNL and Crystal Barrel. Thus, experimental hadron spectroscopy may finally have found the hybrid mesons anticipated by theorists for about 25 years. Of course there is an unresolved concern that these experimental masses are somewhat lighter than the theoretical expectation of $`1.9`$-$`2.0`$ GeV. There are also nonexotic hybrid candidates such as the $`\pi (1800)`$; a recent and reasonably complete review of light meson spectroscopy which discusses hybrid candidates in more detail was recently completed by Godfrey and Napolitano..
Hybrid strong decays are in a confused state. The flux-tube model predicts that the dominant modes should be S+P two-body combinations such as $`\pi f_1`$ and $`\pi b_1`$. The reported observations of hybrids however have for the most part been in the more familiar S+S modes such as $`\pi \eta `$, $`\pi \eta ^{}`$ and $`\pi \rho `$, although there is some evidence for $`\pi b_1`$ and $`\pi f_1`$. VES has reported relative branching fractions for the $`\pi _1(1600)`$ exotic hybrid candidate that actually suggest comparable branching fractions to S+S and S+P modes. Clearly the modelling of strong decays of hybrids is at an early stage, and the experimental determination of relative $`\pi _1`$ hybrid branching fractions will be a very useful contribution (assuming that these states persist with improved statistics!).
Since $`q\overline{q}g`$ hybrids span flavor nonets, there should be many more hybrids near 1.5 GeV if the reports of $`\pi _1`$ exotic hybrids near this mass are correct. Specific models of hybrids such as the flux-tube and bag models find that the majority of light hybrids have nonexotic $`J^{PC}`$. In the flux tube model the lightest hybrid multiplet contains five nonexotic quantum numbers,
$$J^{PC}(\mathrm{lightest}\mathrm{flux}\mathrm{tube}\mathrm{hybrid}\mathrm{nonexotics})=0^+,1^{},1^{++},1^+,2^+$$
(3)
whereas in the bag model the lightest hybrid multiplet contains just three nonexotics,
$$J^{PC}(\mathrm{lightest}\mathrm{bag}\mathrm{model}\mathrm{hybrid}\mathrm{nonexotics})=0^+,1^{},2^+.$$
(4)
Note that both models include a $`1^{}`$ flavor nonet in the set of lowest-lying hybrid mesons. Thus the $`1^{}`$ sector should show evidence of overpopulation relative to the naive quark model, which can be tested at DAFNE. We shall return to this topic in the next section.
### 2.3 Multiquarks and Molecules
In the 1970s it was thought that the existence of many basis states in the $`q^2\overline{q}^2`$ sector implied a very rich spectrum of multiquark resonances. Calculations in specific models such as the MIT bag model and color-truncated potential models appeared to support this picture. However it was subsequently realized that the overlap of these multiquark basis states with the continuum of two color-singlet $`(q\overline{q})(q\overline{q})`$ mesons implied that the multiquark systems need not appear as resonances; they might instead simply be components of nonresonant two-meson continua.
An exception to this absence of four-quark resonances can occur if the multiquark system lies well below all two-body decay thresholds, or if there is a strongly suppressed coupling to the open decay channels; in these cases we might still expect to identify a bag-model “cluster” multiquark resonance.
Nature appears to favor a different type of multiquark system, in which largely unmodified color-singlet $`q\overline{q}`$ or $`qqq`$ hadrons are weakly bound by the residual nuclear forces between color singlets. Examples of such quasinuclear multiquark systems abound; the table of nuclear species gives far more examples than we have of individual hadrons, and hypernuclei extend these systems into strangeness. In the mesonic sector, however, just two possible examples are widely cited, the scalar mesons $`f_0(980)`$ and $`a_0(980)`$.
These scalars are candidates for weakly bound $`K\overline{K}`$ nuclei, “molecules”, due to their masses and quantum numbers (which are those of an S-wave $`K\overline{K}`$ pair), and also because their hadronic couplings appear bizarre for $`n\overline{n}`$ states, which should be very broad and for $`I=0`$ should couple strongly to $`\pi \pi `$. Another problem with a conventional assignment is the two-photon widths of these states, which are much smaller than expected for $`q\overline{q}`$ but are rather similar to predictions for $`K\overline{K}`$ bound states or $`ns\overline{n}\overline{s}`$ four-quark clusters. An interesting test of the nature of these states was proposed by Close, Isgur and Kumano; the theoretical radiative branching fractions from the $`\varphi `$ depend rather strongly on the quark model assignments, and for $`q\overline{q}`$ versus $`K\overline{K}`$ states are
$$B(\varphi \gamma f_0(980),\gamma a_0(980))=\{\begin{array}{cc}\hfill 410^5:& K\overline{K}(\mathrm{both}\mathrm{states})\hfill \\ \hfill 110^5:& f_0(980)=s\overline{s}\hfill \\ \hfill 10^6:& f_0(980),a_0(980)=n\overline{n}.\hfill \end{array}$$
(5)
Close et al. note that the ratio $`\varphi \gamma a_0(980)/\gamma f_0(980)`$ is also of interest, since it can distinguish between different multiquark spatial wavefunctions. For a $`K\overline{K}`$ molecule this ratio is 1, whereas for an $`(ns)(\overline{n}\overline{s})`$ system it is 9.
At this meeting we have heard that the new experimental results from VEPP are not far from the Close et al. predictions for a $`K\overline{K}`$ molecule. (The VEPP experimental branching fractions $`B(\varphi \gamma f_0(980),\gamma a_0(980))`$ are somewhat larger than $`410^5`$, but are roughly consistent with Close et al. given the current errors.) Earlier experimental indications of much larger branching fractions to the 980 MeV states were biased by large nonresonant contributions well below 980 MeV, which had not clearly been identified.
Presumably there are many meson-meson bound states, since many other meson pairs experience attractive residual nuclear interactions. Unlike glueballs and hybrids, the spectrum of molecular states beyond $`K\overline{K}`$ and the nuclei and hypernuclei has received little theoretical attention. There are quark model and meson-exchange model predictions that some vector meson pairs may bind, but to date there has been little systematic investigation of the expected spectrum. As our understanding of residual hadronic forces improves, we can expect this to be one of the interesting areas of development in hadron spectroscopy in the coming years.
### 2.4 Conventional $`q\overline{q}`$ Mesons
As a background to these various hadronic exotica we have a spectrum of conventional $`q\overline{q}`$ states, which must be identified if we are to isolate non-$`q\overline{q}`$ states. Since many of the light non-$`q\overline{q}`$ states predicted by theorists have masses and quantum numbers that allow confusion with excited $`q\overline{q}`$ states, it is important to establish the light $`q\overline{q}`$ spectrum below 2.5 GeV as completely as possible.
Identification of the $`q\overline{q}`$ and non-$`q\overline{q}`$ states in the spectrum will require that we clarify meson spectroscopy to a mass of at least 2.5 GeV, so that the pattern of glueballs, hybrids and multiquarks can be established through the identification of sufficient examples of each type of state.
There has been impressive experimental progress in the identification of the (presumably $`q\overline{q}`$) light meson spectrum in recent years. In Fig.4 we show the masses of the relevant radially- and orbitally-excited multiplets for which candidate states were reported at the WHS99 hadron spectroscopy meeting in Frascati earlier this year, from a review by Barnes. It appears that almost all the $`q\overline{q}`$ multiplets expected below 2.5 GeV have now been identified. These multiplet masses and some representative candidates reported at the WHS99 meeting are given in Table 1.
Surprisingly, these orbital+radial multiplets lie at rather lower masses than predicted by Godfrey and Isgur; compare the predicted and observed 2P and 2D multiplet masses:
$$M(2P)|_{\mathrm{GI}}1.80\mathrm{GeV},$$
(6)
$$M(2P)|_{\mathrm{expt}.}1.7\mathrm{GeV}.$$
(7)
$$M(2D)|_{\mathrm{GI}}2.14\mathrm{GeV},$$
(8)
$$M(2D)|_{\mathrm{expt}.}2.0\mathrm{GeV}.$$
(9)
Evidently, experiment is finding the 2P and 2D multiplets about 0.1-0.2 GeV lower in mass than predicted by the Godfrey-Isgur model. If this discrepancy is confirmed it will be important to determine whether this requires some important modification of the model.
Thus far it has been possible to identify these $`q\overline{q}`$ multiplets largely by the systematics of masses. This is possible because multiplet splittings decrease rapidly with increasing $`L`$, so we are fortunate to find the members of a given higher-$`L`$ multiplet at very similar masses. In principle one might also distinguish between $`q\overline{q}`$ states and non-$`q\overline{q}`$ exotica such as glueballs and hybrids through their strong decay branching fractions and amplitudes. Detailed predictions are now available for these branching fractions for all $`n\overline{n}`$ states expected up to 2.1 GeV, and for a few specific cases at higher mass. If our decay models are accurate, these higher quarkonia often have very characteristic branching fractions, which should be quite distinct from glueball or hybrid decays. Unfortunately, the <sup>3</sup>P<sub>0</sub> decay model and the closely related flux-tube decay model have not been tested carefully, except in a few cases such as $`b_1\omega \pi `$ and $`a_1\rho \pi `$. (These transitions have both S and D amplitudes, and their D/S ratios are sensitive tests of the decay models and are in good agreement with experiment.) A new and very important test of the decay models was recently reported by VES. In both the <sup>3</sup>P<sub>0</sub> and flux-tube decay models, transitions of the type $`(S_{q\overline{q}}=0)(S_{q\overline{q}}=0)+(S_{q\overline{q}}=0)`$ are forbidden, due to the spin-1 character of the decay model pair creation operator. This implies for example that $`\pi _2(1670)b_1\pi `$ should vanish, although it is nominally an allowed D-wave strong decay. VES finds a rather tight upper limit on this transition,
$$B(\pi _2(1670)b_1\pi )<0.19\%(2\sigma c.l.).$$
(10)
This null result is very reassuring, but does not uniquely confirm a <sup>3</sup>P<sub>0</sub>-type decay model; the same theoretical zero follows for example from OGE pair production. A second test due to VES which also involves the $`\pi _2(1670)`$ does not agree with the expectations of the decay models: $`B(\pi _2(1670)\omega \rho )`$ should be about $`16\%`$, and the spin-1 decay operator implies that the $`\omega \rho `$ final state should have spin-1, with the <sup>3</sup>P<sub>2</sub> $`\omega \rho `$ amplitude dominant. VES instead finds
$$B(\pi _2(1670)\omega \rho (S=2))=1.9(0.4)(1.0)\%$$
(11)
and
$$B(\pi _2(1670)\omega \rho (S=1))=0.9(0.2)(0.3)\%,$$
(12)
which suggests that strong decays in this sector may not agree with the decay models.
Until such time as we can test the predictions of the decay models against a wide range of accurately determined experimental decay amplitudes and branching fractions, it will remain unclear whether the predictions are indeed reliable, or accidentally happen to work well for a few special cases. For this reason it would be extremely useful to determine the relative branching fractions of all two-body modes of higher-mass states such as the excited vectors $`\rho (1465)`$ and $`\rho (1700)`$. The current situation, with most modes unmeasured or reported only as “seen” (Tables 2-4) does not allow one to make progress in the very important subject of strong decays. An accurate determination of excited vector decay amplitudes would be an extremely useful DAFNE contribution, as we shall now discuss.
## 3 Exotica and excited vector mesons at DAFNE
The “vector sector” affords a very interesting subject for future investigation at DAFNE. This topic was studied at Frascati in the past at ADONE, albeit with much lower luminosity. $`e^+e^{}`$ annihilation is of course the ideal technique for making these states, since single photons make $`1^{}`$ uniquely. At the time this appeared to be a rather straightforward problem in hadron spectroscopy, since the quark model predictions of excited $`J^{PC}=1^{}`$ vector mesons with radial and orbital excitations (with both 2<sup>3</sup>S<sub>1</sub> and <sup>3</sup>D<sub>1</sub> expected at about $`1\frac{1}{2}`$ GeV) were uncontroversial. Unfortunately it was found that the excited vectors were rather broad, overlapping states, so the radial and orbital excitations could not be clearly separated. This subject has been reviewed by Donnachie, who discusses it in more detail in these proceedings.
The subject advanced somewhat with studies of the $`\rho `$-sector in both $`2\pi `$ and $`4\pi `$ modes, which lead the PDG in 1988 to distinguish two states, the $`\rho (1450)`$ and the $`\rho (1700)`$. (The current status of light vector spectroscopy according to the PDG is shown in Fig.5.) These are usually indentified with the 2S (radial) and D (orbital) excitations respectively, since the masses correspond approximately to quark potential model expectations. (There are problems with this simple assignment, such as the surprisingly large $`e^+e^{}`$ coupling of the nominally $`L=2`$ $`\rho (1700)`$, which has a vanishing wavefunction at contact.)
There are analogous states reported in the isosinglet sector, the $`\omega (1420)`$ and $`\omega (1600)`$. (The situation may be more complicated. See in particular the recent results on $`e^+e^{}\pi ^+\pi ^{}\pi ^o`$ from VEPP, which show a very low mass peak at about 1220 MeV.) In the $`\varphi `$ sector we have evidence for only a single excitation, the $`\varphi (1680)`$.
Interest in the excited vectors has increased with the realization that the lightest hybrid meson multiplet includes a $`1^{}`$ (in both the flux-tube and bag models), and that these hybrid vectors are predicted to be rather narrow. Indeed, in the hybrid meson decay calculations of Close and Page (using the Isgur-Kokoski-Paton flux-tube model) the narrowest hybrid found was the $`\omega `$-flavor $`1^{}`$. (See Table 3 for the predicted partial widths of this vector hybrid.) The Close-Page calculations assumed a mass of 2.0 GeV for the $`\rho _H`$ and $`\omega _H`$ hybrid vectors, but given the reports of the $`\pi _1(1400)`$ and $`\pi _1(1600)`$ hybrid candidates, one should also consider the possibility that the lowest hybrid multiplet lies at about 1.5 GeV. This would give us a third $`1^{}`$ level roughly degenerate with the quark model 2S and D levels, and such light vector hybrids could be very narrow (see Table 3); a hypothetical $`\omega _H(1500)`$ is predicted to have a total width of only about 20 MeV! If the $`1^{}`$ hybrid states are not found at this low mass, one might question the reports of $`\pi _1`$ $`1^+`$ exotics near 1.5 GeV.
The topic of vector meson spectroscopy in this mass region was recently reviewed by Donnachie and Kalashnikova, who concluded that additional vectors beyond the expected $`q\overline{q}`$ states are indeed required to fit the data in both $`I=0`$ and $`I=1`$ channels. In $`I=1`$ in particular, the weakness of $`e^+e^{}\pi ^+\pi ^{}\pi ^o\pi ^o`$ relative to $`e^+e^{}\pi ^+\pi ^{}\pi ^+\pi ^{}`$ cannot be explained by the expected $`\rho (1700)`$ decays alone.
In principle one should be able to separate the 2S, D and H (hybrid) states by studies of their relative decay branching fractions. In Tables 2-4 we show theoretical predictions for the different types of states, compared to 1998 PDG results for experimental branching fractions. The relative strength of the broad $`4\pi `$ modes $`h_1\pi `$ and $`a_1\pi `$ is quite sensitive to the type of parent resonance, and could serve as a useful discriminator if the decay models are accurate. The theoretical expectation is that the D state should populate both modes, H should only populate $`a_1\pi `$, and the 2S state does not couple significantly to either of these modes. How well do these theoretical predictions agree with experiment?
The experimental branching fractions of excited vectors, as reported in the 1998 PDG, are also shown in Tables 2-4. It is clear at a glance that experiment is in a woeful state. Almost nothing is known about the decays of excited $`\omega `$ and $`\varphi `$ states. (Note that excited $`\varphi `$ vectors can be isolated by studying the $`s\overline{s}`$-filter mode $`\varphi \eta `$, which was apparently not attempted previously.) In the $`\rho `$ sector, there are promising indications that the $`\rho (1700)`$ may be observed in many of the expected channels, but there is almost no quantitative information about relative branching fractions which we require for tests of the decay models. In contrast, there are strong limits claimed for $`\rho (1465)`$ branching fractions, which appear to be very different from expectations for a simple 2S radial excitation. Note especially the tight limit $`B(\rho (1465)\omega \pi )<2\%`$. Taken literally, this result is very interesting in that it argues strongly against a 2S assignment for the $`\rho (1465)`$. (Compare the $`\rho _{2S}`$ and $`\rho (1465)`$ entries in Table 2.) Unfortunately it is difficult to reconcile this number with the reported dominance of $`\omega (1419)\rho \pi `$ (Table 3), since that decay differs from $`\rho (1465)\omega \pi `$ only by a flavor factor of 3 (favoring $`\omega (1419)\rho \pi `$) and minor changes in phase space.
Recently the Crystal Barrel Collaboration attempted to separate the contributions of the $`\rho (1465)`$ and $`\rho (1700)`$ to the various $`4\pi `$ final states. Initially the results appeared consistent with the usual quark model assignments 2S and D, but the most recent work has found that essentially all broad $`4\pi `$ modes ($`a_1\pi ,h_1\pi ,\pi (1300)\pi ,\rho \rho `$ and $`\rho \sigma `$) are important in the decays of both the $`\rho (1465)`$ and the $`\rho (1700)`$! Unfortunately the statistical errors of this many-parameter fit are rather large, so each mode typically has a fitted branching fraction about $`2\sigma `$ from zero. The excited vectors would evidently benefit from a study at an $`e^+e^{}`$ facility such as DAFNE, where the complication of competing amplitudes in many other $`J^{PC}`$ channels is not present.
In view of the poorly constrained and perhaps inconsistent branching fractions evident in the PDG, the most reasonable approach in future would probably be to study as many of the quasi-two-body decay modes in Tables 2-4 as possible, determine numerical values for the relative branching fractions, and carry out a global fit of each flavor sector with an assumed two versus three parent resonances in each flavor.
## 4 Two-photon couplings
In the opinion of at least two LEAR experimentalists, using $`\gamma \gamma `$ collisions to clarify the scalar sector is the most interesting contribution DAFNE could make to spectroscopy.
Two-photon couplings of resonances can be inferred by measurement of the cross section
$$\sigma (e^+e^{}e^+e^{}R)$$
(13)
which is proportional to the two-photon width $`\mathrm{\Gamma }\gamma \gamma `$ of the resonance $`R`$, as discussed in Sec.36.3 of the 1998 PDG. Two-photon widths of $`C=(+)`$ resonances have been measured at several $`e^+e^{}`$ facilities in the past, most recently at LEP. These are especially interesting quantities because they show considerable variation between $`q\overline{q}`$ and non-$`q\overline{q}`$ states, and if determined with sufficient accuracy they could be used for example to solve the problem of the assignments of the various light scalars. This subject attracted considerable interest and effort previously, but as $`e^+e^{}e^+e^{}R`$ is an $`O(\alpha ^4)`$ process and the cross section falls rapidly with $`M_R`$, it was not possible to obtain adequate statistics for a definitive analysis.
The two-photon partial widths of $`q\overline{q}`$ states within a flavor multiplet in the SU(3) limit are in the ratio
$$\mathrm{\Gamma }\gamma \gamma f:a:f^{}=25:9:2,$$
(14)
so if a candidate $`q\overline{q}`$ state such as the $`2^{++}`$ $`f_2(1270)`$ is reported, one should also observe its flavor partners at about this relative strength. For example, the $`\mathrm{\Gamma }\gamma \gamma `$ widths of the $`2^{++}`$ multiplet are
$$\mathrm{\Gamma }\gamma \gamma (2^{++})f_2(1270):a_2(1310):f^{}(1525)=2.8(4)\mathrm{keV}:1.00(6)\mathrm{keV}:0.1\mathrm{keV}.$$
(15)
(Moderate suppression of the $`s\overline{s}`$ coupling is expected theoretically due to the heavier strange quark mass.)
Scalars are predicted to have very characteristic two-photon couplings. The largest $`\mathrm{\Gamma }\gamma \gamma `$ width expected for any $`q\overline{q}`$ meson is for the <sup>3</sup>P<sub>0</sub> $`f_0`$ scalar; in the nonrelativistic quark model it has a $`\mathrm{\Gamma }\gamma \gamma `$ width 15/4 times that of the $`f_2`$, and with relativistic corrections the ratio is reduced to $`2`$. Thus for a scalar $`n\overline{n}`$ partner of the $`f_2(1270)`$ we expect a two-photon width of about 5 keV. An $`f_0(1250)`$ scalar signal of about this strength was observed by the Crystal Ball Collaboration in $`\gamma \gamma \pi ^o\pi ^o`$ at DESY, and may be the long-sought and still obscure $`n\overline{n}`$ scalar. In contrast, a pure scalar glueball should have a much smaller two-photon width, since it has no direct coupling to photons. The recent ALEPH results on $`\gamma \gamma `$ couplings of resonances appear to support the $`f_0(1500)`$ as a glueball candidate, since their upper limit
$$\mathrm{\Gamma }\gamma \gamma (f_0(1500))<0.17\mathrm{keV}(95\%c.l.)$$
(16)
is far below the ca. 5 keV expected for an $`n\overline{n}`$ scalar.
The various $`n\overline{n}Gs\overline{s}`$ mixing models in contrast would predict $`\mathrm{\Gamma }\gamma \gamma `$ widths roughly proportional to each state’s $`n\overline{n}`$ amplitude squared, and so could be tested by the relative strength of each scalar resonance in $`\gamma \gamma \pi ^o\pi ^o`$. Finally, $`K\overline{K}`$ molecules and multiquark states are predicted to have much smaller $`\mathrm{\Gamma }\gamma \gamma `$ widths than the corresponding $`n\overline{n}`$ states, which is in agreement with the sub-keV $`\mathrm{\Gamma }\gamma \gamma `$ values reported for the $`f_0(980)`$ and $`a_0(980)`$.
In contrast with the non-observation of the scalar glueball candidate $`f_0(1500)`$ in $`\gamma \gamma `$, we now have clear evidence for the pseudoscalar $`\eta (1440)`$ in $`\gamma \gamma K_sK^\pm \pi ^{}`$, reported by the L3 Collaboration. Once a glueball candidate (this assignment is now implausible due to the high mass predicted for the pseudoscalar glueball by LGT), this state appears most likely to be a radially-excited $`q\overline{q}`$. Similarly there is a possible observation of the scalar glueball candidate $`f_0(1710)`$ by L3 in $`\gamma \gamma K_sK_s`$, although this is preliminary. If the $`f_0(1710)`$ appears clearly in $`\gamma \gamma `$ at the rate expected for a radially-excited $`2^3`$P<sub>0</sub> $`n\overline{n}`$ state, we may be able to eliminate it as a glueball candidate in favor of the $`f_0(1500)`$. Clearly, accurate measurements of scalar $`\mathrm{\Gamma }\gamma \gamma `$ couplings show great promise as a technique for solving the long standing problem of the nature of the various $`f_0`$ scalar resonances.
## 5 Acknowledgements
It is a pleasure to acknowledge the kind invitation of the organizers of the DAFNE meeting to discuss the status of light meson spectroscopy. I would also like to thank my colleagues for discussions of various aspects of hadron physics in the preparation of this report, in particular N. Achasov, F.E. Close, A. Donnachie, U. Gastaldi, S. Godfrey, N. Isgur, Yu. Kalashnikova, E. Klempt, S. Krewald, A.I.Milstein, C.J. Morningstar, P.R. Page, M.R. Pennington, B. Pick, J. Speth, E.S. Swanson, U. Thoma and N. Törnqvist. Research at the Oak Ridge National Laboratory was supported by the U.S. Department of Energy under contract DE-AC05-96OR22464 with Lockheed Martin Energy Research Corp., and additional support was provided by the Deutsche Forschungsgemeinschaft under contract Bo 56/153-1. |
warning/0001/hep-th0001091.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In statistical mechanics and quantum field theory (QFT), duality plays an important role for exploring strong coupling regime from knowledge of the weak coupling behavior. The sine-Gordon and massive Thirring models are famous examples . In this context, affine Toda field theories (ATFTs) are one-parameter families of quantum integrable massive field theories possessing such duality properties .
Recently, we introduced and studied a new $`n1`$ parameter $`(\{\alpha _a\},\{\beta _a\})`$ family of quantum field theories, called multisine-Gordon (MSG) models. The general construction of the MSG in terms of the extended trigonometry associated to the extended complex numbers , namely multicomplex (MC) numbers, was obtained starting from a generator $`e`$ such that $`e^n=1`$. For $`(\{\alpha _a\},\{\beta _a\})`$, we have shown that these MSG models provide a unifying representation for a wide variety of integrable QFTs possessing dual representations. In particular, integrable models studied a few years ago , like integrable deformations of non-linear $`\sigma `$ models or massive Thirring coupled with ATFTs were recovered subject to restrictions on parameter space.
However, a breakthrough comes in from the writing of the extended trigonometric functions (multisine functions), which appeared in the MSG potentials, in terms of the natural multicomplex extension of vertex operators, namely MC-vertex operators. Using this framework, we investigated the existence of non-local conserved charges. In lower-dimensional QFTs, it is well-known that non-local conserved charges may appear, which generate symmetries characterized by braiding relations . They provide a powerful tool for studying non-perturbative effects . A set of equations associated with the conservation of these non-local charges and their algebraic structure was thus obtained.
In this paper, we investigate the dual relationship between two Lagrangian representations generalizing the sine-Gordon model in the previous meaning. Up to restrictions on the parameter space, we show that lowest-rank ATFTs are generated by MC-algebras. Next, description of various integrable perturbations of conformal field theories as different projections over the usual complex space of the same MC-vertex operator is studied.
In section 2, we briefly recall and reinterpret conveniently some results obtained in which constitute the basic ingredient of the following sections. For $`(\{\alpha _a\},\{\beta _a\}i)`$ (or $`(\{\alpha _a\}i,\{\beta _a\})`$, a generic solution of the equations associated to non-local currents conservation is given. Due to their structure, it naturally emerges a dual relation between the parameters of the model (which appear in the potential through MC-vertex operators) and those involved in the conserved currents. To first order in conformal perturbation theory (CPT), it is then possible to introduce a “dual” family of Lagrangian representations : the “dual” potential is built using the expression in MC-space of conserved currents associated to the original one.
In Section 3, imposing a quantum algebraic structure to the non-local conserved charges, we solve the previous equations. Each solution is associated with a multicomplex space of dimension $`n`$, a specific MC-algebra and different ratios of the parameters. For $`(\{\alpha _a\},\{\beta _a\}i)`$, the underlying hidden symmetry of MSG model in each case is identified to a quantum universal enveloping algebra (QUEA) based on an affine Lie algebra $`\widehat{𝒢}`$. This approach provides a unifying “parametrized” description of lowest-rank QUEAs based on $`A_r^{(1)}`$ for $`r3`$, $`A_{2r}^{(2)}`$ for $`r2`$, $`D_4^{(1)}`$, $`(B_r^{(1)},A_{2r1}^{(2)})`$ for $`r3`$, $`(C_2^{(1)},D_3^{(2)})`$ and $`(G_2^{(1)},D_4^{(3)})`$.
We show in section 4 how ATFTs are related to MSG models. In particular, simple relations between the multicomplex dimension $`n`$ and the rank of $`\widehat{𝒢}`$, the kind of MC-algebra (characterized by $`m_a`$) and the Kac labels (denoted $`n_a`$) are obtained. For $`n=3`$ and $`n=4`$, we describe the (dual-)MC-algebras generating $`A_2^{(1)}`$, $`C_2^{(1)}`$, $`D_3^{(2)}`$ and $`A_3^{(1)}`$ ATFTs.
In section 5, we show that whereas the multisine-potential do not depend on the projections of MC-algebras over the usual complex space, its understanding in terms of perturbed conformal field theory (CFT) does. We identify two kinds of perturbed CFT through the introduction of a (real) multicomplex charge at infinity. For each projection, a “parametrized” central charge is computed.
Some conclusions and perspectives are drawn in section 6.
## 2 Dual conserved currents in extended sine-Gordon
In , we introduced the natural extension of the sine-Gordon field theory in the $`n`$-dimensional multicomplex space. This model, generated by the fundamental multicomplex number <sup>1</sup><sup>1</sup>1More precisely, $`e`$ is denoted $`e_{(n|m)}`$ in ref. due to its eventual substructure. $`e`$ , such that $`e^n=1`$, describes a family of $`n1`$ parameter quantum field theories with $`n1`$ scalar fields which interact through a multisine potential . Its Euclidian action can generally be expressed in terms of the extension in MC-space of standard vertex operators :
$`𝒜^{(n|m)}(\eta )={\displaystyle \frac{1}{4\pi }}{\displaystyle d^2z_z\mathrm{\Phi }_{\overline{z}}\mathrm{\Phi }}+{\displaystyle \frac{\lambda }{n\pi }}{\displaystyle d^2z\left(x^{(0)}+\mathrm{}+x^{(n1)}\right)}`$ (2.1)
with :
$`x^{(l)}=\mathrm{exp}(\eta ^{(l)}.\mathrm{\Phi }(z,\overline{z})),`$ (2.2)
where we define :
$`m=2{\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}m_a\text{for }n\text{ even and }m=2{\displaystyle \underset{a=0}{\overset{\frac{n1}{2}1}{}}}m_a+m_{(n1)/2}\text{for }n\text{ odd}`$ (2.3)
characterize the kind of MC-algebra and $`\mathrm{\Phi }(z,\overline{z})`$ is the fundamental ($`n1`$-components)-field of the theory. Note that the index $`l`$ denotes the $`l`$-th multicomplex conjugation of any MC-number. In , we considered only unimodular MC-numbers. The unimodularity condition was implemented on MC-vertex operators (2.2) through :
$`x^n={\displaystyle \underset{l=0}{\overset{n1}{}}}x^{(l)}=1,`$ (2.4)
where $`||||`$ is the pseudo-norm associated to MC-algebra (see refs. for details). Depending on the value of $`n`$ (even or odd cases), expression (2.2) differs. For $`n`$ even, the MC-vertex operator (2.2) is defined as follow. Firstly, we introduce the bilinear relation :
$`.:\left(𝕄_{(n|m)}\right)^{n1}\times \left(𝕄_{(n|m)}\right)^{n1}𝕄_{(n|m)}`$ (2.5)
which can be understood as an extension of the standard scalar product in the multicomplex valued vector space. MC-vertex operators depend on parameters $`(\{\alpha _a\},\{\beta _a\})`$ through the relation :
$`\eta ^{(l)}`$ $`=`$ $`[\alpha _0P_l,\mathrm{},\alpha _{\frac{n}{2}1}P_{\frac{n}{2}1+l};\beta _0Q_{0;l},\mathrm{},\beta _{\frac{n}{2}2}Q_{\frac{n}{2}2;l}]\left(𝕄_{(n|m)}\right)^{n1},`$ (2.6)
where the parameters $`(\{\alpha _a\},\{\beta _a\})`$ and $`\{P_{a+l},Q_{a;l}\}`$ generate the MC-algebra <sup>2</sup><sup>2</sup>2For $`n`$ even and “mimimal” representations (i.e $`m_a=1`$ for all $`a`$), the generators $`P_a`$ are expressed in terms of the fundamental multicomplex element $`e`$ as : $`P_a=\frac{2}{n}_{j=0}^{n/21}\mathrm{sin}[(2a+1)j\frac{\pi }{n}]e^j`$. . Consequently, if $`\varphi (z)`$ denotes the holomorphic part of the fundamental field $`\mathrm{\Phi }(z,\overline{z})`$, e.g $`\mathrm{\Phi }(z,\overline{z})=\varphi (z)+\overline{\varphi (z)}`$, expression (2.2) is defined (since $`𝕄_{(n|m)}`$) with :
$`\eta ^{(l)}.\varphi (z)={\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}\left[\alpha _aP_{a+l}\varphi _a(z)\right]+{\displaystyle \underset{a=0}{\overset{\frac{n}{2}2}{}}}\left[\beta _aQ_{a;l}\phi _a(z)\right]`$ (2.7)
and :
$`Q_{a;l}=P_{a+l}^2+{\displaystyle \frac{m_a}{m_{n/21}}}P_{n/21+l}^2,Q_{a;l}=Q_{a;l+n/2}=Q_{a+n/2;l}=Q_{a+n;l}=Q_{a;l+n},`$
for $`a\{0,\mathrm{},n/22\}`$, $`l\{0,\mathrm{},n/21\}`$ with the conventions $`\alpha _{a+n/2}=\alpha _a,\beta _{a+n/2}=\beta _a`$ and $`m_{a+n/2}=m_a`$.
In , the MSG models were studied for real parameters. Let us here focus on the parameter space restricted to :
$`\alpha _a\text{and}\beta _ai,`$ (2.8)
or $`\alpha _ai\text{and}\beta _a,`$
for all $`a`$. In these cases, we introduce the dual multicomplex element $`\eta ^{(k)}=2\eta ^{(k)}(\eta ^{(k)}.\eta ^{(k)})^1`$ (supposing that $`(\eta ^{(k)}.\eta ^{(k)})`$ is invertible which will be always satisfied in the sequel). Using MC-algebra and equation (2.6), it leads to the dual parameter space :
$`\alpha _a^{}={\displaystyle \frac{2\alpha _a}{(\alpha _{a}^{}{}_{}{}^{2}\beta _a^2)}}\text{for}a\{0,\mathrm{},n/22\},`$
$`\beta _a^{}={\displaystyle \frac{2\beta _a}{(\alpha _a^2\beta _a^2)}}\text{for}a\{0,\mathrm{},n/22\},`$ (2.9)
$`\alpha _{n/21}^{}={\displaystyle \frac{2\alpha _{n/21}}{(\alpha _{n/21}^2_{a=0}^{n/22}\frac{m_a^2}{m_{\frac{n}{2}1}^2}\beta _a^2)}},`$
and the dual MC-algebra generated by $`e^{}`$ with :
$`{\displaystyle \frac{m_a^{}}{m_{\frac{n}{2}1}^{}}}={\displaystyle \frac{m_a}{m_{\frac{n}{2}1}}}{\displaystyle \frac{(\alpha _a^2\beta _a^2)}{(\alpha _{n/21}^2_{a=0}^{n/22}\frac{m_a^2}{m_{\frac{n}{2}1}^2}\beta _a^2)}}\text{for}a\{0,\mathrm{},n/22\}`$ (2.10)
In the following, it is then more convenient to rewrite the action of the MSG (2.1) in MC-space as :
$`𝒜^{(n|m)}(\eta )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle d^2z_z\mathrm{\Phi }_{\overline{z}}\mathrm{\Phi }}+{\displaystyle \frac{\lambda }{2\pi }}{\displaystyle d^2z\mathrm{\Phi }_{pert}(\eta )}.`$ (2.11)
By analogy, we can now consider the MSG model with action $`𝒜^{(n|m^{})}(\eta ^{})`$, where $`m^{}`$ is defined as in eq. (2.3) with eqs. (2.10). The perturbing operator for each model reads respectively :
$`\mathrm{\Phi }_{pert}(\eta )={\displaystyle \frac{2}{n}}{\displaystyle \underset{l}{}}J_{\eta ^{(l)}}\overline{J}_{\eta ^{(l)}}\text{and}\mathrm{\Phi }_{pert}(\eta ^{})={\displaystyle \frac{2}{n}}{\displaystyle \underset{l}{}}J_{\eta ^{(l)}}\overline{J}_{\eta ^{(l)}}.`$ (2.12)
Considering now action (2.1) as a perturbed $`(n1)`$ free field CFT and following Zamolodchikov approach , it is possible to construct $`2n`$ non-local conserved charges to first order in conformal perturbation theory (CPT). Let us consider the holomorphic MC-vertex operator $`J^{(k)}=J_{\eta _{}^{}{}_{}{}^{(k)}}=e^{\eta _{}^{}{}_{}{}^{(k)}.\varphi (z)}`$ (and respectively, for the antiholomorphic part, $`\overline{J}^{(k)}=\overline{J}_{\eta _{}^{}{}_{}{}^{(k)}}=e^{\eta _{}^{}{}_{}{}^{(k)}.\overline{\varphi (z)}}`$). Let us also suppose that its OPE with the $`k^{th}`$-perturbing term of the potential leads to a derivative term : $`\overline{}J^{(k)}=H^{(k)}`$ and similarly for the antiholomorphic part (whereas OPEs with any other $`lk^{th}`$-perturbing term yields to regular terms), e.g. the OPE reads :
$`J_{\eta _{}^{}{}_{}{}^{(k)}}(z)x^{(l)}(w)(zw)^{C^{k,l}(\eta ^{},\eta )}e^{(\eta _{}^{}{}_{}{}^{(k)}+\eta ^{(l)}).\varphi (w)}+\mathrm{}`$ (2.13)
with $`\eta _{}^{}{}_{}{}^{(k)}`$ defined as in (2.6) but with $`\alpha _a\alpha _a^{}`$, $`\beta _a\beta _a^{}`$ and $`m_am_a^{}`$. Since MC-exponents $`C^{k,l}(\eta ^{},\eta )=\eta _{}^{}{}_{}{}^{(k)}.\eta ^{(l)}`$ (reported in ) are of the form $`C^{k,l}(\eta ^{},\eta )=_{a=0}^{\frac{n}{2}1}C_a^{k,l}(\eta ^{},\eta )(P_a^2)`$ in terms of MC-algebra, where $`C_a^{k,l}(\eta ^{},\eta )`$ depend on $`\alpha _a`$, $`\beta _a`$, $`\alpha _a^{}`$, $`\beta _a^{}`$, then using MC-algebra, one easily shows that :
$`J_{\eta ^{{}_{}{}^{}(k)}}(z)x^{(l)}(w)\left[{\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}(zw)^{C_a^{k,l}(\eta ^{},\eta )}(P_a^2)\right]e^{(\eta _{}^{}{}_{}{}^{(k)}+\eta ^{(l)}).\varphi (w)}+\mathrm{}`$ (2.14)
It clearly implies that the conservation condition of the holomorphic current $`J^{(k)}`$ does not depend on the choice of any multicomplex representation. It remains to solve :
$`C_a^{k,k}(\eta ^{},\eta )=2,C_a^{k,l}(\eta ^{},\eta )_{}\text{for}kl,\text{for all}a,`$ (2.15)
with $`(k,l)\{0,\mathrm{},n/21\}`$. Whereas this system of constraints looks overdeterminate, invariance under multicomplex conjugation of action (2.1) leads to redundancies. In parameter space, the constraints (2.15) reads :
$`\left(\alpha _a^{}\alpha _a+\beta _a^{}\beta _a\right)=2,\left(\alpha _{\frac{n}{2}1}^{}\alpha _{\frac{n}{2}1}+{\displaystyle \underset{a=0}{\overset{\frac{n}{2}2}{}}}\beta _a^{}\beta _a{\displaystyle \frac{m_a^{}m_a}{m_{\frac{n}{2}1}^{}m_{\frac{n}{2}1}}}\right)=2,`$ (2.16)
$`\left(\alpha _a^{}\alpha _a+\beta _a^{}\beta _a\right)_{},\left(\alpha _{\frac{n}{2}1}^{}\alpha _{\frac{n}{2}1}+{\displaystyle \underset{a=0}{\overset{\frac{n}{2}2}{}}}\beta _a^{}\beta _a{\displaystyle \frac{m_a^{}m_a}{m_{\frac{n}{2}1}^{}m_{\frac{n}{2}1}}}\right)_{},`$ (2.17)
$`{\displaystyle \frac{m_a}{m_{\frac{n}{2}1}}}\beta _a^{}\beta _a_+^{}\text{and}{\displaystyle \frac{m_a^{}}{m_{\frac{n}{2}1}^{}}}\beta _a^{}\beta _a_+^{},`$ (2.18)
for all $`a\{0,\mathrm{},n/22\}`$. If eqs. (2.16), (2.17) and (2.18) are satisfied, the non-local currents $`J^{(k)}`$ are conserved to first order in CPT for all $`k`$ (and similarly for the antiholomorphic part). They generate $`2n`$ non-local conserved charges :
$`Q^{(k)}`$ $`=`$ $`{\displaystyle \frac{1}{2i\pi }}\left({\displaystyle _z}𝑑zJ^{(k)}+{\displaystyle _{\overline{z}}}𝑑\overline{z}H^{(k)}\right),`$
$`\overline{Q}^{(k)}`$ $`=`$ $`{\displaystyle \frac{1}{2i\pi }}\left({\displaystyle _{\overline{z}}}𝑑\overline{z}\overline{J}^{(k)}+{\displaystyle _z}𝑑z\overline{H}^{(k)}\right)\text{for}k[0,\mathrm{},n1].`$ (2.19)
The non-locality is due to the fact that (anti-)chiral components $`\varphi `$ and $`\overline{\varphi }`$ of the MSG model are non-local with respect to the fundamental field $`\mathrm{\Phi }`$. From this property, braiding relations may arised between these charges. It is now obvious to see that the parameters $`\alpha _a^{}=\alpha _a^{}`$$`\beta _a^{}=\beta _a^{}`$ and $`m_a^{}/m_{n/21}^{}=m_a^{}/m_{n/21}^{}`$ satisfy eqs. (2.16), whereas eqs. (2.17), (2.18) still constraint the parameter space $`(\{\alpha _a\},\{\beta _a\})`$ of MSG (2.1). Under these additional conditions on parameter space, conservation of $`J_{\eta _{}^{}{}_{}{}^{(k)}}=J_{\eta ^{(k)}}`$ in the model associated to action $`𝒜^{(n|m)}(\eta )`$ is ensured. It possesses $`2n`$ non-local conserved currents $`\{J_{\eta ^{(k)}},\overline{J}_{\eta ^{(k)}}\}`$. Consequently, since $`C^{k,l}(\text{(}\eta ^{}\text{)}^{},\eta ^{})=C^{l,k}(\eta ^{},\eta )`$, the whole set of non-local currents $`\{J_{\eta ^{(k)}},\overline{J}_{\eta ^{(k)}}\}`$ are similarly conserved to first order in CPT, in the model $`𝒜^{(n|m^{})}(\eta ^{})`$. It allows to define a duality relation between these two models, at least to order $`𝒪(\lambda )`$. Note that the transformations (2.9), (2.10) involving $`n1`$ parameters, resulting of the current conservation, is completly independent of any particular structure of non-local conserved charge algebra.
However, to define a consistent QFT, we have to consider carefully the renormalization group flows in such models. The crucial point is that MC-vertex operators $`J_{\eta ^{(k)}}`$ (and resp. $`\overline{J}_{\eta ^{(k)}}`$) do not form a closed algebra by themselves for general values of parameters. If this condition is not satisfied (which is generally the case), renormalization requires that counterterms have to be added in such a way that this algebra closes. Consequently, the non-local charges are obviously no longer symmetries of this modified action. A well-known example is provided by simply-laced ATFTs with parameter $`\beta `$ : for smaller values than the critical value $`\beta ^2=1`$ , only tadpole renormalization is necessary. However, in the Kosterliz-Thouless region, ATFTs are $`\mathrm{𝑛𝑜𝑡}`$ renormalizable and exponential of $`\mathrm{𝑎𝑙𝑙}`$ the roots <sup>3</sup><sup>3</sup>3For non-simply laced case, exponential operators associated to short or long roots have drastically different dimensions : one must introduce fermions in such a way as to increase the conformal dimensions of the exponential operators associated to short roots . must be added in order to render them renormalizable. This situation will be relevant in further analysis.
## 3 Quantum algebraic structure and lowest-rank affine Lie algebras parametrization
Known integrable models generally exhibit connections with Hopf algebras like Lie algebras and their quantum deformations. The aim of this section is to clarify in which sense the previous parameter space of the multisine-Gordon model is restricted by imposing this kind of structure to the non-local conserved charge algebra. The principal motivation to find such a structure in the MSG models comes from the powerful framework that these algebras provide : they can be sufficiently restrictive to allow a non-perturbative solution of the theory, determining the $`S`$ matrices for instance . From eqs. (2.15), i.e. (2.16), (2.17) and (2.18) with solutions (2.9), (2.10), it is convenient to introduce the elements :
$`C_{a,a}(\eta ^{},\eta )=2,C_{a,a+\frac{n}{2}}(\eta ^{},\eta )=C_{a+\frac{n}{2},a}=2{\displaystyle \frac{A_a+B_a}{A_aB_a}}=n_1^a`$
$`C_{a,\frac{n}{2}1}(\eta ^{},\eta )=C_{a+\frac{n}{2},n1}=C_{a+\frac{n}{2},\frac{n}{2}1}=C_{a,n1}=2{\displaystyle \frac{B_a\delta _a}{A_aB_a}}=n_2^a`$ (3.20)
$`C_{\frac{n}{2}1,a}(\eta ^{},\eta )=C_{n1,a+\frac{n}{2}}=C_{\frac{n}{2}1,a+\frac{n}{2}}=C_{n1,a}=2{\displaystyle \frac{B_a\delta _a}{A_{\frac{n}{2}1}B}}=n_3^a`$
$`C_{n/21,n1}(\eta ^{},\eta )=C_{n1,\frac{n}{2}1}=2{\displaystyle \frac{A_{\frac{n}{2}1}+B}{A_{\frac{n}{2}1}B}}=n_4`$
with $`A_a=\alpha _a^2,B_a=\beta _a^2`$, $`\delta _a=\frac{m_a}{m_{\frac{n}{2}1}}`$, and :
$`B={\displaystyle \underset{a=0}{\overset{\frac{n}{2}2}{}}}B_a\delta _a^2.`$ (3.21)
As detailed in , the braiding relations between $`(J^{(k)},\overline{J}^{(l)})`$, $`(J^{(k)},\overline{H}^{(l)})`$ and $`(H^{(k)},\overline{J}^{(l)})`$ arising from the non-local property of the currents are identical iff :
$`\{n_1^a,n_4\}\text{and}\{n_2^a,n_3^a\}^{}\text{for}a\{0,1,\mathrm{},{\displaystyle \frac{n}{2}}2\}.`$ (3.22)
Under this assumption, to first order in CPT non-local conserved charges (2.19) obey a $`q`$-deformed structure in MC-space . Furthermore it is possible, in the usual complex space and independently of any representation, to define one other basis of non-local conserved currents; the expansion of the currents $`J^{(k)}`$ in MC-basis $`\{P_a,P_a^2\}`$ reads :
$`J_{\eta _{}^{(k)}{}_{}{}^{}}={\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}\left[\mathrm{sin}(\alpha _a^{}\varphi _a)e^{\beta _a^{}\phi _a}P_{a+k}+\mathrm{cos}(\alpha _a^{}\varphi _a)e^{\beta _a^{}\phi _a}\left(P_{a+k}^2\right)\right]`$ (3.23)
with $`\beta _{\frac{n}{2}1}^{}\phi _{\frac{n}{2}1}=_{a=0}^{\frac{n}{2}2}\frac{m_a^{}}{m_{\frac{n}{2}1}^{}}\beta _a^{}\phi _a`$. Since $`J_{\eta ^{(k)}}`$ is conserved, then each one of its components (and any linear combinations of them) is a non-local conserved current, particulary :
$`𝒥^{(a)}=e^{i\alpha _a^{}\varphi _a+\beta _a^{}\phi _a},`$ (3.24)
$`𝒥^{(a+\frac{n}{2})}=e^{i\alpha _a^{}\varphi _a+\beta _a^{}\phi _a}`$
and similarly for the antiholomorphic part. Non-local conserved charges $`(𝒬_a,\overline{𝒬}_a)`$ associated to these conserved currents can then be obtained as in (2.19). Analogously to the multicomplex case, these charges obey to a $`q`$-deformed algebra iff eqs. (3.20), (3.21) with (3.22) are satisfied. For $`\alpha _a`$ and $`\beta _ai`$, the resulting structure <sup>4</sup><sup>4</sup>4For this parameter space, the normalization of the field is chosen such that $`𝒯^{(a)}`$ takes integer values . is nothing else than a “parametrized” quantum universal envelopping algebra $`𝒰_q(\widehat{𝒢})`$ :
$`𝒬^{(a)}\overline{𝒬}^{(b)}q_a^{C_{a,b}(\eta ,\eta ^{})}\overline{𝒬}^{(b)}𝒬^{(a)}`$ $`=`$ $`\delta _{a,b}{\displaystyle \frac{\lambda }{in\pi }}\left[1q_a^{2𝒯^{(a)}}\right].`$ (3.25)
where $`q_a=\mathrm{exp}(i\frac{\pi }{2}C_{a,a}(\eta ^{},\eta ^{}))`$ is the deformation and $`𝒯^{(a)}`$ is a parametrized topological charge . Here, coefficients $`C_{a,b}(\eta ,\eta ^{})=C_{b,a}(\eta ^{},\eta )`$ given by eqs. (3.20) correspond to the extended Cartan matrix elements of this “parametrized” $`q`$-deformed algebra. As we will see later, the type of affine Lie algebra is encoded in the ratios of the parameters and the MC-algebra structure (e.g. $`A_a`$, $`B_a`$ and $`\delta _a`$). As was shown in , the matrix elements of the extended Cartan matrix for $`n`$ odd can be obtained similarly from eqs. (3.20) by the substitution :
$`nn+1;A_{\frac{n1}{2}}0;\delta _a2\delta _a,`$ (3.26)
while the last equation of (3.20) does not appear. It is more convenient to formulate this substitution in terms of the $`n_i^a`$ so that solutions for the $`n`$ odd case can be directly read off from solutions of the even case (with $`A_{\frac{n}{2}1}=0`$). It reads :
$`nn+1;n_42;n_2^a{\displaystyle \frac{n_2^a}{2}};n_3^a{\displaystyle \frac{n_3^a}{2}},`$ (3.27)
and (3.21) becomes for the odd case :
$`B=4{\displaystyle \underset{a=0}{\overset{\frac{n1}{2}}{}}}B_a\delta _a^2.`$ (3.28)
In the following, the system (3.20) (and correspondingly for $`n`$ odd) is solved <sup>5</sup><sup>5</sup>5For completeness, we also report the cases $`C_2^{(1)}`$, $`D_3^{(2)}`$, $`A_2^{(1)}`$ and $`A_3^{(1)}`$ already obtained in .. First, using (3.20), $`A_a`$, $`\delta _a`$, $`B_a`$, $`A_{\frac{n}{2}1}`$ are expressed in terms of $`n_i^a`$ and $`B`$. Then, equation (3.21), (or (3.28) for $`n`$ odd) appear as a further constraint that fix the $`n_i^a`$. Two generic cases are now studied : $`i`$) no restriction $`(\{\alpha _a\},\{\beta _a\})0`$ , $`ii`$) restrictions on parameters space.
Case $`n`$ even without restriction
The case with no restrictions corresponds to keeping $`A_a0`$, $`B_a0`$, $`A_{\frac{n}{2}1}0`$. The general solution for the system (3.20) reads (with $`B<0`$) :
$`{\displaystyle \frac{A_a}{B_a}}={\displaystyle \frac{2+n_1^a}{2n_1^a}};`$ $`A_{\frac{n}{2}1}={\displaystyle \frac{2+n_4}{2n_4}}B;`$ (3.29)
$`B_a={\displaystyle \frac{n_3^a(2n_1^a)}{n_2^a(2n_4)}}B;`$ $`\delta _a={\displaystyle \frac{2n_2^a}{2n_1^a}},`$
with $`n_1^a=0`$ or $`1`$, $`n_2^a>0`$, $`n_3^a>0`$ and $`n_4=0`$ or $`1`$. Since $`n_i^a^{}`$, $`\delta _a`$ are rational and the $`m_a`$ can then be integers. Here, $`B`$ is a negative real number, and due to (3.21), the $`n_i^a`$ must satisfy :
$`{\displaystyle \underset{a=0}{\overset{\frac{n}{2}2}{}}}{\displaystyle \frac{n_3^an_2^a}{2n_1^a}}={\displaystyle \frac{2n_4}{4}}.`$ (3.30)
Since $`\frac{n_3^an_2^a}{2n_1^a}\frac{1}{2}^{}`$, $`n_4=0`$ is the only consistent case. It corresponds to the line $`d/`$ of Table 1.1.
Case $`n`$ odd without restriction
The solution in this case can be obtained from (3.29) with substitutions (3.27) ($`B<0`$):
$`{\displaystyle \frac{A_a}{B_a}}={\displaystyle \frac{2+n_1^a}{2n_1^a}};B_a={\displaystyle \frac{n_3^a(2n_1^a)}{4n_2^a}}B;\delta _a={\displaystyle \frac{n_2^a}{2n_1^a}},`$ (3.31)
with $`n_1^a=0`$ or $`1`$, $`n_2^a>0`$, $`n_3^a>0`$, and the condition (3.28) :
$`{\displaystyle \underset{a=0}{\overset{\frac{n1}{2}1}{}}}{\displaystyle \frac{n_3^an_2^a}{2n_1^a}}=1.`$ (3.32)
Bounds on the $`n_i^a`$ imply that (3.32) has solutions only if $`n5`$. There are three solutions for $`n=3`$ corresponding to lines a/, b/, c/ of Table 1.1, and one for $`n=5`$ (line e/ of Table 1.1).
| | $`n`$ | a | $`n_1^a`$ | $`n_2^a`$ | $`n_3^a`$ | $`A_a/B_0`$ | $`A_{\frac{n}{2}1}/B_0`$ | $`B_a/B_0`$ | $`\delta _a`$ | Algebra |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| a/ | 3 | 0 | 0 | 1 | 2 | -1 | * | 1 | 1/2 | $`C_2^{(1)}`$ |
| b/ | 3 | 0 | 0 | 2 | 1 | -1 | * | 1 | 1 | $`D_3^{(2)}`$ |
| c/ | 3 | 0 | 1 | 1 | 1 | -3 | * | 1 | 1 | $`A_2^{(1)}`$ |
| d/ | 4 | 0 | 0 | 1 | 1 | -1 | -1 | 1 | 1 | $`(n_4=0)A_3^{(1)}`$ |
| | | 0 | 0 | 1 | 1 | -1 | | 1 | 1/2 | |
| e/ | 5 | | | | | | * | | | $`D_4^{(1)}`$ |
| | | 1 | 0 | 1 | 1 | -1 | | 1 | 1/2 | |
Table 1.1 \- Solutions without restriction, and the corresponding
hidden symmetry algebras.
Figure 1.1 \- Small arrows symbolise the exchange under multicomplex conjugation
$`aa+n/2`$ (or $`aa+(n+1)/2`$ for $`n`$ odd) . The number of links is
$`\left|C_{a,b}(\eta ,\eta ^{})\right|\left|C_{a,b}(\eta ^{},\eta )\right|`$ where big arrows goes from $`a`$ to $`b`$ if $`|C_{a,b}(\eta ,\eta ^{})|>|C_{a,b}(\eta ^{},\eta )|`$.
Case $`n`$ even with restrictions
We consider $`s+1`$ restrictions :
$`A_a`$ $`=`$ $`0\text{ for }a\{0,1,\mathrm{},s\},`$
$`A_a`$ $``$ $`0\text{ for }a\{s+1,\mathrm{},{\displaystyle \frac{n1}{2}}1\}(\mathrm{}\text{ if }s={\displaystyle \frac{n1}{2}}1).`$ (3.33)
The solution of (3.20) for $`as`$ can be obtained from (3.29) by setting $`n_1^a=2`$ :
$`A_a=0;A_{\frac{n}{2}1}={\displaystyle \frac{2+n_4}{2n_4}}B;B_a={\displaystyle \frac{4n_3^a}{n_2^a(2n_4)}}B;\delta _a={\displaystyle \frac{n_2^a}{2}},`$ (3.34)
while (3.29) still holds when $`s+1a\frac{n}{2}1`$. One can also set the restriction $`A_{\frac{n}{2}1}=0`$, in which case solutions of (3.20) are obtained from (3.29) and (3.34) by setting $`n_4=2`$, (while $`n_4=0`$ or $`1`$ if $`A_{\frac{n}{2}1}0`$). In all cases, one still have $`n_2^a>0,n_3^a>0`$. (3.21) yields the condition :
$`{\displaystyle \underset{a=0}{\overset{s}{}}}{\displaystyle \frac{n_2^an_3^a}{4}}+{\displaystyle \underset{a=s+1}{\overset{\frac{n}{2}2}{}}}{\displaystyle \frac{n_2^an_3^a}{2n_1^a}}={\displaystyle \frac{2n_4}{4}}.`$ (3.35)
$``$ For $`A_{\frac{n}{2}1}0`$, this last equation admits solutions only if $`n6`$ : For $`n=4`$, there are three solutions written on lines a/, b/, c/ of Table 1.2, and for $`n=6`$ there is one solution line n/ of Table 1.2, reported in Appendix A.
$``$ For $`A_{\frac{n}{2}1}=0`$, the upper bound on $`n`$ is $`10`$. Solutions are written on lines d/ to m/ and o/ to r/ of Table 1.2 in Appendix A.
Case $`n`$ odd with restrictions
In this case, one can only consider $`s+1`$ restrictions :
$`A_a`$ $`=`$ $`0\text{ for }a\{0,1,\mathrm{},s\},`$
$`A_a`$ $``$ $`0\text{ for }a\{s+1,\mathrm{},{\displaystyle \frac{n1}{2}}\}(\mathrm{}\text{ if }s={\displaystyle \frac{n1}{2}}).`$ (3.36)
Solutions for $`as`$ can be obtained from (3.31) by setting $`n_1^a=2`$ :
$`A_a=0;B_a={\displaystyle \frac{n_3^a}{n_2^a}}B;\delta _a={\displaystyle \frac{n_2^a}{4}},`$
while for $`s+1a\frac{n1}{2}`$ solution is given by (3.31), with $`n_1^a=0`$ or $`1`$. In any case $`n_2^a>0`$ and $`n_3^a>0`$. The condition (3.28) becomes :
$`{\displaystyle \underset{a=0}{\overset{s}{}}}{\displaystyle \frac{n_2^an_3^a}{4}}+{\displaystyle \underset{s+1}{\overset{\frac{n1}{2}1}{}}}{\displaystyle \frac{n_2^an_3^a}{2n_1^a}}=1.`$ (3.37)
Owing to the conditions on $`n_i^a`$, and the upper bound of $`s`$, (3.37) have only solutions for $`n9`$. Solutions are summarized on lines a/ to n/ of Table 1.3 in Appendix A.
To resume, imposing a quantum algebraic structure to the non-local charges restricts the ($`n1`$)-parameter space to a discrete set of one dimensional submanifolds. Each one is characterized by ratios of the parameters ($`\{\alpha _a\},\{\beta _a\}`$) and ratios of $`m_a`$ which determinate $`n_1^a,n_2^a,n_3^a,n_4`$ in eqs. (3.20). Consequently, since (3.20) are identified to Cartan matrix elements, any of these submanifolds is in one-to-one correspondance with an affine Lie algebra (see tables. 1.1-1.3). We see that any lowest-rank affine Lie algebra and its dual appear for each fixed value of the multicomplex space of dimension $`n`$. This is easily understood from the substitutions (2.9) and (2.10) in (3.25), resulting from the duality property.
## 4 Integrability in parameter space : affine Toda field theories and perturbed WZNW models
Toda and affine Toda field theories are generally understood as the simple Lie group extension of the Liouville and sine-Gordon models respectively. From our previous results, we intend to show in this section that some of these theories, more precisely those based on lowest rank simply-laced algebras, can be generated starting from a multicomplex number $`e`$ satisfying $`e^{10}=1`$. Let us now consider the set of MC-vertex operators :
$`y^{(l)}=n^{(l)}x^{(l)}\text{with}n^{(l)}={\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}n_a(P_{a+l}^2).`$ (4.38)
The pseudo-norm associated to these operators is :
$`y^m=nx^m={\displaystyle \underset{a=0}{\overset{\frac{n}{2}1}{}}}n_a^{2m_a}`$ (4.39)
as $`x^m=1`$. It is then straightforward to see that among the models which can be built in terms of these MC-vertex operators, the ones with action (for $`n`$ even) :
$`𝒜^{(n|m)}(\eta )={\displaystyle \frac{1}{4\pi }}{\displaystyle d^2z_z\mathrm{\Phi }_{\overline{z}}\mathrm{\Phi }}+{\displaystyle \frac{\lambda }{n\pi }}{\displaystyle d^2z\underset{k=0}{\overset{n1}{}}n^{(k)}x^{(k)}}`$ (4.40)
possess exactly the same underlying algebraic structure than those considered in previous sections. Except for the explicit expression of the non-local conserved charges and r.h.s of (3.25), which are modified by the presence of the extra-parameters $`n_a`$ appearing through the changes in the non-diagonal terms :
$`H^{(k)}n^{(k)}H^{(k)},`$ (4.41)
in (2.19), the whole analysis concerning the non-local conserved charge algebraic structure is preserved. In the usual complex space (4.40) for $`n`$ even writes :
$`𝒜^{(n|m)}(\eta )={\displaystyle \frac{1}{4\pi }}{\displaystyle d^2z_z\mathrm{\Phi }_{\overline{z}}\mathrm{\Phi }}+{\displaystyle \frac{2\lambda }{n\pi }}{\displaystyle d^2z\underset{a=0}{\overset{\frac{n}{2}1}{}}n_a\mathrm{cos}(\alpha _a\varphi _a)e^{\beta _a\phi _a}}`$ (4.42)
with $`\beta _{\frac{n}{2}1}\phi _{\frac{n}{2}1}=_{a=0}^{\frac{n}{2}2}\frac{m_a}{m_{\frac{n}{2}1}}\beta _a\phi _a`$. This action can be generally put into the form :
$`𝒜^{(n|m)}(\eta )={\displaystyle \frac{1}{4\pi }}{\displaystyle d^2z_z\mathrm{\Phi }_{\overline{z}}\mathrm{\Phi }}+{\displaystyle \frac{\lambda }{n\pi }}{\displaystyle d^2z\left[\underset{a=0}{\overset{n1}{}}n_ae^{\beta _0𝐫_a.\mathrm{\Phi }}\right]}.`$ (4.43)
with $`n_a=n_{a+\frac{n}{2}}`$ for $`a\{0,\mathrm{},\frac{n}{2}1\}`$, and where we introduce some “parametrized” roots $`𝐫_a`$, reported in Appendix B.
Similarly, to obtain the multisine-Gordon action for $`n`$ odd , we do successively the substitutions (3.26) in (4.40), (4.42) and (4.43), then change $`n_{\frac{n1}{2}}n_{\frac{n1}{2}}/2`$. For $`n`$ odd, “parametrized” roots are deduced using the same method. As we see, the form of the action recalls the standard one of affine Toda field theories. Taking the Lagrangian based on the multisine potential, we expand the interaction term around the minimum at $`\mathrm{\Phi }=\mathrm{𝟎}`$ :
$`{\displaystyle \underset{a=0}{\overset{n1}{}}}n_ae^{\beta _0𝐫_a.\mathrm{\Phi }}{\displaystyle \underset{a=0}{\overset{n1}{}}}n_a+\beta _0{\displaystyle \underset{a=0}{\overset{n1}{}}}n_a𝐫_a^i\mathrm{\Phi }^i`$ $`+`$ $`{\displaystyle \frac{\beta _0^2}{2!}}{\displaystyle \underset{a=0}{\overset{n1}{}}}n_a𝐫_a^i𝐫_a^j\mathrm{\Phi }^i\mathrm{\Phi }^j`$
$`+`$ $`{\displaystyle \frac{\beta _0^3}{3!}}{\displaystyle \underset{a=0}{\overset{n1}{}}}n_a𝐫_a^i𝐫_a^j𝐫_a^k\mathrm{\Phi }^i\mathrm{\Phi }^j\mathrm{\Phi }^k+\mathrm{}`$
Stabilization of the classical vacuum implies cancellation of the linear term. Using (4.42), for each $`n`$ even or $`n`$ odd (with (3.26)) case this is ensured iff :
$`{\displaystyle \frac{n_a}{n_{E(\frac{n+1}{2})1}}}={\displaystyle \frac{m_a}{m_{E(\frac{n+1}{2})1}}}\text{for}a\{0,\mathrm{},E({\displaystyle \frac{n+1}{2}})2\},`$ (4.45)
where $`E(x)`$ stands for the integer part of $`x`$. Consequently, we obtain the linear relation among the “parametrized” roots :
$`{\displaystyle \underset{a=0}{\overset{n1}{}}}m_a𝐫_a=0.`$ (4.46)
Such kind of relation is characteristic of ATFTs. One of the roots is generally identified with the negative of the highest root of the finite Lie algebra $`𝒢`$ (with rank $`n1`$) considered and the set of different numbers $`m_a`$ in (4.46) are proportional to the Kac labels. From eq. (4.45), it is interesting to see that ratios of these labels are completly defined by the kind of MC-algebra choosed. Furthermore, using eq. (4.45), the dual Kac labels associated with the dual ATFTs transform as in eq. (2.10). Under these considerations, to each ATFT (and its dual one) corresponds the pseudo-norm (4.39). As the ratios of $`m_a`$ take rational values (see the previous section), there exists one faithful representation $`\pi `$ , given by $`(n\times n)`$ dimensional diagonal matrices:
$`\pi \left[P_a\right]=Diag(0,\mathrm{},0,i,0,\mathrm{},0,i,0,\mathrm{},0)\text{for}a\{0,\mathrm{},n/21\},`$ (4.47)
where $`i`$ (resp. $`i`$) is in the $`a`$ (resp. $`n1a`$) position, e.g. :
$`(\pi \left[P_a\right])_{jj}=(\pi \left[P_{a}^{}{}_{}{}^{}\right])_{n1j,n1j}\text{where }\text{ denotes the ordinary hermitian}`$
$`\text{conjugate, for}(a,k)\{0,\mathrm{},n/21\}\text{and}j\{0,\mathrm{},n1\}.`$
For instance, consider the multisine-Gordon model in multicomplex dimension $`n=3`$ denoted by $`MSG_{(3|m)}(\alpha _0,\beta _0)`$, and its dual denoted by $`MSG_{(3|m^{})}(\alpha _0^{},\beta _0^{})`$. From eqs. (2.10) with substitutions (3.26), we have :
$`{\displaystyle \frac{m_0^{}}{m_1^{}}}={\displaystyle \frac{m_1}{m_0}}\left[{\displaystyle \frac{1}{2n_1^0}}\right].`$ (4.48)
As $`m=2m_0+m_1`$ and similarly $`m^{}=2m_0^{}+m_1^{}=m+\mathrm{\Delta }`$ with $`\mathrm{\Delta }=m_1m_0n_1^0`$, we obtain three solutions :
$``$ $`n_1^0=0`$ : for $`m_0=m_1=1(\mathrm{\Delta }=1)m=n=3`$ we obtain the model denoted $`𝒜^{(3|3)}(\beta _0,i\beta _0)`$ and for $`2m_0^{}=m_1^{}=2n=3m^{}=4`$ we obtain the model denoted $`𝒜^{(3|4)}(\frac{1}{\beta _0},\frac{i}{\beta _0})`$. These two models describe respectively the non-simply laced $`D_3^{(2)}`$ and $`C_2^{(1)}`$ ATFTs.
$``$ $`n_1^0=1`$ : for $`m_0=m_1=1(\mathrm{\Delta }=0)m=m^{}=n=3`$ we obtain the model denoted $`𝒜^{(3|3)}(\sqrt{\frac{3}{2}}\beta _0,\frac{i}{\sqrt{2}}\beta _0)`$. It is self-dual and describes the simply laced $`A_2^{(1)}`$ ATFT.
These three models are all generated by a fundamental multicomplex number $`e`$ satisfying $`e^3=1`$ (or its dual $`e^{}`$) which possesses the ($`m\times m`$) matrix representation $`\pi ^{}`$ :
$`𝒜^{(3|m)}(\alpha _0,\beta _0)\text{ with respect to }\pi ^{}[e]=`$ $`\text{Diag}[\underset{}{e^{i\frac{\pi }{3}}\mathrm{}e^{i\frac{\pi }{3}}},\underset{}{1\mathrm{}1},\underset{}{e^{i\frac{\pi }{3}}\mathrm{}e^{i\frac{\pi }{3}}}]`$
($`m_0`$ times) ($`m_1`$ times) ($`m_0`$ times)
$`\text{ dual with}`$
$`𝒜^{(3|m^{})}(\alpha _0^{},\beta _0^{})\text{ with respect to }\pi ^{}[e^{}]=`$ $`\text{Diag}[\underset{}{e^{i\frac{\pi }{3}}\mathrm{}e^{i\frac{\pi }{3}}},\underset{}{1\mathrm{}1},\underset{}{e^{i\frac{\pi }{3}}\mathrm{}e^{i\frac{\pi }{3}}}].`$
($`m_0^{}`$ times) ($`m_1^{}`$ times) ($`m_0^{}`$ times)
Consider now the multisine-Gordon model for $`n=4`$. The $`q`$-deformed structure of the non-local charge algebra is ensured for the choices $`n_1^0=0`$ and $`n_4=0`$. The resulting model $`𝒜^{(4|4)}`$ corresponds to the simply laced $`A_3^{(1)}`$ ATFT. Since its hidden symmetry is self-dual under transformations (2.9) and (2.10), the weak-strong coupling regimes (with respect to the parameter $`\beta _0`$) of the two dual actions are identical . However, while the fundamental multicomplex number representation associated with one model is :
$`𝒜^{(4|m)}(\alpha _0,\alpha _1;\beta _0)\text{ generated by :}`$
$`\pi ^{}[e]=\text{Diag}[\underset{}{e^{i\frac{\pi }{4}}\mathrm{}e^{i\frac{\pi }{4}}},\underset{}{e^{i\frac{3\pi }{4}}\mathrm{}e^{i\frac{3\pi }{4}}},\underset{}{e^{i\frac{3\pi }{4}}\mathrm{}e^{i\frac{3\pi }{4}}},\underset{}{e^{i\frac{\pi }{4}}\mathrm{}e^{i\frac{\pi }{4}}}],`$
($`m_0`$ times), ($`m_1`$ times), ($`m_1`$ times), ($`m_0`$ times)
its dual representation, associated with the dual model is obtained from $`\frac{m_0^{}}{m_1^{}}=\frac{m_1}{m_0}`$, which corresponds to :
$`𝒜^{(4|m^{})}(\alpha _0^{},\alpha _1^{};\beta _0^{})\text{ generated by :}`$
$`\pi ^{}[e^{}]=\text{Diag}[\underset{}{e^{i\frac{\pi }{4}}\mathrm{}e^{i\frac{\pi }{4}}},\underset{}{e^{i\frac{3\pi }{4}}\mathrm{}e^{i\frac{3\pi }{4}}},\underset{}{e^{i\frac{3\pi }{4}}\mathrm{}e^{i\frac{3\pi }{4}}},\underset{}{e^{i\frac{\pi }{4}}\mathrm{}e^{i\frac{\pi }{4}}}].`$
(
=m0m1
times
)
=m0m1
times
\left({\raisebox{-2.84544pt}{~{}\shortstack{ $m_{0}^{\vee}=m_{1}$ \\
times}}}\right) , (
=m1m0
times
)
=m1m0
times
\left({\raisebox{-2.84544pt}{~{}\shortstack{ $m_{1}^{\vee}=m_{0}$ \\
times}}}\right) , (
=m1m0
times
)
=m1m0
times
\left({\raisebox{-2.84544pt}{~{}\shortstack{ $m_{1}^{\vee}=m_{0}$ \\
times}}}\right)(
=m0m1
times
)
=m0m1
times
\left({\raisebox{-2.84544pt}{~{}\shortstack{ $m_{0}^{\vee}=m_{1}$ \\
times}}}\right)
Then, self-duality of $`A_3^{(1)}`$ translates into the following exchange :
$`ee^{}={\displaystyle \frac{1}{e}}`$ (4.49)
in the multicomplex space.
Action (4.40) in the multicomplex space (and its dual multicomplex) and action (4.43) in the usual complex space provide a unified representation of all lowest-rank affine Toda field theories. If the perturbation is relevant, no new operator will be generated under renormalization flows. Non-local MC-currents are conserved to all order in CPT and consequently quantum duality is satisfied. Tables 1.1-1.3 give the list of ATFTs described in this formalism whereas the exchange of the node under multicomplex conjugation in the corresponding Dynkin diagrams is depicted in fig. 1.1-1.3.
However, if the perturbation becomes marginal, MC-vertex operators associated to non-simple roots are generated under renormalization. For instance, let us consider the simply-laced cases ($`A_2^{(1)}`$, $`A_3^{(1)}`$ and $`D_4^{(1)})`$ in Table 1.1. The perturbation is marginal if $`\beta _0^2=1`$. The perturbing (self-dual) operator in MC-space can be written in the standard complex space as a current-current perturbation of a level-one WZNW model. We obtain respectively $`su(3)`$, $`su(4)`$ and $`so(8)`$ current-current perturbations <sup>6</sup><sup>6</sup>6The resulting action is of the Gross-Neveu type, obtained from the bosonization of the $`𝒢`$-invariant Gross-Neveu models.. This suggests a possible representation of WZNW or Gross-Neveu models in MC-space.
## 5 MC-algebra and perturbed conformal field theories
Instead of considering action (2.1) (or its dual counterpart to first order in CPT) as a perturbed $`(n1)`$ free field CFT, we can proceed differently. Let us consider the MSG potential $`_lx^{(l)}`$ for $`l\{0,\mathrm{},n1\}`$. If we truncate this potential by supressing one of the MC-operators, say $`x^{(0)}`$, the resulting form is no longer invariant under multicomplex conjugation. However, conformal invariance can be realized by adding some specific MC-charge at infinity coupled to the fundamental field of the theory. $`x^{(0)}`$ is then identified as the perturbation in MC-space. To show that, we define the holomorphic part of the MC-stress-energy tensor :
$`T(z)={\displaystyle \frac{1}{2}}(\mathrm{\Phi })^2+\sqrt{2}\beta _0𝐐.^2\mathrm{\Phi },`$ (5.50)
where $`𝐐𝕄_{(n|m)}`$ and similarly for the antiholomorphic part. For further convenience <sup>7</sup><sup>7</sup>7It is also possible to expand $`𝐐`$ over generators $`P_b`$ instead of $`P_b^2`$. In any case, the central charge takes real values., we write $`𝐐`$ as :
$`𝐐={\displaystyle \underset{b=0}{\overset{n/21}{}}}𝐐_b(P_b^2).`$ (5.51)
For $`n`$ even, using the representation (4.47) we define the $`n/2`$ projections $`\pi _a`$ over the usual complex space with :
$`\pi _a(x^{(0)})=\left(\pi [x^{(0)}]\right)_{aa}=e^{\beta _0𝐫_a.\mathrm{\Phi }}\text{for}a\{0,\mathrm{},n/21\}.`$ (5.52)
From eq. (5.52), we see that although the MSG Lagrangian representation does not depend on the choice of a particular projection, the Lagrangian representation in the usual complex space of the CFT part (i.e. equivalently the expression of the perturbation in the complex space) depends on this choice. In fact, due to the permutation symmetry of the “parametrized” roots $`𝐫_a`$ for $`a\{0,\mathrm{},n/22\}`$ (see Appendix B) the differences between all possible perturbations reduce to two distinct cases. In the first case, we use the projection <sup>8</sup><sup>8</sup>8In fact for $`n`$ even (and similarly for $`n`$ odd) all projections $`\pi _a`$ for $`a\{0,\mathrm{},n/22\}`$ are equivalent under multicomplex conjugation up to a permutation of the roots. For $`n`$ odd, using substitutions (3.26), it is obvious to see that the nodes associated respectively to $`n/21`$ and $`n1`$ “collapse” together as $`\alpha _{n/21}0`$. $`\pi _{n/21}`$ and the corresponding QFT is denoted $`𝒫`$. In the second case, we proceed similarly and use the projection $`\pi _0`$. The corresponding QFT is then denoted $`\overline{𝒫}`$. Each charge $`𝐐_b`$ is then associated to the CFT obtained from the projection $`\pi _b`$. Consequently, an appropriate choice of the charge Q, e.g. of the set $`\{𝐐_b\}`$ ensures simultaneously the conformal invariance of all CFTs. From these previous remarks, it reduces to study only $`𝒫`$ and $`\overline{𝒫}`$, i.e. to calculate $`𝐐_{n/21}`$ and $`𝐐_0`$. It translates into a condition on the conformal dimensions $`\mathrm{\Delta }_{n/21}`$ and $`\mathrm{\Delta }_0`$ of the vertex operators :
$`\mathrm{\Delta }_b\left(e^{\beta _0𝐫_a.\mathrm{\Phi }}\right)=1`$ $`\text{for}𝒫(\pi _{n/21}),\text{i.e.}\text{for all}a\{0,\mathrm{},{\displaystyle \frac{n}{2}}2,{\displaystyle \frac{n}{2}},\mathrm{},n2\},`$
or $`\text{for}\overline{𝒫}(\pi _0),\text{i.e.}\text{for all}a\{1,\mathrm{},n1\}.`$ (5.53)
Using eq. (5.50) and (5.51), the holomorphic conformal dimension of each vertex operator is :
$`\mathrm{\Delta }_b\left(e^{\beta _0𝐫_a.\mathrm{\Phi }}\right)={\displaystyle \frac{\beta _0^2}{2}}𝐫_a^2+\sqrt{2}\beta _0𝐐_b.𝐫_a.`$ (5.54)
For further convenience and by analogy with ATFTs approach, let us introduce :
$`𝐐_b={\displaystyle \frac{1}{\sqrt{2}}}\left[\beta _0\rho _b+\beta _0^{}\rho _b^{}\right]\text{where}\rho _b={\displaystyle \underset{\{c\}}{}}\omega _{c;b}\text{and}\rho _b^{}={\displaystyle \underset{\{c\}}{}}\omega _{c;b}^{}`$ (5.55)
with $`cn/21`$ for $`𝒫`$ and $`c0`$ for $`\overline{𝒫}`$. Eqs. (5) can be satisfied if $`\omega _{c;b},\omega _{c;b}^{}`$ are choosed to obey :
$`\omega _{c;b}^{}.𝐫_a=\delta _{ac}{\displaystyle \frac{1}{\beta _0\beta _0^{}}}.`$ (5.56)
Similarly, the same approach can be applied to the $`n`$ odd case and the corresponding results are obtained using the substitutions (3.26). In each case ($`𝒫`$ or $`\overline{𝒫}`$) and any value of $`n`$, it is straightforward to compute the MC-central charge of the conformally invariant part:
$`c=n1+24|𝐐|^2.`$ (5.57)
We notice that expression (5.55) is self-dual under the duality transformation (2.9). From the above analysis, we have computed the “parametrized” central charges $`c_b=\pi _b(c)`$ for $`𝒫`$ and $`\overline{𝒫}`$, expressed in terms of $`n`$, $`(m_a,m_a^{})`$, $`(m,m^{})`$, $`\{𝐫_a\}`$ and $`\{\beta _a\}`$. The “parametrized” co-weights are given in Appendix B and their associated weights are obtained using transformations (2.9). Using eq. (5.57) with eqs. (5.51) and (5.55) we obtain :
* For $`𝒫`$ :
$`c_{n/21}=n1+12\left[\mathrm{\Gamma }_{n;n/21}^{(1)}+\mathrm{\Gamma }_{n;n/21}^{(2)}\right]\text{for }n\text{ even},`$ (5.58)
$`c_{(n1)/2}=n1+12\left[\mathrm{\Gamma }_{n;(n1)/2}^{(1)}\right]\text{for }n\text{ odd}.`$
* For $`\overline{𝒫}`$ and $`n`$ even :
$`c_0=n1+12\left[\mathrm{\Gamma }_{n;0}^{(1)}+\mathrm{\Gamma }_{n;0}^{(2)}+\mathrm{\Gamma }_{n;0}^{(3)}\right]`$ (5.59)
and where we define for $`n`$ even :
$`\mathrm{\Gamma }_{n;b}^{(1)}={\displaystyle \underset{\text{a=0a≠b}}{\overset{n/22}{}}}\left[\left({\displaystyle \frac{𝐫_a^2}{2}}\beta _0+{\displaystyle \frac{1}{\beta _0}}\right)^2\left({\displaystyle \frac{\beta _0}{\beta _a}}\right)^2\right],`$
$`\mathrm{\Gamma }_{n;b}^{(2)}=\left({\displaystyle \frac{_{a=0}^{n/21}m_a𝐫_a^2}{2m_b}}\beta _0+{\displaystyle \frac{m}{2m_b}}{\displaystyle \frac{1}{\beta _0}}\right)^2\left({\displaystyle \frac{\beta _0}{\alpha _b}}\right)^2,`$ (5.60)
$`\mathrm{\Gamma }_{n;0}^{(3)}=\left({\displaystyle \frac{_{a=1}^{n/21}m_a𝐫_a^2}{2m_0}}\beta _0+{\displaystyle \frac{m2m_0}{2m_0}}{\displaystyle \frac{1}{\beta _0}}\right)^2.`$
For $`\overline{𝒫}`$ and $`n`$ odd, we use substitution (3.26).
The conformal dimension of each vertex operator, exept the perturbation, is one. Then, their integrals appear naturally as “screening charges” whereas the conformal dimension of the perturbation for any projection $`\pi _b`$ is :
$`\mathrm{\Delta }_b\left(\pi _b(x^{(0)})\right)=\mathrm{\Delta }_b\left(e^{\beta _0𝐫_b.\mathrm{\Phi }}\right)=1\left[{\displaystyle \frac{m^{}}{m_b^{}}}\left(\beta _0^2{\displaystyle \frac{𝐫_b^2}{2}}\right)+{\displaystyle \frac{m}{m_b}}\right].`$ (5.61)
For $`(\{\alpha _a\}`$, $`\{\beta _a\}i)`$, the perturbation is relevant ($`\mathrm{\Delta }_b<1`$) or marginal ($`\mathrm{\Delta }_b=1`$) iff:
$`(\beta _0^{})^2{\displaystyle \frac{𝐫_b^2}{2}}{\displaystyle \frac{\frac{m}{m_b}}{\frac{m^{}}{m_b^{}}}}.`$ (5.62)
The renormalizability property of the model is then encoded in the MC-algebra.
As we saw previously, depending on the dimension, the specific structure of the MC-algebra (e.g. the values of $`m_a`$) and ratios of the parameters, the resulting model possesses a $`q`$-deformed symmetry. For each projection which leads to $`𝒫`$ or $`\overline{𝒫}`$, we have computed for the ratios reported in Table 1.1 the central charge associated with the truncated MSG model, i.e. without the perturbing term $`x^{(0)}`$. For instance, for projection $`𝒫`$ for case a/ and b/ the truncated model is identified to an $`su(2)su(2)`$ theory (two decoupled Liouville theories), whereas case e/ corresponds to an $`so(4)so(4)`$ theory (four decoupled Liouville theories). For projection $`\overline{𝒫}`$, Toda field theories are obtained. For any projection, case c/ and d/ are respectively identified with $`A_2`$ and $`A_3`$ Toda field theories. Using the results reported in Table 1.1, it can be checked that agreement is obtained with the results found in . However, for imaginary values of the parameters $`\beta _a`$, a truncation of the Hilbert space is necessary in order to obtain a unitary theory. As ratios are fixed, the values of $`\beta _0`$ are fixed in each case, corresponding to a quantum group restriction of the model. This last restriction corresponds to points in parameter space.
For generic values of the parameters the situation is much more complicated. Although the central charge associated with the truncated MSG model can be computed, it is not clear if unitary <sup>9</sup><sup>9</sup>9However, many statistical systems do not satisfy the unitarity condition (nonself intersecting polymer chains, magnetics with stochastic interactions, etc…). representations of the Virasoro algebra always exist . However, for $`𝐫_a^2=2`$ and $`a\{0,\mathrm{},\frac{n}{2}1\}`$ (which is the standard convention for simply laced affine Lie algebra), we notice that any “truncation” ($`𝒫`$ or $`\overline{𝒫}`$) of the basis of MC-vertex operators in action (2.1) leads to a self-dual CFT (ratios of parameters are fixed) and condition (5.62) becomes $`(\beta _0^{})^21`$.
## 6 Concluding remarks
Consider the MSG model with action $`𝒜^{(n|m)}(\eta )`$ generated by the MC-algebra $`𝕄_{(n|m)}`$. The parameter space can be described as follow. In the deep ultraviolet, it exists a discret set of one-dimensional submanifolds $`𝒮_0=(\{\alpha _a^{(0)}\},\{\beta _a^{(0)}\})`$ (see Tables 1.1-1.3) which corresponds to a scale invariant theory, i.e a CFT. This CFT with MC-central charge (5.57) possess a Lagrangian representation in MC-space in terms of a truncated basis of MC-vertex operators, $`_{l=1}^{n1}x^{(l)}`$ for instance. In usual complex space, it admits two unequivalent Lagrangian representations $`𝒫`$ and $`\overline{𝒫}`$. A parametrized central charge $`c_b=c_b(\{\alpha _a^{(0)}\},\{\beta _a^{(0)}\})`$ for $`b=0`$ or $`E((n+1)/2)1`$ was obtained for each representations, describing the fixed points of different known CFTs like decoupled Toda or Toda field theories.
We are now interested in the neighbourhood region of these fixed points where we loose the scale invariance of the model. As the CFT in MC-space is not invariant under multicomplex conjugation, a natural perturbation is then provided by imposing the MC-conjugation symmetry to the resulting model. Indeed this perturbation, say $`x^{(0)}`$, does not correspond to an arbitrary deformation of the CFT. The perturbed model always corresponds to a one-parameter family of integrable massive QFTs. For $`𝒮_0^+=(\{\alpha _a^{(0)}\},\{\beta _a^{(0)}\}i)`$, the model is identified to lowest-rank imaginary coupling ATFTs which possess soliton solutions. However, unitarity is only assumed at specific points on $`𝒮_0^+`$. The theory gets truncated and the hidden symmetry acts as a parametrized quantum group with the deformation parameters $`q_b(\{\alpha _a^{(0)}\},\{\beta _a^{(0)}\})`$ being a root of unity. For $`𝒮_0^{}=(\{\alpha _a^{(0)}\}i,\{\beta _a^{(0)}\})`$, lowest-rank real coupling ATFTs are obtained.
Exept in the Kosterlitz-Thouless region of $`𝒮_0`$, no new terms are generated under the renormalization group flow. Non-local MC-conserved currents $`\{J_{\eta ^{(k)}},\overline{J}_{\eta ^{(k)}}\}`$ exist which are conserved to all orders in CPT. There, the MSG model (2.11) with (2.12) admits a dual Lagrangian representation with action $`𝒜^{(n|m^{})}(\eta ^{})`$, generated by the dual multicomplex algebra $`𝕄_{(n|m^{})}`$. The weak coupling regime of one MSG and the strong coupling regime of its dual are simply related through the parameter exchange $`\beta _01/\beta _0`$.
The neighbourhood region of $`𝒮_0`$ $`(\lambda <<1)`$ is much more complicated but a few remarks can be done from our previous analysis. First, the non-local currents (and dual ones) are still conserved, at least to first order in CPT, for a larger parameter space $`𝒮_1`$ described by eqs. (2.16), (2.17), (2.18). It is easy to see that the symmetry group of $`𝒮_1^+`$ (and similarly for $`𝒮_1^{}`$ with the opposite signature) is the pseudo-rotational group $`SO(\frac{n}{2}R,\frac{n}{2}1)`$ for $`n`$ even and $`SO(\frac{n1}{2}R,\frac{n1}{2})`$ for $`n`$ odd. However, the conserved charges (2.19) do not satisfy the previous $`q`$-deformed symmetry since braiding relations between MC-operators strongly depend on the ratios of parameters. A less restrictive algebraic structure, like more general Hopf algebras, may exist. It is known that integrable models can be generated by some more general algebras than the quantum Lie algebras : the integrability condition is dictated by the Yang-Baxter equation itself . Secondly, the non-local conserved currents does not necessary form a closed algebra. Then, renormalization group flow may generated counterterms which can spoiled the current conservation to higher order in CPT.
To conclude, we would like to mention that all the integrable (dual-)QFTs in and here are generated by the MC-algebra $`𝕄_{(n|m)}`$ for specific values of $`(n,m)`$. The difference between them simply reduces to the parameter space. Particulary, the non-local conserved charges (2.19), if they satisfy a $`q`$-deformed algebraic structure , can be simply expressed in terms of a Lie algebra based on the field $`𝕄_{(n|m)}`$ (multicomplex are of characteristic 0 as for the usual complex numbers), where parameters $`(\alpha _a,\beta _a)`$ play the role of deformations. Furthermore, it is believed that more general MC-algebras exist . Many other QFTs should then be described within this formalism.
### Aknowledgements
We are very grateful to C. Ahn, J.L. Kneur, A. Neveu, M. Rausch de Traubenberg, F. Smirnov, G. Takács and J. Thierry-Mieg for useful discussions. P.B. thanks for the hospitality of KIAS where part of this work was done. Work supported in part by the EU under contract ERBFMRX CT960012.
Appendix A
| | $`n`$ | $`R`$ | a | $`n_1^a`$ | $`n_2^a`$ | $`n_3^a`$ | $`A_a/B_0`$ | $`A_{\frac{n}{2}1}/B_0`$ | $`B_a/B_0`$ | $`\delta _a`$ | Algebra |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| a/ | 4 | 1 | 0 | * | 2 | 1 | 0 | -1 | 1 | 1 | ($`n_4=0`$) $`C_2^{(1)}`$ |
| b/ | 4 | 1 | 0 | * | 1 | 2 | 0 | -1/4 | 1 | 1/2 | ($`n_4=0`$) $`D_3^{(2)}`$ |
| c/ | 4 | 1 | 0 | * | 1 | 1 | 0 | -3/4 | 1 | 1/2 | ($`n_4=1`$) $`A_2^{(1)}`$ |
| d/ | 4 | 2 | 0 | * | 4 | 1 | 0 | 0 | 1 | 2 | $`A_2^{(2)}`$ |
| e/ | 4 | 2 | 0 | * | 1 | 4 | 0 | 0 | 1 | 1/2 | $`A_2^{(2)}`$ |
| f/ | 4 | 2 | 0 | * | 2 | 2 | 0 | 0 | 1 | 1 | $`A_1^{(1)}`$ |
| | | | 0 | * | 1 | 2 | 0 | | 1 | 1/2 | |
| g/ | 6 | 2 | | | | | | 0 | | | $`A_5^{(2)}`$ |
| | | | 1 | 0 | 1 | 1 | -1/4 | | 1/4 | 1 | |
| | | | 0 | * | 2 | 1 | 0 | | 1 | 1 | |
| h/ | 6 | 2 | | | | | | 0 | | | $`B_3^{(1)}`$ |
| | | | 1 | 0 | 1 | 1 | -1 | | 1 | 1 | |
| | | | 0 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| i/ | 6 | 3 | | | | | | 0 | | | $`D_4^{(3)}`$ |
| | | | 1 | * | 1 | 3 | 0 | | 3 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| j/ | 6 | 3 | | | | | | 0 | | | $`G_2^{(1)}`$ |
| | | | 1 | * | 3 | 1 | 0 | | 1/3 | 3/2 | |
| | | | 0 | * | 2 | 1 | 0 | | 1 | 1 | |
| k/ | 6 | 3 | | | | | | 0 | | | $`D_3^{(2)}`$ |
| | | | 1 | * | 2 | 1 | 0 | | 1 | 1 | |
| | | | 0 | * | 1 | 2 | 0 | | 1 | 1/2 | |
| l/ | 6 | 3 | | | | | | 0 | | | $`A_4^{(2)}`$ |
| | | | 1 | * | 2 | 1 | 0 | | 1/4 | 1 | |
| | | | 0 | * | 1 | 2 | 0 | | 1 | 1/2 | |
| m/ | 6 | 3 | | | | | | 0 | | | $`C_2^{(1)}`$ |
| | | | 1 | * | 1 | 2 | 0 | | 1 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| n/ | 6 | 2 | | | | | | -1/2 | | | ($`n_4=0`$) $`A_3^{(1)}`$ |
| | | | 1 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| o/ | 8 | 3 | 1 | * | 1 | 1 | 0 | 0 | 1 | 1/2 | $`D_4^{(1)}`$ |
| | | | 2 | 0 | 1 | 1 | -1/2 | | 1/2 | 1 | |
| | | | 0 | * | 1 | 2 | 0 | | 1 | 1/2 | |
| p/ | 8 | 4 | 1 | * | 1 | 1 | 0 | 0 | 1/2 | 1/2 | $`A_5^{(2)}`$ |
| | | | 2 | * | 1 | 1 | 0 | | 1/2 | 1/2 | |
| | | | 0 | * | 2 | 1 | 0 | | 1 | 1 | |
| q/ | 8 | 4 | 1 | * | 1 | 1 | 0 | 0 | 2 | 1/2 | $`B_3^{(1)}`$ |
| | | | 2 | * | 1 | 1 | 0 | | 2 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| | | | 1 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| r/ | 10 | 5 | | | | | | 0 | | | $`D_4^{(1)}`$ |
| | | | 2 | * | 1 | 1 | 0 | | 1 | 1/2 | |
| | | | 3 | * | 1 | 1 | 0 | | 1 | 1/2 | |
Table 1.2 \- Solutions for the even case with restrictions, and the corresponding algebras
of rank $`r=n1R`$. $`R`$ : number of restrictions ($`R=s+1`$ if $`A_{n/21}0`$, else $`R=s+2`$).
$``$ stands for one of the $`s+1`$ restrictions as defined in (3.33)
Figure 1.2 \- Dynkin diagrams corresponding to affine Lie algebras Table 1.2 .
| | $`n`$ | $`R`$ | a | $`n_1^a`$ | $`n_2^a`$ | $`n_3^a`$ | $`A_a/B_0`$ | $`B_a/B_0`$ | $`\delta _a`$ | Algebra |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| a/ | 3 | 1 | 0 | * | 4 | 1 | 0 | 1 | 1 | $`A_2^{(2)}`$ |
| b/ | 3 | 1 | 0 | * | 1 | 4 | 0 | 1 | 1/4 | $`A_2^{(2)}`$ |
| c/ | 3 | 1 | 0 | * | 2 | 2 | 0 | 1 | 1/2 | $`A_1^{(1)}`$ |
| | | | 0 | * | 1 | 2 | 0 | 1 | 1/4 | |
| d/ | 5 | 1 | | | | | | | | $`A_5^{(2)}`$ |
| | | | 1 | 0 | 1 | 1 | -1/4 | 1/4 | 1/2 | |
| | | | 0 | * | 2 | 1 | 0 | 1 | 1/2 | |
| e/ | 5 | 1 | | | | | | | | $`B_3^{(1)}`$ |
| | | | 1 | 0 | 1 | 1 | -1 | 1 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | 1 | 1/4 | |
| f/ | 5 | 2 | | | | | | | | $`D_4^{(3)}`$ |
| | | | 1 | * | 1 | 3 | 0 | 3 | 1/4 | |
| | | | 0 | * | 1 | 1 | 0 | 1 | 1/4 | |
| g/ | 5 | 2 | | | | | | | | $`G_2^{(1)}`$ |
| | | | 1 | * | 3 | 1 | 0 | 1/3 | 3/4 | |
| | | | 0 | * | 1 | 2 | 0 | 1 | 1/4 | |
| h/ | 5 | 2 | | | | | | | | $`C_2^{(1)}`$ |
| | | | 1 | * | 1 | 2 | 0 | 1 | 1/4 | |
| | | | 0 | * | 1 | 2 | 0 | 1 | 1/4 | |
| i/ | 5 | 2 | | | | | | | | $`A_4^{(2)}`$ |
| | | | 1 | * | 2 | 1 | 0 | 1/4 | 1/2 | |
| | | | 0 | * | 2 | 1 | 0 | 1 | 1/2 | |
| j/ | 5 | 2 | | | | | | | | $`D_3^{(2)}`$ |
| | | | 1 | * | 2 | 1 | 0 | 1 | 1/2 | |
| | | | 0 | * | 1 | 1 | 0 | 1 | 1/4 | |
| k/ | 7 | 2 | 1 | * | 1 | 1 | 0 | 1 | 1/4 | $`D_4^{(1)}`$ |
| | | | 2 | 0 | 1 | 1 | -1/2 | 1/2 | 1/2 | |
| | | | 0 | * | 1 | 2 | 0 | 1 | 1/4 | |
| l/ | 7 | 3 | 1 | * | 1 | 1 | 0 | 1/2 | 1/4 | $`A_5^{(2)}`$ |
| | | | 2 | * | 1 | 1 | 0 | 1/2 | 1/4 | |
| | | | 0 | * | 2 | 1 | 0 | 1 | 1/2 | |
| m/ | 7 | 3 | 1 | * | 1 | 1 | 0 | 2 | 1/4 | $`B_3^{(1)}`$ |
| | | | 2 | * | 1 | 1 | 0 | 2 | 1/4 | |
| | | | 0 | * | 1 | 1 | 0 | 1 | 1/4 | |
| | | | 1 | * | 1 | 1 | 0 | 1 | 1/4 | |
| n/ | 9 | 4 | | | | | | | | $`D_4^{(1)}`$ |
| | | | 2 | * | 1 | 1 | 0 | 1 | 1/4 | |
| | | | 3 | * | 1 | 1 | 0 | 1 | 1/4 | |
Table 1.3 \- Solutions for the odd case with restrictions, and the corresponding algebras.
of rank $`r=n1R`$. $`R`$ : number of restrictions : $`R=s+1`$.
$``$ stands for one of the $`s+1`$ restrictions as defined in (3.36).
Figure 1.3 \- Dynkin diagrams corresponding to affine Lie algebras Table 1.3 .
Appendix B
$`n`$ even.
The set of “parametrized” roots for $`n`$ even reads, with $`a\{0,\mathrm{},\frac{n}{2}2\}`$ :
$`𝐫_a=[0,\mathrm{},0,i{\displaystyle \frac{\alpha _a}{\beta _0}},{\displaystyle \frac{\beta _a}{\beta _0}},0,\mathrm{},0],𝐫_{a+\frac{n}{2}}=[0,\mathrm{},0,i{\displaystyle \frac{\alpha _a}{\beta _0}},{\displaystyle \frac{\beta _a}{\beta _0}},0,\mathrm{},0],`$
$`𝐫_{\frac{n}{2}1}=[0,{\displaystyle \frac{m_0}{m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{\beta _am_a}{\beta _0m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{\beta _{\frac{n}{2}2}m_{\frac{n}{2}2}}{\beta _0m_{\frac{n}{2}1}}},i{\displaystyle \frac{\alpha _{\frac{n}{2}1}}{\beta _0}}],`$ (6.63)
$`𝐫_{n1}=[0,{\displaystyle \frac{m_0}{m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{\beta _am_a}{\beta _0m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{\beta _{\frac{n}{2}2}m_{\frac{n}{2}2}}{\beta _0m_{\frac{n}{2}1}}},i{\displaystyle \frac{\alpha _{\frac{n}{2}1}}{\beta _0}}].`$
For $`𝒫`$, the set of dual “parametrized” co-weights defined by (5.56) reads, with $`a\{0,\mathrm{},\frac{n}{2}2\}`$ :
$`\beta _0^{}\omega _{a;\frac{n}{2}1}^{}`$ $`=`$ $`[0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0,\mathrm{},0,{\displaystyle \frac{m_a}{2im_{\frac{n}{2}1}\alpha _{\frac{n}{2}1}}}],`$
$`\beta _0^{}\omega _{a+\frac{n}{2};\frac{n}{2}1}^{}`$ $`=`$ $`[0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0,\mathrm{},0,{\displaystyle \frac{m_a}{2im_{\frac{n}{2}1}\alpha _{\frac{n}{2}1}}}],`$ (6.64)
$`\beta _0^{}\omega _{n1;\frac{n}{2}1}^{}`$ $`=`$ $`[0,\mathrm{}\mathrm{},0,{\displaystyle \frac{1}{i\alpha _{\frac{n}{2}1}}}].`$
For $`\overline{𝒫}`$, the set of dual “parametrized” co-weights reads, with $`a\{1,\mathrm{},\frac{n}{2}2\}`$ is :
$`\beta _0^{}\omega _{a;0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_a}{2im_0\alpha _0}},{\displaystyle \frac{m_a}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0\mathrm{},0],`$
$`\beta _0^{}\omega _{a+\frac{n}{2};0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_a}{2im_0\alpha _0}},{\displaystyle \frac{m_a}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0\mathrm{},0],`$
$`\beta _0^{}\omega _{\frac{n}{2}1;0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_{\frac{n}{2}1}}{2im_0\alpha _0}},{\displaystyle \frac{m_{\frac{n}{2}1}}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _{\frac{n}{2}1}}}],`$ (6.65)
$`\beta _0^{}\omega _{n1;0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_{\frac{n}{2}1}}{2im_0\alpha _0}},{\displaystyle \frac{m_{\frac{n}{2}1}}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _{\frac{n}{2}1}}}],`$
$`\beta _0^{}\omega _{\frac{n}{2};0}^{}`$ $`=`$ $`[{\displaystyle \frac{1}{i\alpha _0}},0\mathrm{},0].`$
Similarly, $`\omega _a`$ are obtained using substitutions $`\alpha _a\alpha _a^{}`$, $`\beta _a\beta _a^{}`$ and $`m_am_a^{}`$.
$`n`$ odd.
For $`n`$ odd the set of “parametrized” roots is given by ( with $`a\{0,\mathrm{},\frac{n1}{2}1\}`$) :
$`𝐫_a=[0,\mathrm{},0,i{\displaystyle \frac{\alpha _a}{\beta _0}},{\displaystyle \frac{\beta _a}{\beta _0}},0,\mathrm{},0],𝐫_{a+\frac{n+1}{2}}=[0,\mathrm{},0,i{\displaystyle \frac{\alpha _a}{\beta _0}},{\displaystyle \frac{\beta _a}{\beta _0}},0,\mathrm{},0],`$
$`𝐫_{\frac{n1}{2}}=[0,{\displaystyle \frac{2m_0}{m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{2\beta _am_a}{\beta _0m_{\frac{n}{2}1}}},\mathrm{},0,{\displaystyle \frac{2\beta _{\frac{n}{2}2}m_{\frac{n}{2}2}}{\beta _0m_{\frac{n}{2}1}}}].`$ (6.66)
For $`𝒫`$, the set of dual “parametrized” co-weights is then, with $`a\{0,\mathrm{},\frac{n1}{2}1\}`$ :
$`\beta _0^{}\omega _{a;\frac{n1}{2}}^{}`$ $`=`$ $`[0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0,\mathrm{},0],`$
$`\beta _0^{}\omega _{a+\frac{n+1}{2};\frac{n1}{2}}^{}`$ $`=`$ $`[0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0,\mathrm{},0].`$ (6.67)
For $`\overline{𝒫}`$, the set of dual “parametrized” co-weights is, with $`a\{1,\mathrm{},\frac{n1}{2}1\}`$ :
$`\beta _0^{}\omega _{a;0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_a}{2im_0\alpha _0}},{\displaystyle \frac{m_a}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0\mathrm{},0],`$
$`\beta _0^{}\omega _{a+\frac{n+1}{2};0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_a}{2im_0\alpha _0}},{\displaystyle \frac{m_a}{2m_0\beta _0}},0,\mathrm{},0,{\displaystyle \frac{1}{2i\alpha _a}},{\displaystyle \frac{1}{2\beta _a}},0\mathrm{},0],`$ (6.68)
$`\beta _0^{}\omega _{\frac{n1}{2};0}^{}`$ $`=`$ $`[{\displaystyle \frac{m_{\frac{n1}{2}}}{2im_0\alpha _0}},{\displaystyle \frac{m_{\frac{n1}{2}}}{2m_0\beta _0}},0,\mathrm{},0],`$
$`\beta _0^{}\omega _{\frac{n1}{2}+1;0}^{}`$ $`=`$ $`[{\displaystyle \frac{1}{i\alpha _0}},0\mathrm{},0].`$ |
warning/0001/nlin0001071.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In accordance with the finite-band theory, the integrable Korteweg-de Vries (KdV) equation can be considered as the compatibility condition of the two auxiliary linear differential matrix equations $`_{x_\pm }𝚽=𝐔_\pm `$ ($`x_+x,x_{}t`$, $`_xd/dx`$) with matrix operators $`𝐔_\pm `$, coefficients of which, as is well known (see ), are expressed through solutions of the KdV equation. The finite-gap solutions of the KdV equation are solutions of the spectral problem for these equations with a finite-gap spectrum of eigenvalues. These are expressed through multi-dimensional Riemann theta functions with implicit parameters, the evaluation of which is the special algebraic geometrical problem (see ). The class of elliptic finite-gap solutions leads to the problem of the reduction of $`n`$-dimensional theta-functions to the one-dimensional Jacobi theta functions (see ).
In the framework of the spectral problem the finite-gap solutions satisfy the matrix finite-gap equation of the form $`_\pm 𝚿=𝐔_\pm 𝚿𝚿𝐔_\pm `$, where $`𝚿`$ is the matrix components which are polynomial in the eigenvalue of the auxiliary equations (such as in the case of the “sine-Gordon” equation ). Usually, solving this equation is realized with help of the well known Abel transformation with a subsequent solving of the inverse Jacobi problem (see ).
However, in the case under consideration the finite-gap solutions of the KdV equation are elliptic functions represented as rational functions of the elliptic Weierstrass function ($`\mathrm{}`$-function) . Solving the finite-gap equations in terms of $`\mathrm{}`$-function can be reduced to solving simple algebraic equations. This yields the straightforward manner for calculating the elliptic finite-gap solutions in an explicit form.
We shown that in the initial time ($`t=0`$) the finite-gap equations can be reduced to algebraic equations with respect to parameters of the above mentioned rational functions. The computation of these parameters gives all possible initial elliptic finite-gap solutions in the form of linear combinations of $`\mathrm{}`$-functions with shifted arguments.
In accordance with the KdV equation the time dependent elliptic solutions are built as linear combinations of $`\mathrm{}`$-functions with time dependent argument shifts ($`\phi _i`$) (its poles) under condition of its compatibility with corresponding initial solutions. Their time evolution is determined by the poles which satisfy the system of linked dynamic equations which follow from the above mentioned auxiliary equations.
Using the finite-gap equations we shown that the dynamic system is transformed in the system of independent differential equations of the first order with separated variables of the form $`_t\phi _i=X_i(\phi _i)`$. Here $`X_i`$-functions represent themselves roots of some polynomial equations which follow from the finite-gap equations the order of which equals the number of poles in the elliptic solutions.
This paper is organized as follows. In Section 2 the approach to the straightforward calculation of the elliptic finite-gap solutions of the KdV equation, based on the auxiliary system of finite-gap and dynamic equations, is formulated. In Section 3 this approach is applied to the calculation of the initial elliptic 1-, 2- and 3-gap solutions of the KdV equation. In Section 4 the time evolution of these elliptic solutions is investigated. We shown that the linked system of auxiliary dynamic equations, for poles of the corresponding time dependent elliptic solutions, can be integrated with the help of the finite-gap equations.
## 2 Finite-band equations and general elliptic solutions
In the class of elliptic functions (which we shall denote as $`U(x,t)`$) with the $`\mathrm{}`$-functional representation under consideration, the finite-gap solutions of the KdV equation can be considered as solutions of the finite-band equations. The latter gives the compatibility condition of the finite-gap and the general solutions of the auxiliary linear differential equations. These finite-band equations represent a system of equalities, obtained by equating coefficients of the power series in the eigenvalue $`E`$ of the finite-gap and general solutions (see ).
The finite band equations (which are also known as “trace-formulae”) can be written in the form
$$A_{n+1}=\frac{(1)^n}{2^{2n+1}}\chi _{2n+1}(x,t)(n=0,1,\mathrm{}),a_0=1$$
(2.1)
$$A_n=\frac{1}{n!}_z^n\left(\frac{\sqrt{_{n=0}^{2g+1}a_nz^n}}{_{n=0}^gb_n(x,t)z^n}\right)|_{(z=0)},b_0=1$$
(2.2)
($`g`$ is the number of a gaps in the spectrum of the eigenvalues $`E`$), where the $`\chi _n`$-functions are determined by the recursion relation
$$\chi _{n+1}=_x\chi _n+\underset{k=1}{\overset{n1}{}}\chi _k\chi _{nk},\chi _1=U(x).$$
(2.3)
Here $`A_n`$ and $`\chi _n`$ are coefficient functions of the power-series expansion in $`E`$ of the general and the finite-gap solutions of the auxiliary equations, respectively. In view of (2.3) the $`\chi `$-functions have the form of polynomials in the solutions $`U(x,t)`$ and its derivatives. As is well known, (2.1) is algebraicly solvable with respect to the coefficient functions $`b_n(x,t)`$. Therefore, the finite-gap equations (2.1) are reduced to the system containing expressions for the coefficient functions $`b_n(x,t),n=(1,\mathrm{},g)`$, which at $`ng`$ is a closed system of equations for the elliptic finite-gap $`U`$-solutions.
The elliptic finite-gap solutions as double periodic functions of the complex variable $`z`$ which admit the rational functions representation namely the elliptic Weierstrass function $`\mathrm{}`$ . In view of the equation (2.1), $`U`$ is determined by the formula
$$U(z,t)=\alpha \mathrm{}(z)+\underset{i}{}\underset{n_i=1}{\overset{2}{}}\left\{\frac{\alpha _{n_i}(t)}{(\mathrm{}(z)h_i(t))^{n_i}}\right\}+\stackrel{~}{R}(z,t)$$
(2.4)
with poles of the secondary order in $`\mathrm{}\mathrm{}(z|\omega ,\omega ^{})`$ ($`\omega `$ and $`\omega ^{}`$ are real and imaginary half-periods of the $`\mathrm{}`$-function), where $`\stackrel{~}{R}(z,t)`$ means an odd function of $`z`$.
At the initial time ($`t=0`$) the expression (2.4) describes all possible so-called initial elliptic finite-gap solutions, which in the case of even functions ($`\stackrel{~}{R}=0`$) under consideration, takes the form
$$U(z)=\alpha \mathrm{}(z)+\underset{i}{}\underset{n_i=1}{\overset{2}{}}\frac{\alpha _{n_i}}{(\mathrm{}(z)h_i)^{n_i}},$$
(2.5)
that reduces the finite-gap equations (2.1) to simple algebraic equations with respect to the $`\alpha `$\- and $`h_i`$-parameters.
The time evolution of the elliptic finite-gap solutions (2.4) are determined by the time dependence of their poles, which are described with the help of the known dynamic auxiliary equation
$$_t\mathrm{\Psi }(x,t,E)=\left(4+3(U_x+_xU)\right)\mathrm{\Psi }(x,t,E),$$
(2.6)
$`\mathrm{\Psi }=\mathrm{\Psi }(\{a_n)\},\{b_n\})`$, $`b_n=b_n(U(z,t)`$, $`U^{(n)}(z,t))`$ can be reduced to the system of dynamic equations with respect to the poles of the $`U`$-functions. This is a system of coupled differential equations of the first order which, as will be shown below, can be reduced to independent equations with separated variables.
## 3 A calculation of the initial elliptic solutions
The calculation of the initial finite-gap elliptic solutions of the KdV equation is based on the use of the finite-band equations (2.1) in the representation of the rational functions (2.5). On equating the corresponding coefficients of the Laurent expansion in $`\mathrm{}`$ of the left- and right-hand sides (2.1), yields simple algebraic equations which determine the parameters of the expression (2.5) for the initial elliptic finite-gap solutions of the KdV equation.
### 3.1 One-gap elliptic solutions
In the one-gap case the parameters of the elliptic solution $`U(z)`$ are determined by the system of three finite-gap equations of the form (2.1) at $`n=\overline{0,2}`$. A substitution in these equations, of the explicit expressions (2.2) for $`A_n`$ and polynomials in $`U`$ expressions, for the $`\chi _n`$-functions which follow from (2.3), yields the system
$$\begin{array}{cc}& \frac{1}{2}a_1b_1=\frac{1}{2}U;\hfill \\ & \frac{1}{2}\{(a_2\frac{1}{4}a_1^2)+2(b_1^2b_2)a_1b_1\}=\frac{1}{2^3}\{U^2U^{(2)}\};\hfill \\ & \frac{1}{3!}\{(\frac{3}{8}a_1^3\frac{3}{2}a_1a_2+\frac{3}{2}a_3)+3(\frac{1}{4}a_1^2a_2)b_1+3a_1(b_1^2b_2)\hfill \\ & +(12b_1b_24!b_36b_1^3)\}=\frac{1}{2^5}\{U^{(4)}5U^{(1)^2}+6UU^{(2)}2U^3\}\hfill \end{array}$$
(3.1)
in which $`b_n|_{n2}=0`$ (in view of the relation $`b_n|_{ng+1}=0`$), where $`g`$ is the number of gaps in the spectrum of the auxiliary linear differential equation. Excluding $`b_n`$ from the system (3.1) we can obtain the equations
$$\begin{array}{cc}\hfill b_2=0=& \frac{1}{8}(3U^2U^{(2)})+\frac{1}{4}a_1U+\frac{1}{2}a_2\frac{1}{8}a_1^2;\hfill \\ \hfill b_3=0=& \frac{1}{32}(U^{(4)}+10U^35U^{(1)^2}10UU^{(2)})\frac{1}{16}a_1(3U^2U^{(2)})\hfill \\ & +\frac{1}{16}U(a_1^24a_2)+\frac{1}{2}a_3+\frac{1}{4}a_1a_2\frac{1}{16}a_1^3,\hfill \end{array}$$
(3.2)
which is known as the one-gap “trace formulae”. These equations form a closed system which determines the initial elliptic solutions (2.5). Inserting the rational expression for $`U(z)`$ (2.5) into the system (3.2) and equating coefficients of the Laurent expansion in $`\mathrm{}`$ in the right-hand sides to zero, we obtain algebraic relations which lead to the equalities
$$\begin{array}{cc}\hfill 1)& \alpha =2\alpha _{1_i}=\alpha _{2_i}=0,a_1=0,a_2=\frac{1}{4}g_2E,a_3=\frac{1}{4}g_3;\hfill \\ \hfill 2)& \alpha =2,\alpha _{(1,2)_i}=\frac{1}{4}\beta _{(1,2)},a_1=0,a_2=\frac{1}{4}(11g_2120\mathrm{}^2(\phi _1)),\hfill \\ & a_3=\frac{1}{4}g_3+4\mathrm{}(\phi _1)(12\mathrm{}^2(\phi _1)g_2)4\mathrm{}^{}{}_{}{}^{2}(\phi _1).\hfill \end{array}$$
(3.3)
Here and below $`\beta _1=12h^2g_2,\beta _2=2h^{}^2`$, $`\phi _1`$ is the argument of the function $`h=\mathrm{}(\phi _1)`$ which satisfies the equation $`(12h^2g_2)^2=48hh^{}{}_{}{}^{2},`$ from which $`\phi _1=(2/3)\omega _i|_{i=(1,2,4)}`$, where $`\omega _4=(\omega \omega ^{})`$ (see also ).
The two systems (3.3) determine two types of initial elliptic solutions:
$$\begin{array}{cc}\hfill 1)& U(z)=2\mathrm{}(z);\hfill \\ \hfill 2)& U(z)=2\mathrm{}(z)+2[\mathrm{}(z\phi _1)+\mathrm{}(z+\phi _1)]4\mathrm{}(\phi _1),\hfill \end{array}$$
(3.4)
where the first type is the known Lamé potential and the second type is a new potential obtained in .
### 3.2 2-gap elliptic initial solutions
Coefficients of the initial elliptic 2-gap solutions are determined by the system of the four finite-gap equations of the form (2.1) at $`n=\overline{0,3}`$. In analogy to the one-gap case, the explicit form can be obtained by substituting the expressions (2.2) for $`A_n`$ and the expressions for $`\chi _n`$ (from (2.3)) into (2.1). In doing so, the first two equations, which coincide with the first two equations of the system (3.2), are solvable with respect to $`b_1`$ and $`b_2`$. Excluding the latter from the fourth and fifth equation and taking into account the equality $`b_n|_{n3}=0`$ we obtain the finite-gap system
$$\begin{array}{cc}\hfill b_3=0=& \frac{1}{2^5}(16a_3+8a_2U+10U^35U^{}{}_{}{}^{2}2a_1U^{\prime \prime }\hfill \\ & 10UU^{\prime \prime }+U^{(4)});\hfill \\ \hfill b_4=0=& \frac{1}{2^7}(16a_2^2+64a_4+32a_3U+24a_2U^2+35U^4\hfill \\ & 70UU^{}{}_{}{}^{2}8a_2U^{\prime \prime }70U^2U^{\prime \prime }+21U^{\prime \prime }^2\hfill \\ & +28U^{}U^{(3)}+14UU^{(4)}U^{(6)}).\hfill \end{array}$$
(3.5)
Inserting the rational expression (2.5) for $`U(z)`$ into the system (3.5) and equating coefficients of the Laurent expansion in $`\mathrm{}`$, on the right- and left-hand sides, we obtain algebraic relations which lead to the equalities
$$\begin{array}{cc}\hfill 1)& \alpha =6,\alpha _{1_i}=\alpha _{2_i}=0,a_1=0,a_2=\frac{21}{4}g_2,\hfill \\ & a_3=\frac{27}{4}g_3,a_4=\frac{27}{4}g_2^2,a_4=\frac{81}{4}g_2g_3;\hfill \\ \hfill 2)& \alpha =6,\alpha _{1_i}=\delta _{i,j}(12e_j^2\frac{1}{2}g_2),a_2=7(9e_j^2+2\mathrm{\Lambda }_j),\hfill \\ & a_3=18(3e_j^35e_j\mathrm{\Lambda }_j),a_4=27(36e_j^4+16e_j^2\mathrm{\Lambda }_j+3\mathrm{\Lambda }_j^2),\hfill \\ & a_5=54(36e_j^552e_j^3\mathrm{\Lambda }_j9e_j\mathrm{\Lambda }_j^2),\hfill \\ & \mathrm{\Lambda }_j=3e_j^2\frac{1}{4}g_2;\hfill \\ \hfill 3)& \alpha =6,\alpha _{1_i}=(\delta _{i,j}+\delta _{i,k})(12e_i^2g_2),a_2=161\mathrm{\Lambda }_j378e_j^2,\hfill \\ & a_3=531e_j\mathrm{\Lambda }_j+108e_j^3,a_4=27(240\mathrm{\Lambda }_j^2+1280e_j^21159e_j^2\mathrm{\Lambda }_j),\hfill \\ & a_5=27(1594e_j\mathrm{\Lambda }_j^2+8100e_j^5+120e_j^3\mathrm{\Lambda }_j^2);\hfill \\ \hfill 4)& \alpha =6,\alpha _{1_i}=\frac{1}{4}\beta _1\delta _{i,1},\alpha _{2_i}=\frac{1}{4}\beta _2\delta _{i,1},\hfill \\ & a_2=\frac{7}{2}(18h^2+\frac{7}{2}\alpha _1),a_3=\frac{9}{4}(24h^325h^{}{}_{}{}^{2})\frac{34h}{2}\alpha _1),\hfill \\ & a_4=\frac{27}{4}(144h^4+44hh^{}{}_{}{}^{2}+56h^2\alpha _1+\frac{23}{4}\alpha _1^2),\hfill \\ & a_5=\frac{27}{2}(144h^5+6h^2h^{}{}_{}{}^{2}74h^3\alpha _1\frac{41}{2}h^{}{}_{}{}^{2}\alpha _{1}^{}\frac{75h}{2}\alpha _1^2),\hfill \end{array}$$
(3.6)
where $`\phi _2`$ is determined by the equality $`h=\mathrm{}(\phi _2)`$ which satisfies the equation $`64h^{}{}_{}{}^{4}+48hh^{}{}_{}{}^{2}\alpha _{1}^{}\alpha _1^3=0.`$ These four types of equalities (3.6) lead to the following four expressions
$`1)`$ $`U(z)=6\mathrm{}(z);`$ (3.7)
$`2)`$ $`U(z)=6\mathrm{}(z)+2\mathrm{}(z+\omega _i)2e_i;`$
$`3)`$ $`U(z)=6\mathrm{}(z)+2\mathrm{}(z+\omega _i)+2\mathrm{}(z+\omega _k)2(\mathrm{}(\omega _i)+\mathrm{}(\omega _j));`$
$`4)`$ $`U(z)=6\mathrm{}(z)+2\mathrm{}(z\phi _2)+2\mathrm{}(z+\phi _2))4\mathrm{}(\phi _2).`$
which describe all possible initial two-gap solutions of the KdV equations (see ). The first type coincides with the known two-gap Lamé potential, while the second and third type corresponds to the Treibich-Verdier potential . The fourth solution is the new two-gap potential obtained in .
### 3.3 Initial elliptic three-gap solutions
The parameters of the initial elliptic three-gap solutions (2.5) are determined by the system of the finite-band equations (2.1) at $`b_{n3}0`$ and $`b_{n>3}=0`$. Using the expressions (2.2) and (2.3) this system can be expressed through $`U`$-functions. In doing so, the first three equations coinciding formally with (3.1) and are solvable with respect to $`\overline{b_1,b_3}`$. Therefore, taking into account the equality $`b_n|_{n4}=0`$, we can obtain the equation
$`b_4=0=`$ $`16a_2^2+64a_4+32a_3U+24a_2U^2+35U^4`$
$`70UU^{}{}_{}{}^{2}8a_2U^{\prime \prime }70U^2U^{\prime \prime }`$
$`+21U^{\prime \prime }{}_{}{}^{2}+28U^{}U^{\prime \prime \prime }+14UU^{(4)}U^{(6)},`$ (3.8)
which is solvable with respect to parameters of the rational expression $`U(z)`$ (2.5). Substituting (2.5) into the equation (3.8) and equating coefficients of the Laurent expansion in $`\mathrm{}`$ with its right-hand side to zero we can obtain closed algebraic relations for $`\alpha `$-parameters of the $`U`$-solution (2.5).
It can be shown that there are four possible types of solutions for the $`\alpha `$-parameters. The corresponding four types of initial elliptic solutions can be written in the form
$$\begin{array}{cc}\hfill 1)& U(z)=12\mathrm{}(z);\hfill \\ \hfill 2)& U(z)=12\mathrm{}(z)+2\mathrm{}(z+\omega _i)2e_i,e_i=\mathrm{}(\omega _i);\hfill \\ \hfill 3)& U(z)=12\mathrm{}(z)+2(\mathrm{}(z+\omega _i)+\mathrm{}(z+\omega _j))2(e_i+e_j);\hfill \\ \hfill 4)& U(z)=12\mathrm{}(z)+2(\mathrm{}(z+\phi _3)+\mathrm{}(z\phi _3))4\mathrm{}(\phi _3).\hfill \end{array}$$
(3.9)
Here the argument $`\phi _3`$ is determined by the equation
$$h^6+\frac{101}{196}g_2h^4+\frac{29}{49}g_3h^3\frac{43}{784}g_2^2h^2\frac{23}{196}g_2g_3h\left(\frac{1}{3136}g_2^3+\frac{5}{98}g_3^2\right)=0$$
where $`h=\mathrm{}(\phi _3)`$. The initial solutions 1 and 2, 3 are the 3-gap Lamé () and generated on the 3-gap case Treibich-Verdier potentials, respectively. The initial elliptic solution 4 is the 3-gap generalization of the solution obtained in .
Note that the above described algorithm is general and applicable to computing arbitrary initial elliptic $`n`$-gap solutions of the KdV equation.
## 4 A dynamics of the elliptic finite-gap solutions
The time dependent elliptic finite-gap solutions of the KdV equation have the general form of the rational functions of the $`\mathrm{}`$-function (2.4), parameters of which are functions of time $`t`$. These parameters are described by the system of the auxiliary dynamic equation (2.6) and the finite-gap equation (2.1).
By substituting the expression (2.4) for the elliptic finite-gap solutions $`U(z,t)`$ into the KdV equation $`_tU(z,t)=6U(z,t)U^{}(z,t)U^{\prime \prime \prime }(z,t)`$ and equating the coefficients of the Laurent expansion in $`\mathrm{}`$, in its left- and right-hand sides, lead to the known general formula
$$U(z,t)=2\underset{i=1}{\overset{N}{}}\mathrm{}(z\phi _i(t))+C,$$
(4.1)
in which the integer number $`N`$ and the constant $`C`$ are determined by the condition of the reduction of $`U(z,t)`$ to the corresponding initial elliptic finite-gap solution at $`t0`$. In the cases of the elliptic 1-, 2- and 3-gap solutions the numbers $`N`$ must provide the reduction of the time dependent solutions (4.1) to the corresponding initial elliptic solutions of the systems (3.4), (3.7) and (3.9), respectively.
The substitution of the expression (4.1) into the KdV equation reduces the problem of the time evolution of $`U(z,t)`$ to the time evolution its poles. The latter is described by the system
$$\begin{array}{cc}& _t\phi _i(t)=12X_i(t)+C,X_i(t)=\underset{j=1,ji}{\overset{N1}{}}\mathrm{}(\phi _i(t)\phi _j(t))(g2),\hfill \\ & _t\underset{n=1}{\overset{N}{}}\mathrm{}(z\phi _i(t))=0(g=1),\hfill \end{array}$$
(4.2)
which is a system of coupled equations. However, in view of the symmetry properties of the finite-gap equations, (2.1) determine $`X_i(t)`$ as function $`X_i(\phi _i(t))`$. Then the system (4.2) can be transformed to
$$_{\phi _{0i}}^{\phi _i}\frac{d\phi _i}{X_i(h_i(\phi _i))+C}=12t,$$
(4.3)
which describe the time dynamics of the poles in the expression (4.1). Here, initial values $`\phi _{0i}\phi _i(0)`$ are determined from the expressions for initial elliptic finite-gap solutions.
The functions $`X_i(\phi _i)`$ are determined by the finite-gap equations (2.1) with $`U`$-functions in the form (4.1) as roots of $`N`$th order polynomials in $`X_i`$. These polynomials are followed from the algebraic equations which can be obtained by equating coefficients of the Laurent expansion in $`\mathrm{}`$ of the left- and right-hand sides of the indicated equations (2.1).
Thus, the problem of time evolution of the elliptic finite-gap solutions of the KdV equations is reduced to the solution of the finite-band equations, with respect to the $`X_i`$-functions from relations (4.3). The proposed approach will now be applied in calculating the time evolution of elliptic 1-, 2- and 3-gap solutions.
### 4.1 Elliptic 1-gap solutions
The types of time dependent elliptic 1-gap solutions of the KdV equation are determined by the values of the number $`N`$ in the expression (4.1). The condition of a coincidence in the general expression (4.1) with the initial elliptic 1-gap solutions (3.4) at $`t0`$, yields the values, namely $`N=1`$ and $`N=3`$. The corresponding elliptic 1-gap solutions have the form
$$\begin{array}{cc}\hfill 1)& U(z,t)=2\mathrm{}(z\phi _1(t)),\hfill \\ \hfill 2)& U(z,t)=2\underset{1}{\overset{3}{}}\mathrm{}(z\phi _i(t))4\mathrm{}(\phi _1).\hfill \end{array}$$
(4.4)
The time dependence of poles of the elliptic solutions 1 and 2 is determined by the second equation of the system (4.2) at $`N=1`$ and $`N=2`$, respectively. In accordance with the initial conditions defined by the system (3.4), the solutions of the dynamic equation have the form $`\phi _1=c_1t`$ at $`N=1`$ and
$$\phi _1(t)=c_1t,\phi _2(t)=c_2t+\phi _i^{(1)},\phi _3(t)=c_3t\phi _i^{(1)}$$
at $`N=3`$. The substitution of these expressions into (4.4) yields the following two expressions
$$\begin{array}{cc}\hfill 1)& U(z,t)=2\mathrm{}(zc_1t);\hfill \\ \hfill 2)& U(z,t)=2\{\mathrm{}(zc_1t)+\mathrm{}(zc_2t+\phi _i^{(1)})+\mathrm{}(zc_3t\phi _i^{(1)})4\mathrm{}(\phi _1),\hfill \end{array}$$
which describes two possible types of the elliptic 1-gap solutions in the form of superpositions of one and three independent traveling waves, respectively.
### 4.2 Elliptic 2-gap solutions
The possible types of the time dependent elliptic 2-gap solutions of the KdV equation (4.1) are determined by the values $`N`$ which can be obtained from the compatibility condition between (4.1) and (3.7) as $`t0`$. Under this condition the number $`N`$ equal 3, 4 and 5 in the formula (4.1).
1. The time dependent 2-gap elliptic solution corresponding to the initial condition 1 in the system (3.7), which is described by the expression (4.1) at $`N=3`$, has the form
$$U(z,t)=2\underset{i=1}{\overset{3}{}}\mathrm{}(z\phi _i(t)).$$
(4.5)
The time evolution of poles $`\phi _i(t)`$ is described by the equation (4.3) in which $`X_i=_{ji,j=1}^2\mathrm{}(\phi _i\phi _j)`$. Under initial conditions, the lower limits of the integration in (4.3) are $`\phi _{0i}=0`$.
To comput the $`X_i`$-function we substitute the expression (4.5) into the finite-gap equation (3.2). Then, equating coefficients of the Laurent expansion in $`\mathrm{}`$ of the left- and right-hand sides, we obtain algebraic equations which are reduced to the polynomial equation
$$X^3+c_2X^2+c_1X+c_0=0,$$
(4.6)
where $`c_n=c_n(\phi ),n=\overline{1,3}`$ (here and below the subscript $`i`$ of $`\phi `$ is omitted). Three solutions of (4.6) describe three unknown functions $`X_i(\phi _i),i=\overline{1,3}`$. The dependence of the coefficient functions $`c_n`$ on $`\phi `$ is described by the expressions
$$\begin{array}{cc}& c_0=\frac{36}{125}g_3\frac{2}{125}a_2\frac{1}{250}(149g_2+76a_2)h+\frac{1}{400}(g_2+4a_2)\left(\frac{\beta _1}{h^{}}\right)^2;\hfill \\ & c_1=\frac{29}{250}g_2+\frac{8}{125}a_2;\hfill \\ & c_2=\frac{42}{125}h+\frac{3}{800}\left(\frac{\beta _1}{h^{}}\right)^2,\hfill \end{array}$$
where $`\beta _1=12h^2g_2`$ and $`h=\mathrm{}(\phi )`$ so that $`X_i`$ depend on $`\phi `$ through the $`h`$-function.
2. The time dependent elliptic 2-gap solution corresponding to the condition 2 in the system (3.7) is described by the expression (4.1) at $`N=4`$ and has the form
$$U(z,t)=2\underset{i=1}{\overset{4}{}}\mathrm{}(z\phi _i(t))2\mathrm{}(\phi _{04}).$$
(4.7)
The time evolution of poles $`\phi _i(t)`$ are described via the function $`X_i=_{ji,j=1}^3\mathrm{}(\phi _i\phi _j)`$. In view of the initial conditions, the lower limits of an integration in (4.3) are $`\phi _{0i}|_{i3}=0`$ and $`\phi _{04}=\omega _j`$.
For computing the $`X_i`$-functions we substitute the expression (4.7) into the system (3.5). Then, equating coefficients of the Laurent expansion of the left- and right-hand sides, we obtain the equation
$$X^4+c_3X^3+c_2X^3+c_1X+c_0=0,$$
(4.8)
solutions of which describe the functions $`X_n(\phi ),n=\overline{1,4}`$. The dependence of the coefficient functions $`c_n`$ on $`\phi `$ are described by the expressions
$$\begin{array}{cc}\hfill c_0=& \stackrel{~}{m}_1^0g_2^2+\stackrel{~}{m}_2^0g_2a_2+\stackrel{~}{m}_3^0a_2^2\stackrel{~}{m}_4^0a_4;\hfill \\ \hfill c_1=& \stackrel{~}{m}_1^1g_3+\stackrel{~}{m}_2^1a_3+\stackrel{~}{m}_3^1g_2h+\stackrel{~}{m}_4^1a_2h\stackrel{~}{m}_5^1g_2\left(\frac{\beta _1}{h^{}}\right)^2\stackrel{~}{m}_6^1a_2\left(\frac{\beta _1}{h^{}}\right)^2;\hfill \\ \hfill c_2=& \stackrel{~}{m}_1^2g_2\stackrel{~}{m}_2^2a_2;c_3=\stackrel{~}{m}_1^3h+\stackrel{~}{m}_2^3\left(\frac{\beta _1}{h^{}}\right)^2.\hfill \end{array}.$$
(4.9)
Here $`\stackrel{~}{m}_j^i`$ denotes numerical parameters which have the form of some rational fractions.
3. The time dependent elliptic 2-gap solutions, corresponding to the initial conditions 3 and 4 in the system (3.7) which are described by the expression (4.1) at $`N=5`$, have the form
$$U(z)=2\underset{i=1}{\overset{5}{}}\mathrm{}(z\phi _i(t))2(\mathrm{}(\phi _{04})+\mathrm{}(\phi _{05})).$$
(4.10)
In view of the initial conditions, the lower limits of the integration in (4.3) have the common values $`\phi _{0i}|_{i3}=0`$ and $`\phi _{04}=\omega _i,\phi _{05}=\omega _j`$ at the condition 3 and $`\phi _{04}=\phi _{05}=\phi _2`$ at the condition 4.
For computing the $`X_i`$-functions we substitute the expression (4.10) into the system (3.5). Then, equating coefficients of the Laurent expansion in $`\mathrm{}`$ of the left- and right-hand sides, we obtain the equation
$$X^5+c_3X^3+c_2X^2+c_1X+c_0=0,$$
solutions of which describe functions $`X_n(\phi ),n=\overline{1,5}`$. Coefficient functions $`c_n`$ are described by the expressions
$$\begin{array}{cc}\hfill c_0=& \underset{i=1}{\overset{5}{}}(m_i^0g_3+n_i^0a_3)F_i(\phi )+m_6^0g_2g_3+m_7^0g_3a_2m_8^0g_2a_3\hfill \\ & +m_9^0a_2a_3m_{10}^0a_5;\hfill \\ \hfill c_1=& \underset{i=1}{\overset{5}{}}(m_i^1g_2+n_i^1a_2)F_i(\phi )m_6^1g_2^2m_7^1g_2a_2+m_8^1a_2^2+m_9^1a_4;\hfill \\ \hfill c_2=& m_1^2g_3+m_2^2a_3;c_3=\underset{i=1}{\overset{5}{}}m_i^3F_i(\phi )m_2^3g_2+m_3^3a_2,\hfill \end{array}$$
where
$$F_i(\phi )\left(\delta _{i,1}\beta _1+\delta _{i,2}h^2+\delta _{i,3}\left(\frac{\beta _1}{h^{}}\right)^4+\delta _{i,4}h\left(\frac{\beta _1}{h^{}}\right)^2+\delta _{i,5}\left(\frac{\beta _1h^{(4)}}{h^{}^2}\right)\right).$$
Here $`m_i^j,n_i^j`$ denotes some numerical rational fractions, $`\beta _1=12\mathrm{}^2(\phi )g_2,h=\mathrm{}(\phi )`$, $`a_i=a_i(\mathrm{}(\phi ))`$ denote some complicate functions, the explicit form of which we do not present here.
### 4.3 Elliptic 3-gap solutions
The possible types of the time dependent elliptic 3-gap solutions of the KdV equation (4.1) are determined by the values $`N`$ which are obtained from the compatibility condition between (4.1) and (3.9) as $`t0`$. Under this condition, the number $`N`$ takes values $`\overline{6,8}`$ in the formula (4.1).
1. The values $`N=6`$ and $`N=7`$ determine two elliptic 3-gap solutions of the KdV equations with initial conditions 1 and 2 of the system (3.9), which have the form
$$U(z,t)=2\underset{i=1}{\overset{6}{}}\mathrm{}(z\phi _i(t)),\mathrm{and}U(z,t)=2\underset{i=1}{\overset{7}{}}\mathrm{}(z\phi _i(t))2\mathrm{}(\phi _{07}),$$
(4.11)
respectively. The time evolution of the poles $`\phi _i(t)`$ are described by relations (4.3) with the lower integration limits $`\phi _{0,i}|_{i=\overline{1,6}}=0`$ at the initial condition 1 and $`\phi _{0,i}|_{i=\overline{1,6}}=0,\phi _{0,7}=\omega _i`$ at the initial condition 2.
By substituting the expressions (4.11) into the finite-band equation (3.8) and equating coefficients of the Laurent expansion in $`\mathrm{}`$ of the right-hand sides to zero, lead to the two equations
$$\underset{i=0}{\overset{6}{}}c_{6,i}X^i=0,\mathrm{and}\underset{i=0}{\overset{7}{}}c_{7,i}X^i=0,$$
(4.12)
corresponding to the two 3-gap solutions with the values $`N=6`$ and $`N=7`$, respectively. Coefficient functions $`c_{6,i}`$ and $`c_{7,i}`$ in (4.12) are rational functions on $`\mathrm{}(\phi _i)`$ and $`\mathrm{}^{}(\phi _i)`$. Therefore $`X_i`$ as roots of (4.12) are functions of $`\phi _i`$, i.e. $`X_i=X_i(\phi _i)`$ where $`i=\overline{1,6}`$ and $`i=\overline{1,7}`$ at initial conditions 1 and 2, respectively.
2. The value $`N=8`$ determines two similar elliptic 3-gap solutions of the KdV equations with the initial conditions 3 and 4 of the system (3.9), which have the form
$$U(z,t)=2\underset{i=1}{\overset{8}{}}\mathrm{}(z\phi _i(t))2(\mathrm{}(\phi _{07})+\mathrm{}(\phi _{08})).$$
(4.13)
The poles $`\phi _i(t)`$ of (4.13) are described by relations (4.3), with the lower integration limits $`\phi _{0i}|_{i=\overline{1,6}}=0,\phi _{07}=\omega _i,\phi _{08}=\omega _j`$ and $`\phi _{0i}|_{i=\overline{1,6}}=0,\phi _{07}=\phi _{08}=\phi _3`$ at initial conditions 3 an 4, respectively.
Substituting the expressions (4.13) into the finite-gap equation (3.8) and equating the coefficients of the Laurent expansion in $`\mathrm{}`$ of the right-hand sides to zero, we obtain two equations of the form
$$\underset{i=0}{\overset{8}{}}c_{6,i}X^i=0,$$
(4.14)
where $`c_i=c_i(\phi )`$. Eight solutions of (4.14) coincide with eight functions $`X_i(\phi _i),i=\overline{1,8}`$.
The proposed approach is applicable for computing arbitrary elliptic $`n`$-gap solutions of the KdV equation. It can also be applied for computing finite-gap elliptic solutions for other integrable equations.
## 5 Conclusion
The solution of the KdV equation in the class of elliptic finite-gap functions is reduced to the solution of the system of finite-gap equations and auxiliary dynamic equations. In terms of the elliptic $`\mathrm{}`$-function this system reduces to simple algebraic relations which determin the parameters of the unknown solutions which are represented as rational functions of the $`\mathrm{}`$-function. This approach gives a simple algorithm for calculating arbitrary elliptic finite-gap solutions of the KdV equations at an initial time, which was demonstrated by the example of 1-, 2- and 3-gap solutions. The time evolutions of the unknown solutions with a known $`\mathrm{}`$-functional representation, are determined by the dynamics of their poles, which is described by coupled systems of dynamic equations. The latter always can be integrated with the help of the finite-gap equations. This was demonstrated by the example of elliptic 1-, 2- and 3-gap solutions.
The above approach will also be applied to other integrable nonlinear equations in a future paper.
### Acknowledgements
This research was supported by the CDRF grant UM1-325. |
warning/0001/hep-ph0001109.html | ar5iv | text | # References
Brane-Universe in Six Dimensions with Two Times
Merab GOGBERASHVILI
Institute of Physics, Georgian Academy of Sciences
6 Tamarashvili Str., Tbilisi 380077, Georgia
(E-mail: gogber@hotmail.com)
Abstract
> Brane-Universe model embedded in 6-dimensional space-time with the signature (2+4) is considered. A matter is gravitationally trapped in three space dimensions, but both time-like directions are open. Choosing of the dimension and the signature of the model is initiated with the conformal symmetry for massless particles and any point in our world can be (1+1) string-like object.
PACS number: 98.80.Cq
In conventional Kaluza-Klein’s picture extra dimensions are curled up to unobservable size. Last years models with large extra dimensions become popular (see for example ). Those approaches also do not contradict to present time experiments . Main obstacles in progress of such models were how to confine a matter inside brane-Universe and to explain observed 4-dimensional Newton’s law. In the papers we had introduced possible mechanism of overcoming these problems. Special solution of multi-dimensional Einstein’s equations , responsible for gravitational trapping of a matter, provides effective 4-dimensional Newton’s law on the brane. In this model it can be explained also non-observation of cosmological constant in four dimensions .
In our previous papers for simplicity 5-dimensional model with the signature (1+4) was considered. In general the question about the number of dimensions and signature is open. We consider that a matter is trapped on the brane by gravitation. So it is natural to assume that for massless case (it means weakest coupling with gravity) symmetries of sub-manifold can be restored. It is well known that in the zero-mass limit main equations of physics are invariant under fifteen parameter nonlinear conformal transformations. From the other hand a long time ago it was discovered that conformal group can be written as a linear Lorentz type transformations in 6-dimensional space-time with the signature (2+4) (for these subjects see for example ). Another indication that real number of space-like dimensions of Universe can be four, could be $`O(4)`$ symmetry of the solution of Shrodinger equations for the hydrogen atom.
Theories with extra time-like dimensions have been a subject of interest for some time (the latest article about large extra time-like dimensions is , see also ). Kaluza-Klein 6-dimensional model with two times, but with compact extra dimensions, was investigated before in . Necessity of two times follow also from string theory (F-theory) . There exist another two groups of 6-dimensional schemes, but with three times. The schemes of the first group suffer from internal inconsistencies, second , more sophisticated scheme is internally consistent.
In all these papers multi time dimensions are introduced, however, not much is known about the classical vacuum solutions of two- or more-timing theories. It is known that theories with compact internal time-like dimensions have several pathological features. The most conspicuous may be the fact that excitations of the internal dimensions have negative norm. Experimental lower bounds on possible violation of unitarity put a limit on the maximum radius of the internal time-like directions .
After unification of space and time coordinates (but not dimensions) in Lorents transformations many physicists intuitively are accustomed to consider space and time dimensions identically. However, there is principal difference among them. Main difference is that one can easily change position, or stop in the space but everybody follows to time flow which was began from the Big Bang. Decreasing of time flow, or shift of time vector for any system, means backward moving in the time and then disappearance for other observers. Another problem is that we know several space directions and it is easy to add next ones, but we know only one time. In case of two or more time directions the question arises what we are measuring with our clocks. Thus ordinary methods of trapping using in our previous papers can be not applicable for time-like dimensions.
Possible indication of existence of extra time directions can be non-conservation of the energy in four dimensions. This can be happening only after interaction with the matter with another direction of time (or energy) vector. As we mentioned above we only passively follow to time flow, direction of which coincides with cosmological arrow of time (possible source of different asymmetries in our world, see for example ). So it could work mechanism similar to that considered in papers \- after some time from the Big Bang all particles with another directions of time had been disappeared from our view. In this case may be no special trapping mechanism in extra times is necessary.
All points of the space in Universe are equivalent. However, we have global zero for time coordinate, it is the moment of the Big Bang. What we measure for any particle is the difference of two energies - energy of the vacuum and the particle itself. When particle follows to our time flow we notice only this difference, but to shift the time vector we need total value of the energy corresponding to the age of Universe. Thus if the particles with another energy vector had been disappeared the matter of our world is trapped in our time by Plank scale.
In this paper we consider 6-dimensional Kaluza-Klein model with one space-like and one time-like extended extra dimension. It is generalization of 5-dimensional scheme of papers with one brane where all matter is localized. Trapping mechanism considered in these articles works for the case of one extra space-like dimension and is applicable for 6-dimensional model with two times. Time-like dimensions in our scheme are open, similar to the approach of .
We are looking for solution of 6-dimensional Einstein’s equations with the cosmological term $`\mathrm{\Lambda }`$ in two time- and four space-like dimensions
$$^6R_{AB}\frac{1}{2}g_{AB}{}_{}{}^{6}R=\mathrm{\Lambda }g_{AB}+G\tau _{AB}.$$
(1)
Here $`{}_{}{}^{6}R_{AB}^{}`$ and $`G`$ are 6-dimensional Ricci tensor and gravitational constant and big Latin indices $`A,B,\mathrm{}=0,1,2,3,5,6`$.
Energy-momentum tensor for brane-Universe with the signature (2+3) embedded in 6-dimensional space-time with signature (2+4) is taken in the form
$$\tau _{\mu \nu }=g_{\mu \nu }\sigma \delta (\frac{x^6}{ϵ}),\tau _{55}=\sigma \delta (\frac{x^6}{ϵ}),\tau _{66}=0.$$
(2)
Here $`\sigma `$ is brane tension, $`\delta (x^6/ϵ)`$ \- delta function and $`ϵ\sqrt{\mathrm{\Lambda }}`$ \- brane width in extra space-like dimension $`x^6`$. Greek indices $`\alpha ,\beta \mathrm{}=0,1,2,3`$ numerate coordinates in four dimensions. This form of the brane energy-momentum follows from kink-like solution of scalar field equations without introducing of coupling with gravity. In general it is difficult to separate the energy of the brane and gravitation field and we must consider more complicate model.
In (1) we choose negative sign for cosmological constant $`\mathrm{\Lambda }`$. Canceling mechanism of papers for space-like extra dimension works only for the negative $`\mathrm{\Lambda }`$. Solution with positive $`\mathrm{\Lambda }`$ corresponds to trapping in extra time and is considered at the end of the paper.
In many Kaluza-Klein models with large extra dimensions localization of particular quantum fields on the brane is investigated (latest paper in this direction is ). In our model matter is trapped on the brane by gravitation. The source of anti-gravity responsible for this trapping can be negative multi-dimensional cosmological constant. Since gravity is universal field we don’t need any other classical source for localization of a matter on the brane. Exact mechanism for different fields is difficult to find because of problems with consideration of quantum fields in curve space-time (see for example ). Even in four dimensions only few exactly solvable models exist. Our approach is more general. Gravitation field is not localized on the brane, however, by canceling mechanism we have ordinary Newton’s law. Gravitational potential has minimum on the brane and all particles having coupling with gravity are sitting there.
To keep the width of brane-Universe during expansion, it means for the stabile localization of the mater on the brane, as in paper we look for the solution of (1) with zero extra momentum
$$P_i=T_i^A𝑑S_A=0.$$
(3)
Small Latin indices $`i,j,k,\mathrm{}=5,6`$ numerate coordinates of extra dimensions. Here $`T_A^B=t_A^B+\tau _A^B`$ is total energy momentum tensor of matter fields $`\tau _A^B`$ and gravitational field itself
$$t_A^B=\frac{1}{2G}[g^{BD}_A\mathrm{\Gamma }_{DE}^Eg^{ED}_A\mathrm{\Gamma }_{DE}^B+\delta _A^B({}_{}{}^{6}R2\mathrm{\Lambda })].$$
(4)
Type of matter fields is not important now. Some particular cases are considered in the paper .
In this article we don’t want to touch the old problem with the energy of gravitation and as in the paper choose $`T_A^B`$ in the form of so called Lorentz energy-momentum complex
$$T_A^B=\frac{1}{2G\sqrt{g}}_CX_A^{BC},$$
(5)
where
$$X_A^{BC}=X_A^{CB}=\sqrt{g}[g^{BD}g^{CE}(_Dg_{AE}_Eg_{AD})].$$
(6)
This form is convenient, since in this case energy-momentum tensor of gravitational field $`t_A^B`$ coincides with canonical energy-momentum tensor (4) received from Hilbert’s form of the gravitational Lagrangian $`L_g=\sqrt{g}({}_{}{}^{6}R2\mathrm{\Lambda })`$.
To satisfy the stability condition (3), components of the energy-momentum tensor on the solutions must satisfy the condition
$$T_i^A=0.$$
(7)
From $`(i\alpha )`$ component of this relation and (4) we find $`_i\mathrm{\Gamma }_{\mu \nu }^\alpha =0`$. Thus simple solution of (7) is
$$g_{i\alpha }=0,g_{\alpha \beta }=\lambda (x^i)\eta _{\alpha \beta }(x^\nu ),$$
(8)
where $`\eta _{\alpha \beta }(x^\nu )`$ is ordinary 4-dimensional metric tensor and $`\lambda (x^i)`$ is arbitrary function of extra coordinates. Solution (8) which we received from stability conditions, is similar with the anzats of .
Stability condition (7) for the case of diagonal metric tensor of extra dimensions $`(g_{56}=0)`$ has the form
$$_5(g^{55}g^{66}_5g_{66})=_6(g^{55}g^{66}_6g_{55})=0.$$
(9)
One of the solutions of this system is
$$g_{55}=exp(cx^6),g_{66}=1,$$
(10)
where $`c`$ is integration constant.
Using (8) one can find decomposition of Einstein’s equations (see also )
$`R_{\alpha \beta }{\displaystyle \frac{1}{2}}g_{\alpha \beta }(D_iD^i\lambda +{\displaystyle \frac{1}{2\lambda }}D_i\lambda D^i\lambda )={\displaystyle \frac{1}{2}}g_{\alpha \beta }[\mathrm{\Lambda }{\displaystyle \frac{1}{2}}G\sigma \delta ({\displaystyle \frac{x^6}{ϵ}})],`$
$`R_{ij}{\displaystyle \frac{2}{\lambda }}(D_iD_j\lambda {\displaystyle \frac{1}{2\lambda }}D_i\lambda D_j\lambda )={\displaystyle \frac{1}{2}}g_{ij}[\mathrm{\Lambda }{\displaystyle \frac{5}{2}}G\sigma \delta ({\displaystyle \frac{x^6}{ϵ}})].`$ (11)
Here $`R_{ij}`$ and $`D_i`$ are correspondingly Ricci tensor and the covariant derivative in extra space-time with the metric tensor $`g_{ij}`$. Ricci tensor in four dimensions $`R_{\alpha \beta }`$ is constructed from $`g_{\alpha \beta }=\lambda (x^i)\eta _{\alpha \beta }(x^\nu )`$ in a standard way.
Using (10) and the properties of the step function $`H(x^5)`$
$$|x|^{}=H(x)H(x),H(x)^{}=\delta (x)=\frac{1}{|ϵ|}\delta (\frac{x}{ϵ)}$$
(12)
where prime denotes derivative, one can show that system (S0.Ex1) has the trapping solution
$$\lambda =g_{55}=exp(c|x^6|),g_{66}=1.$$
(13)
The integration constant here has the value
$$c=\sqrt{\frac{2\mathrm{\Lambda }}{5}}=\frac{G\sigma ϵ}{4}.$$
(14)
This formula also contains necessary relation between the brane tension $`\sigma `$ and 6-dimensional cosmological constant $`\mathrm{\Lambda }`$.
For this solution in four dimensions we have ordinary Einstein’s equations without the cosmological term
$$R_{\alpha \beta }\frac{1}{2}\eta _{\alpha \beta }R=0,$$
(15)
which is function of only 4-dimensional metric tensor $`\eta _{\alpha \beta }(x^\nu )`$. After adding of 4-dimensional source at the right hand of (15) one can show that as in 5-dimensional case a matter is trapped in three space with ordinary Newton’s low, while now we have unobservable extra open time direction.
From (14) we see that constant $`c`$ is real only for our choice of sign of $`\mathrm{\Lambda }`$. Also we noticed that for branes with positive tension constant $`c`$ is negative and function $`\lambda `$ decreasing fare from the brane, as in papers . For the case of negative $`\sigma `$ exponential factor in (13) is positive as in papers and gravitational potential has minimum on the brane.
In this paper and our previous articles , in contrast with the approach of , we consider only one brane. Interactions of branes with the negative and positive tensions are often considered in Kaluza-Klein theories with the large extra dimensions. One must be careful in this case. It is known for a long time that system of negative and positive masses began to accelerate till the speed of the light. Acceleration can destroy the branes. Even one brane with positive tension is strange objects, since it is gravitationally repulsive . This can cause change of the time direction on the brane , while $`t^2`$ is still positive. Negative tension can change signature and thus interchanges time and space coordinates.
System (S0.Ex1) with positive $`\mathrm{\Lambda }`$ has similar to (13) solution
$$\lambda =g_{66}=exp(c|x^5|),g_{55}=1,$$
(16)
which corresponds to trapping in extra time, while all four space-like coordinates are open.
At the end of the paper we want to note that in this paper 6-dimensional model with string-like extra dimension was considered. Any point-like particle in our world can have tail in (1+1) dimensions and we have interesting possibility for nontrivial application of string theory. |
warning/0001/hep-ph0001042.html | ar5iv | text | # Non-leptonic hyperon weak decays in the Skyrme model revisited
## I Introduction
Although during the last few decades much progress has been done in the theoretical study of hadron structure, the non-leptonic weak decays of hyperons still remain far from being well understood. This class of decays involve not only weak interactions but also low momentum strong processes, which have made their first principle calculation unfeasible so far. In this situation, different hadron models have been used to get the corresponding theoretical predictions. The available experimental information allows to determine both $`S`$\- and $`P`$-wave decay amplitudes separately for various processes. In the case of the $`S`$-wave decays, the predictions given by quark models with QCD enhancement factors turn out to be quite successful (see Ref. and references therein). These models, however, have serious difficulties in reproducing at the same time the empirical results in the case of the $`P`$-wave amplitudes. Indeed, this seems to be a problem (so called “$`S/P`$ wave puzzle”) which is common to other approaches as e.g. heavy baryon chiral perturbation theory, QCD sum rules, etc.
A possible solution to the $`S/P`$ wave puzzle has been suggested some years ago within the context of chiral topological soliton models. It was shown that in these models the $`P`$-wave amplitudes receive, in addition to the standard pole diagrams, extra contributions from contact terms. Then, it was speculated that these extra terms could provide a clue to this issue. Unfortunately, at the time this suggestion was made such models were hampered by several serious problems, such as very poor predictions for the hyperon spectrum, far too small results for the $`S`$-wave non-leptonic decay amplitudes, etc. Therefore, it was hard to draw definite conclusions about the real relevance of the contact terms. With the introduction of more refined methods to treat the chiral symmetry breaking terms in the effective action the situation has significantly improved. Indeed, a scheme in which one introduces $`SU(3)`$ collective coordinates and the hamiltonian —including a symmetry breaking piece— is diagonalized exactly leads to very good results for the hyperon spectra, as well as reasonable predictions for different baryon properties (for a review see Ref.). Moreover, it has been recently shown that within such scheme, and using a simple Cabbibo current-current form for the weak interaction, one can obtain the correct $`S`$-wave absolute values. Thus, we are now in position to verify whether chiral soliton models can provide a unified and consistent description of the hyperon non-leptonic decays. In this work we will describe the weak interactions by means of an effective weak chiral lagrangian —which is more general than the Cabbibo current-current coupling often used in previous Skyrme model calculations —, where the low energy coupling constants will be fixed to reproduce the known $`2\pi `$ and $`3\pi `$ $`K`$-meson decays. We will concentrate mostly in the dominant octet-like piece of this $`\mathrm{\Delta }S=1`$ lagrangian, considering terms up to order $`p^4`$. In addition, possible non-octet contributions will be considered for the particular case of the $`𝒜(\mathrm{\Sigma }_+^+)`$ amplitude which, as well known, vanishes in the pure octet approximation. The non-leptonic hyperon decay amplitudes will be obtained by evaluating the corresponding matrix elements using the topological soliton model wave functions.
The article is organized as follows: in Sec. II we give a brief overview of the $`SU(3)`$ soliton model and introduce the octet-like weak effective chiral lagrangian to be used in the following two sections. In Sec. III and IV we describe the calculation of the $`S`$-wave and $`P`$-wave amplitudes respectively, and present the corresponding results. The impact of the non-octet-like components of the weak effective lagrangian is discussed in Sec. V and in Sec. VI we state our conclusions. Finally, in the Appendix we give some details concerning the evaluation of the matrix elements of the collective operators which appear in the calculation of the decay amplitudes.
## II The model
In chiral topological soliton models baryons are described as topological excitations of a chiral effective action which depends only on meson fields. We use the form
$$\mathrm{\Gamma }=\mathrm{\Gamma }_{SK}+\mathrm{\Gamma }_{WZ}+\mathrm{\Gamma }_{SB},$$
(1)
where $`\mathrm{\Gamma }_{SK}`$ and $`\mathrm{\Gamma }_{WZ}`$ stand for the Skyrme and Wess-Zumino actions respectively, and $`\mathrm{\Gamma }_{SB}`$ is an $`SU(3)`$ symmetry breaking piece. The Skyrme action has the usual form
$$\mathrm{\Gamma }_{SK}=d^4x\left\{\frac{f_\pi ^2}{4}\text{Tr}\left[_\mu U(^\mu U)^{}\right]+\frac{1}{32ϵ^2}\text{Tr}\left[[U^{}_\mu U,U^{}_\nu U]^2\right]\right\},$$
(2)
where the chiral field $`U`$ is a non–linear realization of the pseudoscalar octet, $`f_\pi =93`$ MeV is the pion decay constant and $`ϵ`$ is the dimensionless Skyrme parameter. The Wess-Zumino action reads
$`\mathrm{\Gamma }_{WZ}`$ $`=`$ $`{\displaystyle \frac{iN_c}{240\pi ^2}}{\displaystyle d^5xϵ^{\mu \nu \rho \sigma \tau }\text{Tr}[L_\mu L_\nu L_\rho L_\sigma L_\tau ]},`$ (3)
where $`L_\mu =U^{}_\mu U`$ and $`N_c=3`$ is the number of colors. Finally the symmetry breaking piece $`\mathrm{\Gamma }_{SB}`$ is given by
$`\mathrm{\Gamma }_{SB}`$ $`=`$ $`{\displaystyle }d^4x\{{\displaystyle \frac{f_\pi ^2m_\pi ^2+2f_K^2m_K^2}{12}}\text{Tr}[U+U^{}2]+\sqrt{3}{\displaystyle \frac{f_\pi ^2m_\pi ^2f_K^2m_K^2}{6}}\text{Tr}\left[\lambda _8(U+U^{})\right]`$ (5)
$`+{\displaystyle \frac{f_K^2f_\pi ^2}{12}}\text{Tr}\left[(1\sqrt{3}\lambda _8)(U(_\mu U)^{}^\mu U+U^{}_\mu U(^\mu U)^{})\right]\},`$
where $`f_K`$ is the kaon decay constant and $`m_\pi `$ and $`m_K`$ are the pion and kaon masses, respectively. In our numerical calculations below we will set these parameters to their empirical values and take $`ϵ=4.1`$ which is suitable for a good description of many baryon properties in this model.
In the soliton picture we are using the strong interaction properties of the low–lying $`\frac{1}{2}^+`$ and $`\frac{3}{2}^+`$ baryons are computed following the standard $`SU(3)`$ collective coordinate approach to the Skyrme model. We introduce for the chiral field the ansatz
$$U_0(𝐫,t)=A(t)\left(\begin{array}{cc}c𝟙+𝕚𝝉\widehat{𝒓}𝕤& 0\\ 0& 1\end{array}\right)A^{}(t),$$
(6)
where we have used the abbreviations $`c=\mathrm{cos}F(r)`$ and $`s=\mathrm{sin}F(r)`$, $`F(r)`$ being the chiral angle that parameterizes the soliton. The collective rotation matrix $`A(t)`$ is $`SU(3)`$ valued. Substituting the configuration Eq. (6) into $`\mathrm{\Gamma }`$ yields (upon canonical quantization of $`A`$) the collective Hamiltonian. Its eigenfunctions are identified as the baryon wavefunctions $`\mathrm{\Psi }_B`$. Due to the symmetry breaking piece $`\mathrm{\Gamma }_{SB}`$, the Hamiltonian is obviously not $`SU(3)`$ symmetric. However, as shown by Yabu and Ando , it can be diagonalized exactly. The diagonalization essentially amounts to admixtures of states from higher $`SU(3)`$ representations into the octet ($`J=\frac{1}{2}`$) and decouplet ($`J=\frac{3}{2}`$) states. This procedure has proven to be quite successful in describing the hyperon spectrum and static properties .
In order to describe the non-leptonic hyperon decays we have to introduce an effective weak $`\mathrm{\Delta }S=1`$ lagrangian. The latter is constrained by weak interactions to transform either as $`\underset{¯}{8}`$ or $`\underset{¯}{27}`$ under the chiral group $`SU(3)_L`$. Here, we will take into account only the dominant octet-like couplings, which lead to pure $`\mathrm{\Delta }I=1/2`$ transitions. The remaining $`\underset{¯}{27}`$ piece includes both $`\mathrm{\Delta }I=1/2`$ and $`\mathrm{\Delta }I=3/2`$ operators and turns out to be suppressed in view of the yet not completely understood “$`\mathrm{\Delta }I=1/2`$ rule”. Further considerations about these non-octet-like terms will be given in Sec. V. We consider the effective lagrangian given by
$$_w^{(8)}=g\text{Tr}[\lambda _6_\mu U^\mu U^{}]+g^{}\text{Tr}[\lambda _6_\mu U^\mu U^{}_\nu U^\nu U^{}]+g^{\prime \prime }\text{Tr}[\lambda _6_\mu U_\nu U^{}^\mu U^\nu U^{}].$$
(7)
It should be noticed that this is not the most general octet-like $`\mathrm{\Delta }S=1`$ interaction one can write down up to $`𝒪(p^4)`$ in the momentum power expansion. The latter, containing many other terms, has been presented in Ref.. For the decays we are interested in (no external fields), it turns out that the most general lagrangian includes 15 independent terms of order $`p^4`$ leading to pure $`\mathrm{\Delta }I=1/2`$ transitions. However, to this order, it has been shown that the couplings considered in (7) are sufficient to fit the known data on $`K\pi \pi `$ and $`K\pi \pi \pi `$. In the absence of further information from the meson sector, we will just stick to this simple form. In order to give an idea of the uncertainties in our calculations, we will consider two sets of values for the constants $`g,g^{}`$ and $`g^{\prime \prime }`$ which provide fits to the kaon data of similar quality. Set A corresponds to the parameters used in Refs. ,
$$g=3.60\times 10^8m_\pi ^2;g^{}/g=1.50\times 10^1\mathrm{fm}^2;g^{\prime \prime }/g=6.74\times 10^2\mathrm{fm}^2,$$
(8)
while Set B corresponds to the values obtained in Ref.,
$$g=2.98\times 10^8m_\pi ^2;g^{}/g=1.69\times 10^1\mathrm{fm}^2;g^{\prime \prime }/g=1.87\times 10^2\mathrm{fm}^2.$$
(9)
To calculate the hyperon decays in the context of the Skyrme model with $`SU(3)`$ collective coordinates, we include the soft meson fluctuations on top of the soliton background. This is achieved using
$$U=U_MU_0(𝐫,t)U_M,$$
(10)
where $`U_M=1+i\stackrel{}{\tau }\stackrel{}{\pi }/(2f_\pi )+\mathrm{}`$. Inserting this expression into $`_w^{(8)}`$ and taking the appropriate matrix elements one can obtain the desired $`S`$-wave and $`P`$-wave decay amplitudes. This is worked out in the following two sections.
## III $`S`$-wave amplitudes
As mentioned in the Introduction, it has been recently shown that if the Cabbibo current-current form is used to describe the weak interactions, the present soliton model leads to a reasonably good description of the $`S`$-wave hyperon decay amplitudes. In this section we study these amplitudes using the effective weak chiral lagrangian given by Eqs. (79).
As usual, we assume that isospin symmetry is preserved. In such limit, the following relations between the non-leptonic decay amplitudes can be derived:
$`\mathrm{\Sigma }_{}^{}`$ $`=`$ $`\mathrm{\Sigma }_+^+\sqrt{2}\mathrm{\Sigma }_0^+`$ (12)
$`\mathrm{\Lambda }_{}^0`$ $`=`$ $`\sqrt{2}\mathrm{\Lambda }_0^0`$ (13)
$`\mathrm{\Xi }_{}^{}`$ $`=`$ $`\sqrt{2}\mathrm{\Xi }_0^0,`$ (14)
where the lower indices indicate the charge of the outgoing pion. In this way, only four of the seven measurable amplitudes need to be considered as independent. For simplicity, we choose these amplitudes to be $`\mathrm{\Lambda }_0^0`$, $`\mathrm{\Sigma }_0^+`$, $`\mathrm{\Xi }_0^0`$ and $`\mathrm{\Sigma }_+^+`$.
For a process $`BB^{}\pi `$, we can define the amplitudes $`𝒜`$ and $``$, corresponding to $`S`$\- and $`P`$-wave decays respectively, according to
$$iB^{}\pi |_w^{(8)}|B=\overline{u}_B^{}(𝒜+\gamma _5)u_B.$$
(15)
In the soft-pion limit, the octet nature of $`_w^{(8)}`$, together with current algebra relations, lead to an additional constraint for the $`S`$-wave $`\mathrm{\Sigma }`$ decay amplitudes, namely $`𝒜(\mathrm{\Sigma }_{}^{})=\sqrt{2}𝒜(\mathrm{\Sigma }_0^+)`$. Thus from (12) one obtains $`𝒜(\mathrm{\Sigma }_+^+)=0`$. The remaining $`𝒜`$ amplitudes can be calculated by replacing Eqs. (10) and (6) in $`_w^{(8)}`$ and taking the corresponding matrix elements. We find
$$𝒜(BB^{}\pi ^0)=\alpha B^{}|R_{78}|B,$$
(16)
where
$$\alpha =\frac{4\pi i}{\sqrt{3}f_\pi }𝑑rr^2\left[g\left(F^2+2\frac{\mathrm{sin}^2F}{r^2}\right)g^{}\left(F^2+2\frac{\mathrm{sin}^2F}{r^2}\right)^2g^{\prime \prime }\left(F^44\frac{F^2\mathrm{sin}^2F}{r^2}\right)\right].$$
(17)
In Eq. (16), $`R_{78}`$ stands for an $`SU(3)`$ rotation matrix, $`R_{ab}=1/2\text{Tr}\left[\lambda _aA\lambda _bA^{}\right]`$. As explained in the Appendix, its matrix elements between the collective wavefunctions describing the different baryon states can be calculated as linear combinations of $`SU(3)`$ Clebsch-Gordan coefficients.
Our results for the $`𝒜`$ amplitudes are summarized in Table I, where we also quote the values corresponding to the quark model (QM) . Following the usual convention, the overall phase has been fixed to obtain $`𝒜(\mathrm{\Lambda }_{}^0)`$ real and positive. It can be seen from the table that the predictions obtained with Set A are about 15 % higher than those arising from Set B. In both cases, our results are somewhat below the experimental values. However, since the deviation is in the same direction for all processes (notice that this is not the case for the QM values), the agreement is significantly improved if one considers the ratios between the different amplitudes. In general, it could be said that our results and those corresponding to the QM are of similar quality.
It is also interesting to compare the present results with those of previous soliton calculations. As already mentioned, in some of them the Cabbibo current-current has been used. Within such scheme the best agreement with empirical data has been obtained in Ref. where, as done here, the baryon wavefunctions arise from an exact diagonalization of the SU(3) collective hamiltonian. Generally speaking the results reported in Table I are somewhat smaller (in absolute value) than those of Ref.. However, it should be stressed that the present calculation is free from the uncertainties related to the question of whether (and to which extent) QCD enhancement factors have to be included in soliton calculations. Here, such factors are already accounted for in the value of the low energy constants that appear in the weak lagrangian. On the other hand, if our results are compared with previous calculations based on effective weak chiral lagrangians, we see that the use of empirical input parameters in the strong effective action, together with the exact diagonalization of the SU(3) collective hamiltonian, lead to a significant improvement in the predictions.
A final comment concerns Ref. and its Addendum . In Ref., the author studies weak decays of hyperons in a chiral topological soliton model, starting with a $`\mathrm{\Delta }S=1`$ lagrangian including six $`𝒪(p^4)`$ terms. As shown in the Addendum, the Cayley-Hamilton theorem can be used to reduce these six terms to only four independent ones. Moreover, following the same steps as in our calculation, one arrives to further relations between the corresponding hyperon decay amplitudes. At the end one is effectively left with only two terms, which —as done in the present work— can be chosen to be those proportional to $`g^{}`$ and $`g^{\prime \prime }`$ in Eq. (7). Following Ref., one might try to see whether it is possible to find a set of coefficients for those terms capable to reproduce the empirical values for both kaon and hyperon decays. For this purpose, and in order to relate the coefficients with the hyperon amplitudes, the author of Ref. makes use of the relation (16) and takes the $`SU(3)`$ symmetric limit. In this way all decay amplitudes can be expressed in terms of one of them and the $`f/d`$ ratio. The use of empirical values for the latter leads to an incompatible system of equations, so that the author argues that this precludes a successful application of the chiral lagrangian model. We believe that the results found in the present work provide some ground to disagree with such strong conclusion. In fact, the relations in Ref. are expected to be modified by $`SU(3)`$ breaking effects and by next-to-leading order corrections in $`N_c`$, which in general introduce modifications to Eq. (16). From the results displayed in Table I it is seen that if the chiral soliton approach is used together with an effective weak chiral lagrangian consistent with $`K`$ meson decays, one can obtain a reasonably good description of the $`S`$-wave hyperon decays, already at leading order in $`N_c`$.
## IV $`P`$-wave amplitudes
We turn now to the evaluation of the $`P`$-wave amplitudes. In the non-relativistic limit, they can be calculated using
$$\overline{u}_B^{}\gamma _5u_B\frac{}{2\overline{M}}\chi ^{}(\lambda ^{})\stackrel{}{\sigma }\stackrel{}{q}\chi (\lambda ),$$
(18)
where $`\overline{M}`$ is the average of the $`B`$ and $`B^{}`$ empirical masses. As stated in the Introduction, in the present model $``$ receives contributions of two different kinds. One of them arises from contact terms in $`_w^{(8)}`$ (see Fig. 1a), whereas the other one is given by the pole diagrams shown in Figs. 1b-1d. Our evaluation of the contact contribution in the case of $`\pi ^0`$ emission leads to
$$_{contact}(BB^{}\pi ^0)=\frac{2\overline{M}}{3f_\pi }B^{}|d^3r\widehat{𝒫}_c|B,$$
(19)
where<sup>*</sup><sup>*</sup>*Eq. (24) shows some differences with respect to Eq. (24) of the erratum of Ref. which, we believe, contains errors and/or misprints.
$`\widehat{𝒫}_c`$ $`=`$ $`g\left\{(1+c)(F^{}+2{\displaystyle \frac{s}{r}})R_{63}{\displaystyle \frac{2}{\sqrt{3}}}(1c)(F^{}2{\displaystyle \frac{s}{r}})R_3^{(+)}\right\}`$ (24)
$`+g^{}\left(F^2+2{\displaystyle \frac{s^2}{r^2}}\right)\left\{(1+c)(F^{}+2{\displaystyle \frac{s}{r}})R_{63}+{\displaystyle \frac{2}{\sqrt{3}}}\left[F^{}(3c)+2{\displaystyle \frac{s}{r}}(3c1)\right]R_3^{(+)}\right\}`$
$`+g^{\prime \prime }\{(1+c)F^{}(F^2+2{\displaystyle \frac{s}{r}}F^{}+2{\displaystyle \frac{s^2}{r^2}})R_{63}`$
$`+{\displaystyle \frac{2}{\sqrt{3}}}\left[F^3(3c)+2{\displaystyle \frac{s}{r}}F^{}(1c)(F^{}{\displaystyle \frac{s}{r}})+4{\displaystyle \frac{s^3}{r^3}}c\right]R_3^{(+)}`$
$`{\displaystyle \frac{8}{\sqrt{3}}}{\displaystyle \frac{s}{r}}[cF^2+F^{}{\displaystyle \frac{s}{r}}+c{\displaystyle \frac{s^2}{r^2}}]R_3^{()}\},`$
with
$$R_3^{(\pm )}=R_{68}R_{33}\pm R_{63}R_{38}.$$
(25)
Similar expressions can be found in the case of charged outgoing pions. On the other hand, from the pole diagrams we obtain
$$_{pole}(BB^{}\pi ^0)=2\sqrt{2}\overline{M}\underset{B^{\prime \prime }}{}\left[g_A^{B^{}B^{\prime \prime }}\frac{𝒜(BB^{\prime \prime }\pi ^0)}{M_BM_{B^{\prime \prime }}}+g_A^{BB^{}}\frac{𝒜(B^{\prime \prime }B^{}\pi ^0)}{M_B^{}M_{B^{\prime \prime }}}\right],$$
(26)
where we have neglected the small $`K`$ pole term contribution in Fig. 1d, and we have made use of the generalized Goldberger-Treiman relations to write the strong coupling constants in terms of the axial charges. Notice that Eq. (26) includes a sum over a set of intermediate states $`B^{\prime \prime }`$. In our calculations, we have included the $`J^\pi =\frac{1}{2}^+`$ collective eigenfunctions that arise from the exact diagonalization of the strong Hamiltonian. It turns out that only a few excited states are needed to find convergence and their contribution represents, at most, $`15\%`$ of the total values of $`_{pole}`$.
For the sake of consistency, in order to estimate the size of the $`_{pole}`$ amplitudes we will take into account both the axial charges and the $`𝒜`$ amplitudes obtained within our model. Therefore, we consider the axial charge operator $`\widehat{g}_A`$ arising from the action (1), which reads
$`{\displaystyle \frac{\widehat{g}_A}{\sqrt{2}}}`$ $`=`$ $`{\displaystyle \frac{1}{3}}{\displaystyle d^3r\left\{f_\pi ^2\left(F^{}+2c\frac{s}{r}\right)+\frac{2}{e^2}\frac{s}{r}\left(c\frac{s^2}{r^2}+\frac{s}{r}F^{}+cF^2\right)\right\}R_{33}}`$ (29)
$`+{\displaystyle \frac{f_K^2f_\pi ^2}{9}}{\displaystyle d^3r(1c)\left\{\left[\frac{2s}{r}(1+2c)+F^{}\right](1R_{88})R_{33}+\left[\frac{2s}{r}F^{}\right]R_{83}R_{38}\right\}}`$
$`+{\displaystyle \frac{N_c}{36\pi ^2\mathrm{\Theta }_K}}{\displaystyle d^3r(1c)\frac{s}{r}\left(\frac{s}{r}2F^{}\right)d_{3kk^{}}R_{3k}J_k^{}},`$
where $`\mathrm{\Theta }_K`$ is the kaonic moment of inertia. It is worth to mention that this operator leads to a low value for the neutron beta decay form factor, $`g_A0.75`$, compared with the experimental result of about 1.25.
Numerical results for both the contact and pole contributions to the $``$ amplitudes are given in Table II. It can be seen that the absolute values for the total amplitudes are far too small in comparison with the experimental results. In the case of the pole contributions, this could be explained in part by an underestimation of the axial form factors, as suggested by the low value in the case of the neutron beta decay. In particular, for $`(\mathrm{\Lambda }_0^0)`$, the contribution obtained from Eq. (26) using the empirical values of the $`𝒜`$ amplitudes and axial charges is in very good agreement with the experimental result (see the value corresponding to the chiral fit in Table II). In our calculation, instead, the suppression arising from the somewhat low predictions for the $`S`$-wave amplitudes, together with the underestimation of the axial form factors, conspire to end up with a reduction factor of about 1/3. In the case of the remaining $`P`$-wave amplitudes, it is well known that the usage of phenomenological values in (26) does not allow to get a good fit of the experimental values. In this sense, the contact contributions have been suggested as a possible novel ingredient to solve the discrepancy. Our results show, however, that the values for $`_{contact}`$ amount at most 1/10 of the empirical $``$ amplitudes. Thus, even if the effect goes in the right direction, the contact contributions appear to be too small to represent a potential solution for the $`S/P`$ problem.
To check the dependence of our results on the Skyrme parameter $`ϵ`$ we have considered departures from the central value $`ϵ=4.1`$. We find that the absolute values of the amplitudes tend to increase as $`ϵ`$ increases. The amplitudes which turn out to be the most sensitive to the variation of $`ϵ`$ are those corresponding to the $`\mathrm{\Lambda }n\pi ^0`$ process. For $`ϵ=4.5`$, which already implies an increase of more that $`25\%`$ for the $`\mathrm{\Delta }N`$ splitting, we find, for Set A, $`𝒜(\mathrm{\Lambda }_0^0)2\times 10^7`$ which is quite close to the empirical value $`𝒜(\mathrm{\Lambda }_0^0)_{emp}=2.37\times 10^7`$. However, even in this case, the predicted $`(\mathrm{\Lambda }_0^0)`$ is still more than a factor 2 below the corresponding empirical value. Thus, we can conclude that the statements above are quite robust under variations of the only adjustable parameter in our calculation.
## V Contributions of non-octet-like components
As stated in Sec. II, the non-leptonic decay amplitudes are dominated by the octet-like components of the weak effective lagrangian, hence only these components have been considered in the previous two sections. On the other hand, we have also mentioned that this approximation leads to a vanishing $`\mathrm{\Sigma }_+^+`$ $`S`$-wave amplitude. In fact, the experimental value of $`𝒜(\mathrm{\Sigma }_+^+)`$, although significantly smaller than the other measured $`S`$-wave amplitudes, is found to be different from zero. In this section we will investigate whether within the present model the standard non-octet-like contributions to the weak effective lagrangian are able to account for this difference. It is clear that such contributions will also modify the results obtained in the previous sections for the other $`S`$-wave and $`P`$-wave decay amplitudes. However, the modifications are, in the worst case, of the same order of magnitude than the uncertainties involved in the parameters $`g`$, $`g^{}`$ and $`g^{\prime \prime }`$ in the weak effective lagrangian. Therefore, in what follows we will concentrate only on the $`S`$-wave $`\mathrm{\Sigma }_+^+`$ decay amplitude.
The lowest order 27-plet contribution to the weak effective lagrangian occurs at $`p^2`$. It can be written as
$$_2^{(27)}=c_3t_{cd}^{ab}\text{Tr}\left(Q_a^cU^{}_\mu U\right)\text{Tr}\left(Q_d^bU^{}^\mu U\right)$$
(30)
where $`\left(Q_a^c\right)_{ij}=\delta _{cj}\delta _{ai}`$, and
$`t_{21}^{31}`$ $`=`$ $`t_{12}^{13}=t_{13}^{12}=t_{31}^{21}={\displaystyle \frac{3}{2}},`$ (31)
$`t_{12}^{31}`$ $`=`$ $`t_{21}^{13}=t_{13}^{21}=t_{31}^{12}=1,`$ (32)
with $`t_{cd}^{ab}=0`$ otherwise. As in the case of the octet-like piece, we will take into account also the effect of next-to-leading order couplings. We consider the $`𝒪(p^4)`$ interaction
$`_4^{(27)}`$ $`=`$ $`g_1t_{cd}^{ab}\text{Tr}\left(Q_a^cL_\mu \right)\text{Tr}\left(Q_d^bL_\nu L^\mu L^\nu \right)`$ (35)
$`+g_2t_{cd}^{ab}\text{Tr}\left(Q_a^cL_\mu \right)\text{Tr}\left(Q_d^b\{L^\mu ,L^2\}\right)`$
$`+g_3t_{cd}^{ab}\text{Tr}\left(Q_a^c[L_\mu ,L_\nu ]\right)\text{Tr}\left(Q_d^b[L^\mu ,L^\nu ]\right).`$
Once again the most general lagrangian allowed by chiral symmetry includes many possible terms, and the coupling constants cannot be fully determined from the available information on the kaon sector. The structure chosen in (35) is, however, sufficient to get a good fit of $`2\pi `$ and $`3\pi `$ $`K`$ decays. From such a fit one obtains $`g_11.0\times 10^{10}`$, $`g_20.4\times 10^{10}`$, $`g_30.1\times 10^{10}`$ together with $`c_3/f_\pi ^20.8\times 10^9`$. This set of values is used in the numerical calculation below.
The desired $`S`$-wave $`\mathrm{\Sigma }_+^+`$ amplitude can be now easily obtained by inserting the explicit form of the chiral field $`U`$ in the effective couplings (30) and (35). By doing this we arrive to
$$𝒜^{(27)}(\mathrm{\Sigma }^+n\pi ^+)=\frac{\sqrt{15}}{8\sqrt{2}}\left(I_2+I_4\right)n|D_{\frac{3}{2},0}^{27}|\mathrm{\Sigma }^+.$$
(36)
Here, the left lower index of the SU(3) Wigner function $`D`$ stands for $`(Y,I,I_3)=(1,3/2,3/2)`$ while the right lower index corresponds to $`(0,0,0)`$. The radial integrals $`I_2`$ and $`I_4`$ are given by
$`I_2`$ $`=`$ $`{\displaystyle \frac{16\pi }{3f_\pi }}c_3{\displaystyle 𝑑rr^2\left(F^2+2\frac{\mathrm{sin}^2F}{r^2}\right)}`$ (37)
$`I_4`$ $`=`$ $`{\displaystyle \frac{16\pi }{3f_\pi }}{\displaystyle }drr^2[g_1F^2(F^24{\displaystyle \frac{\mathrm{sin}^2F}{r^2}})+2g_2(F^2+2{\displaystyle \frac{\mathrm{sin}^2F}{r^2}})^2`$ (39)
$`+g_3{\displaystyle \frac{8\mathrm{sin}^2F}{r^2}}(2F^2+{\displaystyle \frac{\mathrm{sin}^2F}{r^2}})].`$
Even if the integrand of $`I_4`$ is suppressed by the coefficients $`g_i`$ (which are one to two orders of magnitude lower that $`c_3/f_\pi ^2`$), it can be seen that the suppression is compensated by the values of the radial integrals, in such a way that at the end $`I_4`$ dominates over $`I_2`$. Evaluating the matrix element in (36), we finally obtain
$$𝒜(\mathrm{\Sigma }_+^+)0.01\times 10^7,$$
(40)
to be compared with the empirical value $`𝒜_{emp}(\mathrm{\Sigma }_+^+)=0.13\times 10^7`$ given in Table I. We observe that our estimation for $`𝒜(\mathrm{\Sigma }_+^+)`$, although non-vanishing, is roughly one order of magnitude smaller than the empirical result. As in the case of the octet-like contributions this statement remains valid for reasonable variations of the Skyrme parameter around its central value $`ϵ=4.1`$.
## VI Conclusions
In this work we have revisited the problem of the calculation of the non-leptonic hyperon decay amplitudes in the topological soliton models. We have used the approach to the $`SU(3)`$ Skyrme model in which both the isospin and the strange degrees of freedom are treated as collective rotations around the usual $`SU(2)`$ hedgehog ansatz and the symmetry breaking terms in the strong action are diagonalized exactly. To describe the weak interactions we have used a chiral effective action, in which low energy constants are adjusted to describe the known $`2\pi `$ and $`3\pi `$ weak kaon decays. For the $`S`$-wave decay amplitudes we have found that, compared with previous calculations based on effective weak chiral lagrangians, the use of empirical input parameters in the strong effective action, together with the exact diagonalization of the $`SU(3)`$ collective hamiltonian, lead to a significant improvement in the predictions. A similar result has been recently obtained using a Cabbibo current-current type weak interaction. Although our predictions are about $`30\%`$ below the empirical values we consider them as satisfactory in view of the simplicity of the model and the fact that higher order $`N_c`$ corrections of that size are to be expected. On the other hand, our results badly fail to reproduce the empirical $`P`$-wave amplitudes. In soliton models, such amplitudes receive two types of contributions, namely those arising from the usual pole diagrams and those coming from contact terms. The presence of the latter provided some hope that the long standing $`S/P`$ wave puzzle could find a solution within these models. Our results show that, unfortunately, such contact contributions are far too small to close the gap between the predictions coming from the pole terms alone and the empirical values. Although one cannot exclude some corrections to these results due to higher order effects neglected in this work (such as e.g. the kaon induced components which are known to play a significant role in the determination of the parity violating $`\pi N`$ coupling constant), it is difficult to believe that they could lead to a solution of this problem. Finally, we have estimated the contribution to the decay amplitudes coming from non-octet terms in the weak effective action. Since these contributions are generally very small we have concentrated only on the $`S`$-wave $`\mathrm{\Sigma }_+^+`$ decay amplitude which, as well known, vanishes if only octet terms are considered. Our result, although non-zero, turns out to be roughly one order of magnitude smaller than the empirical value. This clearly indicates that, within the Skyrme model, more refined wave functions and/or effective weak interactions are needed to understand the subtle effects related with the small violations of the $`\mathrm{\Delta }I=1/2`$ rule observed in the non-leptonic hyperon $`S`$-wave decays.
###### Acknowledgements.
D.G.D. acknowledges a Reentry Grant and a research fellowship from Fundación Antorchas, Argentina. This work was supported in part by the grant PICT 03-00000-00133 from ANPCYT, Argentina. N.N.S. is fellow of CONICET, Argentina.
## Calculation of the collective matrix elements
As explained in the main text the calculation of the decay amplitudes involves the matrix elements of some collective operators $`\widehat{𝒪}`$ between baryon wavefunctions. To evaluate these matrix elements we proceed as follows.
In general, the wavefunction corresponding to a baryon $`B`$ can be expanded in terms of $`SU(3)`$ Wigner functions $`D_{\alpha \beta }^R`$,
$$\mathrm{\Psi }_B=\underset{R}{}C_B^R\sqrt{\text{dim }(R)}D_{\alpha \beta }^R,$$
(41)
where $`\alpha =(Y,I,I_3)`$ and $`\beta =(1,J,J_3)`$ carry the baryon quantum numbers, and $`R`$ is the corresponding representation. The coefficients $`C_B^R`$ are obtained from the diagonalization procedure described in Sec. II. In addition, the collective operators $`\widehat{𝒪}`$ can always be expressed as
$$\widehat{𝒪}=\underset{\widehat{\alpha },\widehat{\beta },\widehat{R}}{}\gamma _{\widehat{\alpha },\widehat{\beta }}^{\widehat{R}}D_{\widehat{\alpha },\widehat{\beta }}^{\widehat{R}},$$
(42)
where $`\gamma _{\widehat{\alpha },\widehat{\beta }}^{\widehat{R}}`$ are numerical coefficients. These coefficients result from expressing the “cartesian” $`SU(3)`$ indexes $`a=1\mathrm{}8`$ in terms of the “spherical” $`SU(3)`$ indexes $`(Y,I,I_3)`$ and performing suitable Clebsch-Gordan series expansions when needed. For example, for the operator $`R_{78}`$ appearing in Eq. (16) we have
$$R_{78}=\frac{i}{\sqrt{2}}\left[D_{(1,\frac{1}{2},\frac{1}{2}),(0,0,0)}^8D_{(1,\frac{1}{2},\frac{1}{2}),(0,0,0)}^8\right],$$
(43)
while the combination $`R_3^{(+)}`$ in Eq. (25) can be written as
$`R_3^{(+)}`$ $`=`$ $`{\displaystyle \frac{\sqrt{6}}{10}}\left[D_{(1,\frac{1}{2},\frac{1}{2}),(0,1,0)}^8+D_{(1,\frac{1}{2},\frac{1}{2}),(0,1,0)}^8\right]+{\displaystyle \frac{1}{10}}\left[D_{(1,\frac{1}{2},\frac{1}{2}),(0,1,0)}^{27}+D_{(1,\frac{1}{2},\frac{1}{2}),(0,1,0)}^{27}\right]`$ (45)
$`{\displaystyle \frac{1}{\sqrt{5}}}\left[D_{(1,\frac{3}{2},\frac{1}{2}),(0,1,0)}^{27}+D_{(1,\frac{3}{2},\frac{1}{2}),(0,1,0)}^{27}\right]`$
Thus, the matrix element $`B^{}|\widehat{𝒪}|B`$ can finally be expressed in terms of the standard $`SU(3)`$ Clebsch-Gordan coefficients. Namely,
$$B^{}|\widehat{𝒪}|B=\underset{\widehat{\alpha },\widehat{\beta },\widehat{R}}{}\gamma _{\widehat{\alpha },\widehat{\beta }}^{\widehat{R}}\underset{R,R^{},\mu }{}C_B^R(C_B^{}^R^{})^{}\sqrt{\frac{\text{dim }(R)}{\text{dim }(R^{})}}\left(\begin{array}{ccc}\widehat{R}& R& R^{}\\ \widehat{\alpha }& \alpha & \alpha ^{}\end{array}\right)\left(\begin{array}{ccc}\widehat{R}& R& R^{}\mu \\ \widehat{\beta }& \beta & \beta ^{}\end{array}\right),$$
(46)
where the brackets indicate the Clebsch-Gordan coefficients. The sum over $`\mu `$ refers to the situations in which the Clebsch-Gordan expansion of the product of two $`D`$’s includes more than one representation with the same dimension. |
warning/0001/hep-ph0001128.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In this contribution I discuss the possibility of production of high energy cosmic neutrinos ($`E10^6`$ GeV) in cores of Active Galactic Nuclei (AGN) originating from proton acceleration, the effects of three flavour neutrino mixing on these high energy cosmic neutrino fluxes and the prospects for their observations in new km<sup>2</sup> surface area underwater/ice neutrino telescopes.
In addition to AGNs, high energy cosmic neutrinos may also be produced in several currently envisaged other cosmologically distant astrophysical sources. These sources may include, for instance, Gamma Ray Burst fireballs and Topological Defects . For some possible effects of neutrino mixing other than the flavour one on high energy cosmic neutrino fluxes, see . The present study is particularly useful as several high energy neutrino telescopes are now at their rather advanced stage of development and deployment .
I start in Section 2 with a brief description of possibility of production of high energy cosmic neutrinos and discuss in some detail the effects of three flavour neutrino mixing on their subsequent propagation and further discuss the prospects for their detection. In Section 3, I summarize the results.
## 2 High Energy Neutrinos from AGNs
### 2.1 Production
High energy cosmic neutrinos may mainly be produced either in $`p\gamma `$ or in $`pp`$ collisions in a cosmologically distant environment.
In $`p\gamma `$ collisions, high energy $`\nu _e`$ and $`\nu _\mu `$ are mainly produced through $`p+\gamma \mathrm{\Delta }^+n+\pi ^+`$ (typically with $`\nu _e/\nu _\mu 1/2`$). The same collisions will give rise to a greatly suppressed high energy $`\nu _\tau `$ flux ($`\nu _\tau /\nu _{e,\mu }<\mathrm{\hspace{0.17em}10}^5`$) mainly through $`p+\gamma D_S^++\mathrm{\Lambda }^0+\overline{D}^0`$. In $`pp`$ collisions, the $`\nu _\tau `$ flux may be obtained through $`p+pD_S^++X`$. The relatively small cross-section for $`D_S^+`$ production together with the low branching ratio into $`\nu _\tau `$ implies that the $`\nu _\tau `$ flux in $`pp`$ collisions is also suppressed up to 5 orders of magnitude relative to $`\nu _e`$ and/or $`\nu _\mu `$ fluxes (which are mainly produced through $`\pi ^\pm `$) .
### 2.2 Propagation
Matter effects on vacuum neutrino oscillations are relevant if $`G_F\rho /m_N\delta m^2/2E`$. Using $`\rho `$ from Ref. as an example, it turns out that matter effects are absent for $`\delta m^2𝒪(10^{10})`$ eV<sup>2</sup>. Matter effects are not expected to be important in the neutrino production regions around AGN and will not be further discuss here.
In the framework of three flavour analysis, the flavour precession probability from $`\alpha `$ to $`\beta `$ neutrino flavour is
$`P(\nu _\alpha \nu _\beta )P_{\alpha \beta }`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{3}{}}}|U_{\alpha i}|^2|U_{\beta i}|^2`$ (1)
$`+{\displaystyle \underset{ij}{}}U_{\alpha i}U_{\beta i}^{}U_{\alpha j}^{}U_{\beta j}\mathrm{cos}\left({\displaystyle \frac{2\pi L}{l_{ij}}}\right),`$
where $`\alpha ,\beta =e,\mu ,`$ or $`\tau `$. $`U`$ is the 3$`\times `$3 MNS mixing matrix and can be obtained in usual notation through
$`U`$ $``$ $`R_{23}(\theta _1)\text{diag}(e^{i\delta /2},1,e^{i\delta /2})`$ (2)
$`R_{31}(\theta _2)\text{diag}(e^{i\delta /2},1,e^{i\delta /2})R_{12}(\theta _3),`$
thus coinciding with the standard form given by the Particle Data Group . In Eq. (1), $`l_{ij}4\pi E/\delta m_{ij}^2`$ with $`\delta m_{ij}^2|m_i^2m_j^2|`$ and $`L`$ is the distance between the source and the detector. For simplicity, I assume here a vanishing value for CP violating phase $`\delta `$ and $`\theta _{31}`$ in $`U`$.
At present, the atmospheric muon and solar electron neutrino deficits can be explained with oscillations among three active neutrinos . For this, typically, $`\delta m^2𝒪(10^3)`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta 𝒪(1)`$ for the explanation of atmospheric muon neutrino deficit, whereas for the explanation of solar electron neutrino deficit, we may have $`\delta m^2𝒪(10^{10})`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta 𝒪(1)`$ \[just so\] or $`\delta m^2𝒪(10^5)`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta 𝒪(10^2)`$ \[SMA (MSW)\] or $`\delta m^2𝒪(10^5)`$ eV<sup>2</sup> and $`\mathrm{sin}^22\theta 𝒪(1)`$ \[LMA (MSW)\]. The present status of data thus permits multiple oscillation solutions to solar neutrino deficit. I intend to discuss here implications of these mixings for high energy cosmic neutrino propagation.
In the above explanations, the total range of $`\delta m^2`$ is $`10^{10}\delta m^2/`$ eV$`{}_{}{}^{2}10^3`$ irrespective of neutrino flavour. The typical energy span relevant for possible flavour identification for high energy cosmic neutrinos is $`210^6E/`$GeV$`210^7`$ in which currently the neutrino flux from cores of AGNs dominate. Taking a typical distance between the AGN and our galaxy as $`L100`$ Mpc (where 1 pc $`310^{16}`$ m), note that $`\mathrm{cos}`$ term in Eq. (1) vanishes and so Eq. (1) reduces to
$$P_{\alpha \beta }\underset{i=1}{\overset{3}{}}|U_{\alpha i}|^2|U_{\beta i}|^2.$$
(3)
It is assumed here that no relatively dense objects exist between the AGN and the earth so as to effect significantly this oscillations pattern. Note also that since $`P_{\alpha \beta }`$ in above Eq. is symmetric under the exchange of indices $`\alpha `$ and $`\beta `$ implying that no $`T`$ (or $`CP`$) violation effects arise in neutrino vacuum flavour oscillations for high energy cosmic neutrinos .
Let me denote by $`F_\alpha ^0`$, the intrinsic neutrino fluxes. From the discussion in the previous Subsection, it follows that $`F_e^0:F_\mu ^0:F_\tau ^0=1:2:<10^5`$. For simplicity, I take these ratios as 1 : 2 : 0. In order to estimate the final (downward going) flux ratios of high energy cosmic neutrinos reaching on earth, let me introduce a 3$`\times `$3 matrix of vacuum flavour precession probabilities such that
$$F_\alpha =\underset{\beta }{}P_{\alpha \beta }F_\beta ^0,$$
(4)
where the unitarity conditions for $`P_{\alpha \beta }`$ read as
$`P_{ee}+P_{e\mu }+P_{e\tau }`$ $`=`$ $`1,`$
$`P_{e\mu }+P_{\mu \mu }+P_{\mu \tau }`$ $`=`$ $`1,`$
$`P_{e\tau }+P_{\mu \tau }+P_{\tau \tau }`$ $`=`$ $`1.`$ (5)
The explicit form for the matrix $`P`$ in case of just so flavour oscillations as solution to solar neutrino problem along with the solution to atmospheric neutrino deficit in terms of $`\nu _\mu `$ to $`\nu _\tau `$ oscillations with maximal mixing is
$$P=\left(\begin{array}{ccc}1/2& 1/4& 1/4\\ 1/4& 3/8& 3/8\\ 1/4& 3/8& 3/8\end{array}\right).$$
(6)
Using Eq. (6) and Eq. (4), it follows that $`F_e:F_\mu :F_\tau =1:1:1`$ at the level of $`F_e^0`$. Also, Eq. (5) is satisfied. The same flux ratio is obtained in the remaining two cases for which the corresponding $`P`$ matrics are: \[for SMA (MSW)\]
$$P=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1/2& 1/2\\ 0& 1/2& 1/2\end{array}\right),$$
(7)
whereas in case of LMA (MSW),
$$P=\left(\begin{array}{ccc}5/8& 3/16& 3/16\\ 3/16& 13/32& 13/32\\ 3/16& 13/32& 13/32\end{array}\right).$$
(8)
Thus, essentially independent of the oscillation solutions for solar neutrino problem, it follows that $`F_e:F_\mu :F_\tau =1:1:1`$. The deviations from these ratios are estimated to be small .
Summarizing, although intrinsically the downward going high energy cosmic tau neutrino flux is negligibally small however because of vacuum flavour oscillations it becomes comparable to $`\nu _e`$ flux thus providing some prospects for its possible detection.
### 2.3 Prospects for detection
I briefly mention here the prospects for detection of downward going high energy cosmic tau neutrinos through double shower technique . For prospects of observations of high energy cosmic tau neutrinos other than double shower technique, see , whereas for possibility of detection of non tau neutrinos, see .
The downward going tau neutrinos reaching close to the surface of the detector may undergo a charged current deep inelastic scattering with nuclei inside/near the detector and produce a tau lepton in addition to a hadronic shower.
This tau lepton traverses a distance, on average proportional to its energy, before it decays back into a tau neutrino and a second shower most often induced by decay hadrons. The second shower is expected to carry about twice as much energy as the first and such double shower signals are commonly referred to as a double bangs. As tau leptons are not expected to have further relevant interactions (with high energy loss) in their decay timescale, the two showers should be separated by a clean $`\mu `$-like track.
The calculation of downward going contained but separable double shower event rate can be carried out by replacing the muon range expression with the tau range expression and then subtracting it from the linear size of a typical high energy neutrino telescope in the event rate formula while using the expected $`\nu _\tau `$ flux spectrum given by Eq. (4). This ensures that the two separate showers are contained within km of the underwater/ice detector. Here, I restrict myself by mentioning that the expected number of contained but separable double showers induced by downward going high energy tau neutrinos for $`E210^6`$ GeV may be $`𝒪(10)`$/yr$``$sr irrespective of the oscillation solutions of solar neutrino problem, if one uses the $`F_e^0`$ from Ref. as an example. At this energy, the two showers initiated by the downward going high energy cosmic tau neutrinos are well separated ($``$ 70 m) such that the size of the second shower is essentially 2 times the first shower and the two showers are connected by a $`\mu `$like track (see Fig. 1). This identification, if empirically realized, may provide a possibility to isolate $`\nu _\tau `$ flavour from the rest of neutrino flavors. The chance of having double shower events induced by non tau neutrinos is negligibly small for relevant energies.
## 3 Conclusions
1. Intrinsically, the flux of high energy cosmic tau neutrinos is quite small, relative to non tau flavour neutrinos, typically being $`F_\tau ^0/F_{e,\mu }^0<\mathrm{\hspace{0.17em}10}^5`$ (whereas $`F_e^0/F_\mu ^01/2`$) from cosmologically distant astrophysical sources, namely, for instance, cores of Active Galactic Nuclei.
2. Because of neutrino oscillations, this ratio can be greatly enhanced. In the context of three flavour neutrino mixing scheme which can accommodate the oscillation solutions to solar and atmospheric neutrino deficits in terms of oscillations between three active neutrinos, the final ratio of fluxes of downward going high energy cosmic neutrinos on earth is $`F_eF_\mu F_\tau F_e^0`$, essentially irrespective of the oscillation solutions to solar neutrino problem.
3. This enhancement in high energy cosmic tau neutrino flux may lead to the possibility of its detection in km<sup>2</sup> surface area high energy neutrino telescopes. For $`210^6E`$/GeV $`210^7`$, the downward going high energy cosmic tau neutrinos may produce a double shower signature because of charged current deep inelastic scattering followed by a subsequent hadronic decay of the associated tau lepton.
## Acknowledgments
I thank Japan Society for the Promotion of Science (JSPS) for financial support. |
warning/0001/hep-ph0001254.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Scalar quarks (the SUSY partners of quarks) could be produced at appreciable rates at the TESLA collider. A large amount of work is already avaliable concerning the phenomenology and the radiative corrections of the scalar fermion sector relevant for the linear collider (see e.g. and references therein). With the experimental precision expected at TESLA it will be become necessary to include also the quantum corrections in the theoretical investigations. The largest radiative corrections for the squark sector of the MSSM are associated with the strong interaction. Such QCD corrections have been investigated for the production cross-section and for several other squark observables . However, the electroweak (EW) corrections can also be sizeable and are in general not negligible. This applies in particular to the squarks of the third generation, owing to the large Yukawa couplings of the fermions
$$\lambda _t\frac{h_t}{g}=\frac{m_t}{\sqrt{2}M_W\mathrm{sin}\beta },\lambda _b\frac{h_b}{g}=\frac{m_b}{\sqrt{2}M_W\mathrm{cos}\beta },$$
(1)
where $`\mathrm{tan}\beta `$ is the ratio between the vacuum expectation values of the two Higgs bosons doublets, $`\mathrm{tan}\beta =v_2/v_1`$. The Yukawa couplings (1) determine the strength of the interactions between quarks and Higgs particles, squarks and higgsinos (the SUSY partners of the quarks and Higgs particles respectively), and part of the interactions between the squarks and Higgs bosons. However, in the case of squarks there are additional interactions originating from the breaking of SUSY, the so-called soft-SUSY-breaking trilinear terms. Although these terms are bounded by the condition that the vacua do not break charge and color, they can be large enough to provide large quantum corrections. Recently, the Yukawa corrections to the production cross-section have received attention.
We have computed the leading electroweak corrections to the partial decay width of a bottom squark into a top quark and a chargino $`\mathrm{\Gamma }(\stackrel{~}{b}t\chi ^{})`$ . The reasons to choose this concrete channel are at least twofold. First of all it is a process of special interest. The third generation squarks can be the lightest sfermions of the model, owing to the large Yukawa couplings that make their mass decrease when one assumes a common sfermion mass scale at a unification scale, according to the Renormalization Group evolution. Secondly, the third generation squarks are most likely to develop large EW corrections, owing to their large couplings to the Higgs sector. Of course, the neutral channels ($`\stackrel{~}{t}t\chi ^0`$ and $`\stackrel{~}{b}b\chi ^0`$) are equally interesting.
A key point in the computation of observables with $`R`$-odd external particles is that it is no longer possible to separate between SUSY and non-SUSY corrections. This means that, when making the appropriate renormalization, both the $`R`$-even and $`R`$-odd sectors of the theory have to be renormalized. Hence, for the computation of the decay process mentioned above we should perform the renormalization of all the neutralino-chargino sector, together with the gauge and the Higgs sector, which is a rather voluminous task. We have therefore chosen a scenario that allows a simplified treatment: If we assume that the gaugino soft-SUSY-breaking masses are much larger than the higgsino mass parameter $`|\mu |`$ then we can treat the $`\chi ^{}`$ appearing in the process as a purely higgsino particle; in this way we can avoid to deal with all the plethora of gauge and gaugino particles in the electroweak sector. Of course this is only a first approximation, but it is already sufficient to demonstrate the importance of the corrections, before performing the full computation. This specification to the leading Yukawa terms is meaningful only in the higgsino approximation, where the soft-SUSY-breaking gaugino masses have to fulfill the relation
$$M^{},M\{|\mu |,M_W\},$$
with the lightest chargino as pure higgsino, that is
$$\chi _1^{}=\stackrel{~}{h}^{},m_{\chi _1^{}}|\mu |,\chi _2^{}=\stackrel{~}{w}^{},m_{\chi _2^{}}M.$$
(2)
## 2 Tree-level relations
The tree-level Lagrangian for the top-sbottom-chargino interactions reads
$$_{t\stackrel{~}{b}\chi }=g\stackrel{~}{b}_a^{}\overline{\chi }_i^+(A_+^{ai}P_L+ϵ_iA_{}^{ai}P_R)t+\mathrm{h}.\mathrm{c}.,$$
(3)
where $`ϵ_i`$ is the sign of the ith chargino eigenvalue $`M_i(i=1,2`$ with $`|M_1|<|M_2|`$) in the real matrix representation, and the coupling matrices are denoted by<sup>2</sup><sup>2</sup>2See Refs. for full notation niceties.
$$A_+^{ai}=R_{1a}V_{i1}\lambda _bR_{2a}V_{i2},A_{}^{ai}=R_{1a}\lambda _tU_{i2}.$$
(4)
The explicit appearance of the Yukawa couplings (1) in the Lagrangian above requires both the introduction of top and bottom quark mass counterterms (in the on-shell scheme) and also a suitable prescription for the renormalization of $`\mathrm{tan}\beta `$. We denote by $`m_{\stackrel{~}{b}_a}(a=1,2)`$, with $`m_{\stackrel{~}{b}_1}<m_{\stackrel{~}{b}_2}`$, the two sbottom mass eigenvalues. The sbottom mixing angle $`\theta _{\stackrel{~}{b}}`$ is defined by the transformation relating the weak-interaction ($`\stackrel{~}{b}_a^{}=\stackrel{~}{b}_L,\stackrel{~}{b}_R`$) and the mass eigenstate ($`\stackrel{~}{b}_a=\stackrel{~}{b}_1,\stackrel{~}{b}_2`$) squark bases:
$$\stackrel{~}{b}_a^{}=R_{ab}\stackrel{~}{b}_b;R=\left(\begin{array}{cc}\mathrm{cos}\theta _{\stackrel{~}{b}}& \mathrm{sin}\theta _{\stackrel{~}{b}}\\ \mathrm{sin}\theta _{\stackrel{~}{b}}& \mathrm{cos}\theta _{\stackrel{~}{b}}\end{array}\right);$$
(5)
$`R`$ is the matrix appearing in eq. (4). By this basis transformation, the sbottom mass matrix,
$$_{\stackrel{~}{b}}^2=\left(\begin{array}{cc}M_{\stackrel{~}{b}_L}^2+m_b^2+\mathrm{cos}2\beta (\frac{1}{2}+\frac{1}{3}s_W^2)M_Z^2& m_b(A_b\mu \mathrm{tan}\beta )\\ m_b(A_b\mu \mathrm{tan}\beta )& M_{\stackrel{~}{b}_R}^2+m_b^2\frac{1}{3}\mathrm{cos}2\beta s_W^2M_Z^2,\end{array}\right),$$
(6)
becomes diagonal: $`R^{}_{\stackrel{~}{b}}^2R=\mathrm{diag}\{m_{\stackrel{~}{b}_2}^2,m_{\stackrel{~}{b}_1}^2\}.`$
Our aim is to compute the radiative corrections in an on-shell renormalization scheme; hence, the input parameters are physical observables (i.e. the physical masses $`m_{\stackrel{~}{b}_2},m_{\stackrel{~}{b}_1}`$, …) rather than formal parameters in the Lagrangian (i.e. the soft-SUSY-breaking parameters $`M_{\stackrel{~}{b}_L}^2,A_b`$, … in eq. (6)). Specifically, we use the following set of independent parameters for the squark sector:
$$(m_{\stackrel{~}{b}_1},m_{\stackrel{~}{b}_2},\theta _{\stackrel{~}{b}},m_{\stackrel{~}{t}_1},\theta _{\stackrel{~}{t}}).$$
(7)
The value of the other stop mass $`m_{\stackrel{~}{t}_2}`$ is then determined by $`SU(2)_L`$ gauge invariance. For the numerical study, we shall use a range of bottom-squark masses $`300350\text{ GeV}`$, relevant for a $`\sqrt{s}=800\text{ GeV}`$ $`e^+e^{}`$ linear collider. The sbottom and stop trilinear soft-SUSY-breaking terms $`A_b`$ and $`A_t`$ are fixed at the tree-level by the previous parameters as follows:
$$A_b=\mu \mathrm{tan}\beta +\frac{m_{\stackrel{~}{b}_2}^2m_{\stackrel{~}{b}_1}^2}{2m_b}\mathrm{sin}2\theta _{\stackrel{~}{b}};A_t=\mu \mathrm{cot}\beta +\frac{m_{\stackrel{~}{t}_2}^2m_{\stackrel{~}{t}_1}^2}{2m_t}\mathrm{sin}2\theta _{\stackrel{~}{t}}.$$
(8)
We impose the approximate (necessary) condition
$$A_q^2<3(m_{\stackrel{~}{t}}^2+m_{\stackrel{~}{b}}^2+M_H^2+\mu ^2),$$
(9)
where $`m_{\stackrel{~}{q}}`$ is of the order of the average squark masses for $`\stackrel{~}{q}=\stackrel{~}{t},\stackrel{~}{b}`$, to avoid colour-breaking minima in the MSSM Higgs potential . Of course the relation (8) receives one-loop corrections. However, since these parameters do not enter the tree-level expressions, these effects translate into two-loop corrections to the process under study. The bound (9) translates into a stringent constrain to the sbottom-quark mixing angle for moderate and large values of $`\mathrm{tan}\beta \text{ }\stackrel{>}{}\text{ }10`$: with an approximate limit $`|\mu |\text{ }\stackrel{>}{}\text{ }80\text{ GeV}`$ from the negative output of the chargino search at LEP, the condition (9) can only be satisfied by a cancellation of the two terms in (8) which is easily spoiled when $`\theta _{\stackrel{~}{b}}`$ is varied. The right hand side of eq. (9) is not rigorous; so we will present results also when this bound is not satisfied, but we will clearly mark these regions. With the use of the bound (9) also the squark-squark-Higgs-boson couplings are restricted. This is a welcome feature, since these couplings can in general be very large, eventually spoiling perturbativity.
It is clear that the radiative corrections to the process $`\stackrel{~}{b}_at\chi _i^{}`$ will only be of practical interest in the region where it also has a large tree-level branching ratio. There are several channels ($`\stackrel{~}{b}_ab\stackrel{~}{g}`$, $`\stackrel{~}{b}_ab\chi _\alpha ^0`$, $`\stackrel{~}{b}_2\stackrel{~}{b}_1h^0`$, …) that contribute to the sbottom-quark decay width. The gluino channel, if kinematically avaliable, saturates the total width, so in order to have an appreciable branching ratio $`\stackrel{~}{b}_at\chi _i^{}`$ we start out assuming that the gluino is much heavier than the squarks $`m_{\stackrel{~}{g}}>m_{\stackrel{~}{b}_a}`$, $`a`$=1, 2. Neutralino masses, on the other hand, are related to chargino masses; thus, no additional conditions can be imposed on this side.
Let us define the branching ratio for the decay under investigation:
$$BR_0(\stackrel{~}{b}_at\chi _1^{})=\frac{\mathrm{\Gamma }_0(\stackrel{~}{b}_a\chi _1^{}t)}{\mathrm{\Gamma }_0^T(\stackrel{~}{b}_a)},$$
(10)
where $`\mathrm{\Gamma }_0^T(\stackrel{~}{b}_a)`$ is the total $`\stackrel{~}{b}_a`$ decay width. This branching ratio is maximized in a scenario where the lightest chargino is higgsino-dominated and $`\mathrm{tan}\beta `$ is of low–moderate value. For large $`\mathrm{tan}\beta \text{ }\stackrel{>}{}\text{ }40`$, $`\mathrm{\Gamma }_0^T`$ is dominated by the neutral higgsino contribution.
Figure 1 displays the value of the branching ratio (10) as a function of $`\mathrm{tan}\beta `$, $`m_{\stackrel{~}{b}_1}`$ and $`\theta _{\stackrel{~}{b}}`$, for given values of the other parameters. It can be seen that low $`\mathrm{tan}\beta `$ values enhance the branching ratio. From now on we will concentrate in the region of $`\mathrm{tan}\beta 20`$; with this typical value the branching ratio still is appreciably high, whereas the electroweak corrections can be enhanced by means of the bottom Yukawa coupling (1). In Fig. 1(b) we can see the thresholds for opening the Higgs channels, namely $`\stackrel{~}{b}_2\stackrel{~}{b}_1h^0`$ (at the left end of the figure) and $`\stackrel{~}{b}_1\stackrel{~}{t}_1H^{}`$ (at its right end). When these channels are open, they tend to decrease the branching ratio (10) to undetectable small values. The large decay width into Higgs bosons results from large values for the $`A`$ parameters (8) in these kinematical region. Of course one could fix the input parameters (7) in such a way that $`A_{\{t,b\}}`$ are small in one of these regions (say at $`m_{\stackrel{~}{b}_1}`$ light), but at the price of making them large at its central value and even larger at the other end. This effect is also seen in Fig. 1(c), as the $`A`$ parameters are related to the angle trough (8). Note that the allowed range of $`\theta _{\stackrel{~}{b}}`$ is rather narrow, so that the physical bottom squark mass eigenstates basically coincide with the left- and right-handed chiral electroweak eigenstates.
## 3 One-loop corrections
The QCD one-loop corrections were originally computed in Refs.. Our QCD results presented here were computed independently and are in full agreement with those of . We include them in our discussion, for comparison with the residual ones, within the same scenario in which we computed the Yukawa part<sup>3</sup><sup>3</sup>3Our computation of the QCD effects can be found in .. The Yukawa corrections were first presented in <sup>4</sup><sup>4</sup>4The results presented here differ slightly from those of Ref. due to a recently discovered computer bug.. The full analytical results of the corrections can be found in . The QCD corrections contain all the gluon and gluino exchange diagrams, together with the soft and hard gluon bremsstrahlung, and the Yukawa corrections contain the diagrams in which Higgs bosons and higgsinos are exchanged. We use the on-shell renormalization scheme with the input parameters described in (7). The renormalization of the $`\mathrm{tan}\beta `$ parameter (necessary for the weak corrections) is fixed in such a way that the decay width $`\mathrm{\Gamma }(H^+\tau ^+\nu _\tau )`$ does not receive quantum corrections . The $`\mu `$ parameter is renormalized in analogy to fermion mass renormalization, since in our approximation it is the mass of the chargino involved in the decay. The bottom-squark mixing angle has to be renormalized as well. At variance with the other parameters appearing in our process, it is still not clear how this angle could be measured.<sup>5</sup><sup>5</sup>5For the top-squark mixing angle, a recent study has shown that a good precision can be obtained by measuring the production cross-section $`\sigma (e^+e^{}\stackrel{~}{t}_a\stackrel{~}{t}_b)`$, using polarized electrons, with the help of the polarization asymmetry . Hence we treat $`\theta _{\stackrel{~}{b}}`$ as a formal parameter and impose as a renormalization condition that it is not shifted by loop corrections from the mixed $`\stackrel{~}{b}_1\stackrel{~}{b}_2`$ self-energy <sup>6</sup><sup>6</sup>6Several different renormalization conditions for the squark mixing angle have been discussed in the literature, see e.g. and references therein..
The quantity under study will be the relative one-loop correction defined as:
$$\delta ^{ai}=\frac{\mathrm{\Gamma }(\stackrel{~}{b}_at\chi _i^{})\mathrm{\Gamma }_0(\stackrel{~}{b}_at\chi _i^{})}{\mathrm{\Gamma }_0(\stackrel{~}{b}_at\chi _i^{})}.$$
(11)
We start with the QCD corrections shown in Figs. 2-4. For the numerical evaluation we use $`\alpha _s(m_{\stackrel{~}{b}_a})`$, using the one-loop MSSM $`\beta `$-function, but, for the $`m_{\stackrel{~}{b}_a}`$ we use, it is basically the 4-flavour SM $`\beta `$-function, as the scale is almost always below the threshold of coloured SUSY particles (and top quark). In Fig. 2 we can see the evolution with $`\mathrm{tan}\beta `$ and $`\mu `$, which are the most interesting cases. The corrections are large ($`>10\%`$) and vary slowly for large values of $`\mathrm{tan}\beta `$ ($`\text{ }\stackrel{>}{}\text{ }20`$). We remark that for $`\mu <120\text{ GeV}`$ and $`\mathrm{tan}\beta >20`$ the corrections can be very large near the phase space limit of the lightest sbottom decay. However, this effect has nothing to do with the phase space exhaustion, but rather with the presence of a dynamical factor which goes to the denominator of $`\delta `$ in eq. (11). That factor is fixed by the structure of the interaction Lagrangian of the sbottom decay into charginos and top; for the parameters in Fig. 2, it turns out to vanish nearly at the phase space limit in the case of the lightest sbottom ($`\stackrel{~}{b}_1`$) decay. However, this is not so either for the heaviest sbottom ($`\stackrel{~}{b}_2`$) or for $`\mu >120\text{ GeV}`$ as it is patent in the same figure. The different evolution that exhibit the corrections of the two sbottoms has more relation with the electroweak nature of the process than with the purely QCD loops. For small angles $`\theta _{\stackrel{~}{b}}`$ and $`\theta _{\stackrel{~}{t}}`$ the squarks are nearly chiral, namely
$$\stackrel{~}{b}_1\stackrel{~}{b}_R,\stackrel{~}{b}_2\stackrel{~}{b}_L,\stackrel{~}{t}_1\stackrel{~}{t}_R,\stackrel{~}{t}_2\stackrel{~}{t}_L,$$
(12)
and thus their very different couplings to charginos translate into very different behaviours of (11) with $`\mathrm{tan}\beta `$ and $`\mu `$. In fact, the sbottom mixing angle plays a crucial role, as seen in Fig. 3; we also see, however, that its value is highly constrained by the condition (9). We should also comment on the effect of the gaugino mass parameter $`M`$ and the gluino mass in Fig. 4. The gluino evolution is rather flat once the pseudo-thresholds of $`\stackrel{~}{b}_ab\stackrel{~}{g}`$ are passed; thus, even if the gluino were much heavier than the squarks it would have an effect on the sbottom decay while at the same time it would prevent the otherwise dominant decay $`\stackrel{~}{b}_ab\stackrel{~}{g}`$. The correction is saturated for the gaugino mass parameter $`M\text{ }\stackrel{>}{}\text{ }200\text{ GeV}`$. Therefore the effects computed here can be compared with the ones obtained in the higgsino approximation discussed below. Finally, we point out the possible existence of non-decoupling effects in the QCD part. In it is shown that there exist a non-decoupling effect at large gluino masses, however this effect is numerically small and is not the one reflected in Fig. 4(a). The origin of the effect is related to the breaking of SUSY, specifically to the fact that the chargino coupling has a renormalization group evolution which is different to that of the gauge coupling in a non-SUSY world, the difference being sensitive to the splitting among the various SUSY scales – e.g. the scales of the squark and gluino masses.
The other parameters of the model present a rather mild effect on the corrections for squark masses in the ballpark of several hundreds of GeV. In summary the QCD corrections on the decay $`\stackrel{~}{b}_at\chi _i^{}`$ are large ($`20\%`$ for $`\stackrel{~}{b}_2`$, $`60\%`$ for $`\stackrel{~}{b}_1`$) and negative for values of the parameter space relevant to TESLA energies, with a higgsino-like chargino and moderate or large values of $`\mathrm{tan}\beta `$.
We now turn to the discussion of the Yukawa corrections where also non-decoupling effects may come into play. They have a different origin as compared to the pure QCD ones but they are also triggered by SUSY-breaking and can be numerically important. We remind that, in the computation of the Yukawa corrections, the higgsino approximation, eq. (2), was used, and so only the lightest chargino is avaliable for the decay. In the relevant large $`\mathrm{tan}\beta `$ segment under consideration, namely
$$20\text{ }\stackrel{<}{}\text{ }\mathrm{tan}\beta \text{ }\stackrel{<}{}\text{ }40,$$
(13)
the bottom quark Yukawa coupling $`\lambda _b`$ is comparable to the top quark Yukawa coupling $`\lambda _t`$. Even though the extreme interval $`40<\mathrm{tan}\beta <60`$ can be tolerated by perturbation theory, we shall confine ourselves to the moderate range (13). This is necessary to preserve the condition (9) for the typical set of sparticle masses used in our analysis.
The corresponding corrections $`\delta ^{ai}`$ (11) are shown in Figs. 5(a) and 5(b) as a function of the lightest stop and sbottom masses, respectively. The precise value of the lightest stop mass is an important parameter to determine the corrections to the lightest sbottom decay width. On the other hand the lightest sbottom mass does not play a major role for the corrections, aside from the presence of various thresholds. The allowed range for the sbottom and stop mixing angles is conditioned by the upper bound on the trilinear couplings and is obtained from eqs. (8) and (9). In the physical $`\theta _{\stackrel{~}{b}}`$ range, the variation of the correction (11) is shown in Fig. 6(a). The large values of the corrections far away from the allowed region (9) are due to the large values of the soft-SUSY-breaking trilinear coupling $`A_b`$ (8). On the other hand, the permitted range for the stop mixing angle $`\theta _{\stackrel{~}{t}}`$ is much larger, and we have plotted the corrections within the allowed region in Fig. 6(b). Note that the sign of the quantum effects for the lightest sbottom decay width changes within the domain of variation of $`\theta _{\stackrel{~}{t}}`$. Finally, we display the evolution of the SUSY-EW effects as a function of $`\mathrm{tan}\beta `$ (Fig. 7(a)) and of $`\mu `$ (Fig. 7(b)) within the region of compatibility with the constraint (9).
A few words are in order to explain the origin of the leading electroweak effects. One could expect that they come from the well-known large $`\mathrm{tan}\beta `$ enhancement stemming from the chargino-stop corrections to the bottom mass (see e.g. Ref. ). Nonetheless this is only partially true, since in the present case the remaining contributions can be sizeable enough. One can also think on the SUSY counterpart of the bottom mass counterterm corrections, that is, the finite contributions to the sbottom wave function renormalization constants , as an additional leading contribution. Both of these effects are of non-decoupling nature. However the addition of these two kind of contributions does not account for the total behaviour in all of the parameter space. To be more precise, in the region of the parameter space that we have dwelled upon the bottom mass contribution is seen to be dominant only for the lightest sbottom decay and for the lowest values of $`\mathrm{tan}\beta `$ in the range (13). This is indeed the case in Fig. 6(b) where $`\mathrm{tan}\beta =20`$ and therefore the bottom mass effect modulates the electroweak correction in this process and $`\delta ^{11}`$ becomes essentially an odd function of the stop mixing angle. This fact is easily understood since, as noted above, sbottoms are nearly chiral – eq. (12) – and the $`\stackrel{~}{b}_R`$ is the only one with couples with $`\lambda _b`$ – eq. (4). On the other hand, from Fig. 7(a) it is obvious that the (approximate) linear behaviour with $`\mathrm{tan}\beta `$ expected from bottom mass renormalization becomes completely distorted by the rest of the contributions, especially in the high $`\mathrm{tan}\beta `$ end. In short, the final electroweak correction cannot be simply ascribed to a single renormalization source but to the full Yukawa-coupling combined yield.
In general the SUSY-EW corrections to $`\mathrm{\Gamma }(\stackrel{~}{b}_at\chi _i^{})`$ are smaller than the QCD corrections. The reason why the electroweak corrections are smaller is in part due to the condition (9) restricting our analysis within the $`\mathrm{tan}\beta `$ interval (13). From Figs. 6 and 7(a) it is clear that outside this interval the SUSY-EW contributions could be much higher and with the same or opposite sign as the QCD effects, depending on the choice of the sign of the mixing angles. Moreover, since we have focused our analysis to sbottom masses accessible to TESLA, again the theoretical bound (9) severely restricts the maximum value of the trilinear couplings and this prevents the electroweak corrections from being larger. This cannot be cured by assuming larger values of $`\mu `$, because $`\mu `$ directly controls the value of the (higgsino-like) chargino mass, in the final state in the decay under study.
## 4 Conclusions
In summary, the MSSM corrections to squark decays into charginos can be significant and therefore must be included in any reliable analysis. The main corrections arise from the strongly interacting sector of the theory (i.e. the one involving gluons and gluinos), but also non-negligible effects may appear from the electroweak sector (characterized by chargino-neutralino exchange) at large (or very small) values of $`\mathrm{tan}\beta `$. In both cases non-decoupling effects related to the breaking of SUSY may be involved, but it is in the electroweak part where they can be numerically more sizeable. However, for sparticle masses of a few hundred $`GeV`$ a reliable estimate of the correction requires the calculation of the QCD and also of the complete Yukawa-coupling electroweak contribution. The QCD corrections are negative in most of the MSSM parameter space accessible to TESLA. They are of the order
$`\delta _{QCD}(\stackrel{~}{b}_1t\chi _1^{})`$ $``$ $`60\%`$
$`\delta _{QCD}(\stackrel{~}{b}_2t\chi _1^{})`$ $``$ $`20\%`$
for a wide range of the parameter space (Fig. 2). In certain corners of this space, though, they vary in a wide range of values. EW corrections can be of both signs. Our renormalization prescription uses the mixing angle between squarks as an input parameter. This prescription forces the physical region to be confined within a narrow range when we require compatibility with the non-existence of colour breaking vacua. Within this restricted region the typical corrections vary in the range (Figs. 6, 7)
$`\delta _{EW}(\stackrel{~}{b}_1t\chi _1^{})`$ $``$ $`+25\%\text{ to }15\%`$
$`\delta _{EW}(\stackrel{~}{b}_2t\chi _1^{})`$ $``$ $`+5\%\text{ to }5\%,`$
However we must recall that these limits are not rigorous. In the edge of such regions we find the largest EW contributions. We stress that for these decays it is not possible to narrow down the bulk of the electroweak corrections to just some simple-structured leading terms.
The present study has an impact on the determination of squark parameters at TESLA. The squark masses used in it would be available already for TESLA running at a center of mass energy of $`800\text{ GeV}`$. The large corrections found from both the QCD and the EW (Yukawa) sector, make this calculation necessary, not only for prospects of precision measurements in the sbottom-chargino-neutralino sectors, but also for a reliable first determination of their parameters.
## Acknowledgments
This work has been partially supported by the Deutsche Forschungsgemeinschaft and by CICYT under project No. AEN99-0766. |
warning/0001/hep-th0001086.html | ar5iv | text | # Untitled Document
BOUNDARY OPERATORS IN QUANTUM FIELD THEORY
Giampiero Esposito
Istituto Nazionale di Fisica Nucleare, Sezione di Napoli, Complesso Universitario di Monte S. Angelo, Via Cintia, Edificio N’, 80126 Napoli, Italy Università di Napoli Federico II, Dipartimento di Scienze Fisiche, Complesso Universitario di Monte S. Angelo, Via Cintia, Edificio N’, 80126 Napoli, Italy
The fundamental laws of physics can be derived from the requirement of invariance under suitable classes of transformations on the one hand, and from the need for a well-posed mathematical theory on the other hand. As a part of this programme, the present paper shows under which conditions the introduction of pseudo-differential boundary operators in one-loop Euclidean quantum gravity is compatible both with their invariance under infinitesimal diffeomorphisms and with the requirement of a strongly elliptic theory. Suitable assumptions on the kernel of the boundary operator make it therefore possible to overcome problems resulting from the choice of purely local boundary conditions.
PACS numbers: 03.70.+k, 04.70.Gw
1. Introduction
The aim of theoretical physics is to provide a clear conceptual framework for the wide variety of natural phenomena, so that not only are we able to make accurate predictions to be checked against observations, but the underlying mathematical structures of the world we live in can also become sufficiently well understood by the scientific community. What are therefore the key elements of a mathematical description of the physical world? Can we derive all basic equations of theoretical physics from a few symmetry principles? What do they tell us about the origin and evolution of the universe? Why is gravitation so peculiar with respect to all other fundamental interactions?
The above questions have received careful consideration over the last decades, and have led, in particular, to several approaches to a theory aiming at achieving a synthesis of quantum physics on the one hand, and general relativity on the other hand. This remains, possibly, the most important task of theoretical physics. The need for a quantum theory of gravity is already clear from singularity theorems in classical cosmology. Such theorems prove that the Einstein theory of general relativity leads to the occurrence of space-time singularities in a generic way. At first sight one might be tempted to conclude that a breakdown of all physical laws occurred in the past, or that general relativity is severely incomplete, being unable to predict what came out of a singularity. It has been therefore suggested that all these pathological features result from the attempt of using the Einstein theory well beyond its limit of validity, i.e. at energy scales where the fundamental theory is definitely more involved. General relativity might be therefore viewed as a low-energy limit of a richer theory, which achieves the synthesis of both the basic principles of modern physics and the fundamental interactions in the form presently known .
Within the framework just outlined it remains however true that the various approaches to quantum gravity developed so far suffer from mathematical inconsistencies, or incompleteness in their ability of accounting for some basic features of the laws of nature. From the point of view of general principles, the space-time approach to quantum mechanics and quantum field theory \[3–5\], and its application to the quantization of gravitational interactions, remains indeed of fundamental importance . When one tries to implement the Feynman “sum over histories” one discovers that, already at the level of non-relativistic quantum mechanics, a well defined mathematical formulation is only obtained upon considering a heat-equation problem. The measure occurring in the Feynman representation of the Green kernel is then meaningful, and the propagation amplitude of quantum mechanics in flat Minkowski space-time is obtained by analytic continuation. This is a clear indication that quantum-mechanical problems via path integrals are well understood only if the heat-equation counterpart is mathematically well posed. In quantum field theory one then deals with the Euclidean approach, and its application to quantum gravity relies heavily on the theory of elliptic operators on Riemannian manifolds . To obtain a complete picture one has then to specify the boundary conditions of the theory, i.e. the class of Riemannian geometries with their topologies involved in the sum, and the form of boundary data assigned on the bounding surfaces.
In particular, recent work has shown that the only set of local boundary conditions on metric perturbations which are completely invariant under infinitesimal diffeomorphisms is incompatible with the request of a good elliptic theory. More precisely, while the resulting operator on metric perturbations can be made of Laplace type and elliptic in the interior of the Riemannian manifold under consideration, the property of strong ellipticity is violated. This is a precise mathematical expression of the request that a unique smooth solution of the boundary-value problem should exist which vanishes at infinite geodesic distance from the boundary. This opens deep interpretive issues, since only for gravity does the request of complete gauge invariance of the boundary conditions turn out to be incompatible with a good elliptic theory . It is then impossible to make sense even just of the one-loop semiclassical approximation, because the functional trace of the heat operator is found to diverge .
We have been therefore led to consider non-local boundary conditions for the quantized gravitational field at one-loop level . On the one hand, such a scheme already arises in simpler problems, i.e. the quantum theory of a free particle subject to non-local boundary data on a circle . One then finds two families of eigenfunctions of the Hamiltonian: surface states which decrease exponentially as one moves away from the boundary, and bulk states which remain instead smooth and non-vanishing. The generalization to an Abelian gauge theory such as Maxwell theory can fulfill non-locality, ellipticity and complete gauge invariance of boundary conditions providing one learns to work with pseudo-differential operators in one-loop quantum theory . On the other hand, in the application to quantum gravity, since the boundary operator acquires new kernels responsible for the pseudo-differential nature of the boundary-value problem, one might hope to be able to recover a good elliptic theory under a wider variety of conditions.
This is precisely the aim of the present paper. After a survey of operators of Laplace type and of the associated boundary operators in section 2, section 3 introduces integro-differential boundary operators in Euclidean quantum gravity. Strong ellipticity of differential and pseudo-differential boundary-value problems is then defined in detail in section 4, and the application to Euclidean quantum gravity is studied in section 5. Further examples, of simpler nature, are given in section 6, and concluding remarks are presented in section 7.
2. Operators of Laplace type and their boundary operators
In the Euclidean approach to quantum field theory and quantum gravity one studies differentiable manifolds endowed with positive-definite metrics $`g`$, so that space-time is actually replaced by an $`m`$-dimensional Riemannian space $`(M,g)`$. An operator $`𝒫`$ of Laplace type, which acts on gauge fields, maps smooth sections of a vector bundle $`V`$ over $`M`$ into smooth sections of the same bundle, i.e.
$$𝒫:C^{\mathrm{}}(V,M)C^{\mathrm{}}(V,M),$$
and reads
$$𝒫=g^{ab}_a_bE,$$
$`(2.1)`$
where $`g^{ab}`$ is the contravariant form of the Riemannian metric for $`M`$, $``$ is the connection on $`V`$, and $`E`$ is an endomorphism. In Ref. a thorough investigation of boundary operators for elliptic operators of the form (2.1) has been performed. The key elements we need to recall are as follows.
If the manifold $`M`$ has a smooth non-empty boundary $`M`$, two vector bundles over $`M`$, hereafter denoted by $`W`$ and $`W^{}`$, yield a complete description of the problem. The boundary operator $`B`$ maps smooth sections of $`W`$ into smooth sections of $`W^{}`$:
$$B:C^{\mathrm{}}(W,M)C^{\mathrm{}}(W^{},M).$$
For mixed boundary conditions, the operator $`B`$ frequently reads
$$B\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \mathrm{\Lambda }& I\mathrm{\Pi }\end{array}\right),$$
$`(2.2)`$
where $`\mathrm{\Pi }`$ and $`I\mathrm{\Pi }`$ are complementary projectors, and $`\mathrm{\Lambda }`$ is a first-order tangential differential operator
$$\mathrm{\Lambda }(I\mathrm{\Pi })\left[\frac{1}{2}\right(\mathrm{\Gamma }^i\widehat{}_i+\widehat{}_i\mathrm{\Gamma }^i)+S](I\mathrm{\Pi }).$$
$`(2.3)`$
With our notation, $`\mathrm{\Gamma }^i`$ are endomorphism-valued vector fields on the boundary, $`\widehat{}`$ is the induced connection on $`M`$, and $`S`$ is an endomorphism on $`M`$. By virtue of (2.3) one has
$$\mathrm{\Pi }\mathrm{\Lambda }=\mathrm{\Lambda }\mathrm{\Pi }=0,$$
$`(2.4)`$
and hence $`B`$ is a projector, in that $`B^2=B`$. The boundary-value problem is meant to be the pair $`(𝒫,B)`$, where $`𝒫`$ is the operator (2.1) and $`B`$ is given in (2.2). The corresponding mixed boundary conditions read
$$\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \mathrm{\Lambda }& I\mathrm{\Pi }\end{array}\right)\left(\begin{array}{c}[\phi ]_M\\ [\phi _{;N}]_M\end{array}\right)=0,$$
$`(2.5)`$
where <sub>;N</sub> denotes covariant differentiation along the direction normal to the boundary, i.e. $`N^a_a`$. Moreover, the boundary operator (2.2) may be expressed in the form
$$B=PL,$$
$`(2.6)`$
where $`P`$ is the map
$$P:C^{\mathrm{}}(W,M)C^{\mathrm{}}(W^{},M)$$
given by
$$P\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ 0& I\mathrm{\Pi }\end{array}\right),$$
$`(2.7)`$
and $`L`$ is a map
$$L:C^{\mathrm{}}(W,M)C^{\mathrm{}}(W,M)$$
expressed in matrix form as
$$L\left(\begin{array}{cc}I& 0\\ \mathrm{\Lambda }& I\end{array}\right).$$
$`(2.8)`$
Interestingly, the operator $`P`$ is itself a projector: $`P^2=P`$, whereas $`L`$ is non-singular, with inverse
$$L^1=\left(\begin{array}{cc}I& 0\\ \mathrm{\Lambda }^1& I\end{array}\right).$$
$`(2.9)`$
The “column vector” used in Eq. (2.5), i.e.
$$\psi (\phi )\left(\begin{array}{c}[\phi ]_M\\ [\phi _{;N}]_M\end{array}\right),$$
$`(2.10)`$
is a section of the bundle $`W`$ of boundary data, whereas the auxiliary vector bundle $`W^{}`$ has sections given by (see the main diagonal of $`B`$ in (2.2))
$$\psi ^{}(\phi )\left(\begin{array}{c}\mathrm{\Pi }[\phi ]_M\\ (I\mathrm{\Pi })[\phi _{;N}]_M\end{array}\right).$$
$`(2.11)`$
At this stage, a naturally occurring question is under which conditions a projector $`P`$ gives rise to a projector $`B`$ such that $`B=PL`$ as in Eq. (2.6). To obtain equations in a form as general as possible we replace $`P`$ defined in (2.7) by the $`2\times 2`$ matrix
$$P\left(\begin{array}{cc}\alpha & \beta \\ \gamma & \delta \end{array}\right),$$
$`(2.12)`$
where $`\alpha ,\beta ,\gamma ,\delta `$ are, for the time being, some unknown operators to be determined by imposing suitable restrictions (see below). The projector condition $`P^2=P`$ yields therefore four operator equations, i.e.
$$\alpha ^2+\beta \gamma =\alpha ,$$
$`(2.13)`$
$$\alpha \beta +\beta \delta =\beta ,$$
$`(2.14)`$
$$\gamma \alpha +\delta \gamma =\gamma ,$$
$`(2.15)`$
$$\gamma \beta +\delta ^2=\delta .$$
$`(2.16)`$
A particular solution of Eqs. (2.13)–(2.16) is given by the case in which
$$\beta =\gamma =0,$$
$`(2.17)`$
$$\alpha ^2=\alpha ,$$
$`(2.18)`$
$$\delta ^2=\delta ,$$
$`(2.19)`$
$$\alpha +\delta =I.$$
$`(2.20)`$
This yields the operator $`P`$ in the form (2.7) appropriate for the Grubb–Gilkey–Smith boundary-value problem . If the conditions (2.17)–(2.20) are not fulfilled, one gets instead from Eqs. (2.13)–(2.16) the equations
$$\alpha (\alpha I)=\beta \gamma ,$$
$`(2.21)`$
$$\alpha =I\beta \delta \beta ^1,$$
$`(2.22)`$
$$\alpha =I\gamma ^1\delta \gamma ,$$
$`(2.23)`$
$$\delta (\delta I)=\gamma \beta ,$$
$`(2.24)`$
provided that $`\beta `$ and $`\gamma `$ can be inverted.
3. Euclidean quantum gravity
In Euclidean quantum gravity, mixed boundary conditions on metric perturbations $`h_{cd}`$ occur naturally if one requires their complete invariance under infinitesimal diffeomorphisms, as is proved in detail in Ref. . On denoting by $`N^a`$ the inward-pointing unit normal to the boundary, by
$$q_b^a\delta _b^aN^aN_b$$
$`(3.1)`$
the projector of tensor fields onto $`M`$, with associated projection operator
$$\mathrm{\Pi }_{ab}^{cd}q_{(a}^cq_{b)}^d,$$
$`(3.2)`$
the gauge-invariant boundary conditions for one-loop quantum gravity read
$$\left[\mathrm{\Pi }_{ab}^{cd}h_{cd}\right]_M=0,$$
$`(3.3)`$
$$\left[\mathrm{\Phi }_a(h)\right]_M=0,$$
$`(3.4)`$
where $`\mathrm{\Phi }_a`$ is the gauge-averaging functional necessary to obtain an invertible operator $`P_{ab}^{cd}`$ on metric perturbations. When $`P_{ab}^{cd}`$ is chosen to be of Laplace type, $`\mathrm{\Phi }_a`$ reduces to the familiar de Donder term
$$\mathrm{\Phi }_a(h)=^b(h_{ab}\frac{1}{2}g_{ab}g^{cd}h_{cd})=E_a^{bcd}_bh_{cd},$$
$`(3.5)`$
where $`E^{abcd}`$ is the DeWitt supermetric on the vector bundle of symmetric rank-two tensor fields over $`M`$ ($`g`$ being the metric on $`M`$):
$$E^{abcd}\frac{1}{2}(g^{ac}g^{bd}+g^{ad}g^{bc}g^{ab}g^{cd}).$$
$`(3.6)`$
The boundary conditions (3.3) and (3.4) can then be cast in the Grubb–Gilkey–Smith form (2.5), where $`\mathrm{\Lambda }`$ is the first-order operator on the boundary defined in Eq. (2.3). However, the work in Ref. has shown that an operator of Laplace type on metric perturbations is then incompatible with the requirement of strong ellipticity of the boundary-value problem (see section 4), because the operator $`\mathrm{\Lambda }`$ contains tangential derivatives of metric perturbations.
To take care of this serious drawback, the work in Refs. has proposed to consider in the boundary condition (3.4) a gauge-averaging functional given by the de Donder term (3.5) plus an integro-differential operator on metric perturbations, i.e.
$$\mathrm{\Phi }_a(h)E_a^{bcd}_bh_{cd}+_M\zeta _a^{cd}(x,x^{})h_{cd}(x^{})𝑑V^{}.$$
$`(3.7)`$
We now begin by remarking that the resulting boundary conditions can be cast in the form
$$\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \mathrm{\Lambda }+\stackrel{~}{\mathrm{\Lambda }}& I\mathrm{\Pi }\end{array}\right)\left(\begin{array}{c}[\phi ]_M\\ [\phi _{;N}]_M\end{array}\right)=0,$$
$`(3.8)`$
where $`\stackrel{~}{\mathrm{\Lambda }}`$ reflects the occurrence of the integral over $`M`$ in Eq. (3.7). It is convenient to work first in a general way and then consider the form taken by these operators in the gravitational case. On requiring that the resulting boundary operator
$$=\left(\begin{array}{cc}_{11}& _{12}\\ _{21}& _{22}\end{array}\right)\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \mathrm{\Lambda }+\stackrel{~}{\mathrm{\Lambda }}& I\mathrm{\Pi }\end{array}\right)$$
$`(3.9)`$
should remain a projector: $`^2=`$, we find the condition
$$(\mathrm{\Lambda }+\stackrel{~}{\mathrm{\Lambda }})\mathrm{\Pi }\mathrm{\Pi }(\mathrm{\Lambda }+\stackrel{~}{\mathrm{\Lambda }})=0,$$
$`(3.10)`$
which reduces to
$$\mathrm{\Pi }\stackrel{~}{\mathrm{\Lambda }}=\stackrel{~}{\mathrm{\Lambda }}\mathrm{\Pi },$$
$`(3.11)`$
by virtue of (2.4).
In Euclidean quantum gravity at one-loop level, Eq. (3.11) leads to
$$\mathrm{\Pi }_{ac}^{br}(x)_M\zeta _b^{cq}(x,x^{})h_{qr}(x^{})𝑑V^{}=_M\zeta _a^{cd}(x,x^{})\mathrm{\Pi }_{cd}^{qr}(x^{})h_{qr}(x^{})𝑑V^{},$$
$`(3.12)`$
which can be re-expressed in the form
$$_M\left[\mathrm{\Pi }_{ac}^{br}(x)\zeta _b^{cq}(x,x^{})\zeta _a^{cd}(x,x^{})\mathrm{\Pi }_{cd}^{qr}(x^{})\right]h_{qr}(x^{})𝑑V^{}=0.$$
$`(3.13)`$
Since this should hold for all $`h_{qr}(x^{})`$, it eventually leads to the vanishing of the term in square brackets in the integrand. The notation $`\zeta _b^{cq}(x,x^{})`$ is indeed rather awkward, because there is an even number of arguments, i.e. $`x`$ and $`x^{}`$, with an odd number of indices. Hereafter, we therefore assume that a vector field $`T`$ and kernel $`\stackrel{~}{\zeta }`$ exist such that
$$\zeta _b^{cq}(x,x^{})T^p(x)\stackrel{~}{\zeta }_{bp}^{cq}(x,x^{})T^p\stackrel{~}{\zeta }_{bp}^{c^{}q^{}}.$$
$`(3.14)`$
The projector condition (3.11) is therefore satisfied if and only if
$$T^p(x)\left[\mathrm{\Pi }_{ac}^{br}(x)\stackrel{~}{\zeta }_{bp}^{cq}(x,x^{})\stackrel{~}{\zeta }_{ap}^{cd}(x,x^{})\mathrm{\Pi }_{cd}^{qr}(x^{})\right]=0.$$
$`(3.15)`$
4. Strong ellipticity
We are here concerned with the issue of ellipticity of the boundary-value problem of section 3. For this purpose, we begin by recalling what is known about ellipticity of the Laplacian (hereafter $`P`$) on a Riemannian manifold with smooth boundary. This concept is studied in terms of the leading symbol of $`P`$. It is indeed well known that the Fourier transform makes it possible to associate to a differential operator of order $`k`$ a polynomial of degree $`k`$, called the characteristic polynomial or symbol. The leading symbol, $`\sigma _L`$, picks out the highest order part of this polynomial. For the Laplacian, it reads
$$\sigma _L(P;x,\xi )=|\xi |^2I=g^{\mu \nu }\xi _\mu \xi _\nu I.$$
$`(4.1)`$
With a standard notation, $`(x,\xi )`$ are local coordinates for $`T^{}(M)`$, the cotangent bundle of $`M`$. The leading symbol of $`P`$ is trivially elliptic in the interior of $`M`$, since the right-hand side of (4.1) is positive-definite, and one has
$$\mathrm{det}[\sigma _L(P;x,\xi )\lambda ]=(|\xi |^2\lambda )^{\mathrm{dim}V}0,$$
$`(4.2)`$
for all $`\lambda 𝒞𝐑_+`$. In the presence of a boundary, however, one needs a more careful definition of ellipticity. First, for a manifold $`M`$ of dimension $`m`$, the $`m`$ coordinates $`x`$ are split into $`m1`$ local coordinates on $`M`$, hereafter denoted by $`\left\{\widehat{x}^k\right\}`$, and $`r`$, the geodesic distance to the boundary. Moreover, the $`m`$ coordinates $`\xi _\mu `$ are split into $`m1`$ coordinates $`\left\{\zeta _j\right\}`$ (with $`\zeta `$ being a cotangent vector on the boundary), jointly with a real parameter $`\omega T^{}(𝐑)`$. At a deeper level, all this reflects the split
$$T^{}(M)=T^{}(M)T^{}(𝐑)$$
$`(4.3)`$
in a neighbourhood of the boundary .
The ellipticity we are interested in requires now that $`\sigma _L`$ should be elliptic in the interior of $`M`$, as specified before, and that strong ellipticity should hold. This means that a unique solution exists of the differential equation obtained from the leading symbol:
$$[\sigma _L(P;\left\{\widehat{x}^k\right\},r=0,\left\{\zeta _j\right\},\omega i\frac{}{r})\lambda ]\phi (r,\widehat{x},\zeta ;\lambda )=0,$$
$`(4.4)`$
subject to the boundary conditions
$$\sigma _g(B)(\left\{\widehat{x}^k\right\},\left\{\zeta _j\right\})\psi (\phi )=\psi ^{}(\phi )$$
$`(4.5)`$
and to the asymptotic condition
$$\underset{r\mathrm{}}{lim}\phi (r,\widehat{x},\zeta ;\lambda )=0.$$
$`(4.6)`$
In Eq. (4.5), $`\sigma _g`$ is the graded leading symbol of the boundary operator of section 2 in the local coordinates $`\left\{\widehat{x}^k\right\},\left\{\zeta _j\right\}`$, and is given by
$$\sigma _g(B)=\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ i\mathrm{\Gamma }^j\zeta _j& I\mathrm{\Pi }\end{array}\right).$$
$`(4.7)`$
Roughly speaking, the above construction uses Fourier transform and the inward geodesic flow to obtain the ordinary differential equation (4.4) from the Laplacian, with corresponding Fourier transform (4.5) of the original boundary conditions. The asymptotic condition (4.6) picks out the solutions of Eq. (4.4) which satisfy Eq. (4.5) with arbitrary boundary data $`\psi ^{}(\phi )`$ (see (2.11)) and vanish at infinite geodesic distance to the boundary. When all the above conditions are satisfied $`\zeta T^{}(M),\lambda 𝒞𝐑_+,(\zeta ,\lambda )(0,0)`$ and $`\psi ^{}(\phi )C^{\mathrm{}}(W^{},M)`$, the boundary-value problem $`(P,B)`$ for the Laplacian is said to be strongly elliptic with respect to the cone $`𝒞𝐑_+`$.
However, when the gauge-averaging functional (3.7) is used in the boundary condition (3.4), the work in Ref. has proved that the operator on metric perturbations takes the form of an operator of Laplace type $`P_{ab}^{cd}`$ plus an integral operator $`G_{ab}^{cd}`$. Explicitly, one finds (with $`R_{bcd}^a`$ being the Riemann curvature of the background geometry $`(M,g)`$)
$$P_{ab}^{cd}=E_{ab}^{cd}(\text{ }\text{ / }\text{ }+R)2E_{ab}^{qf}R_{qpf}^cg^{dp}E_{ab}^{pd}R_p^cE_{ab}^{cp}R_p^d,$$
$`(4.8)`$
$$G_{ab}^{cd}=U_{ab}^{cd}+V_{ab}^{cd},$$
$`(4.9)`$
where
$$U_{ab}^{cd}h_{cd}(x)=2E_{rsab}^r_MT^p(x)\stackrel{~}{\zeta }_p^{scd}(x,x^{})h_{cd}(x^{})𝑑V^{},$$
$`(4.10)`$
$$h^{ab}V_{ab}^{cd}h_{cd}(x)=_{M^2}h^{ab}(x^{})T^q(x)\stackrel{~}{\zeta }_{pqab}(x,x^{})T^r(x)\stackrel{~}{\zeta }_r^{pcd}(x,x^{\prime \prime })h_{cd}(x^{\prime \prime })𝑑V^{}𝑑V^{\prime \prime }.$$
$`(4.11)`$
We now assume that the operator on metric perturbations, which is so far an integro-differential operator defined by a kernel, is also pseudo-differential. This means that it can be characterized by suitable regularity properties obeyed by the symbol. More precisely, let $`S^d`$ be the set of all symbols $`p(x,\xi )`$ such that (1) $`p`$ is $`C^{\mathrm{}}`$ in $`(x,\xi )`$, with compact $`x`$ support. (2) For all $`(\alpha ,\beta )`$, there exist constants $`C_{\alpha ,\beta }`$ for which
$$\begin{array}{ccc}& \left|(i)^{_{k=1}^m(\alpha _k+\beta _k)}\left(\frac{}{x_1}\right)^{\alpha _1}\mathrm{}\left(\frac{}{x_m}\right)^{\alpha _m}\left(\frac{}{\xi _1}\right)^{\beta _1}\mathrm{}\left(\frac{}{\xi _m}\right)^{\beta _m}p(x,\xi )\right|\hfill & \\ & C_{\alpha ,\beta }\left(1+\sqrt{g^{ab}(x)\xi _a\xi _b}\right)^{d_{k=1}^m\beta _k},\hfill & (4.12)\hfill \end{array}$$
for some real (not necessarily positive) value of $`d`$. The associated pseudo-differential operator, defined on the Schwarz space and taking values in the set of smooth functions on $`M`$ with compact support:
$$P:𝒮C_c^{\mathrm{}}(M)$$
acts according to
$$Pf(x)e^{i(xy)\xi }p(x,\xi )f(y)\mu (y,\xi ),$$
$`(4.13)`$
where $`\mu (y,\xi )`$ is here meant to be the invariant integration measure with respect to $`y_1,\mathrm{},y_m`$ and $`\xi _1,\mathrm{},\xi _m`$. Actually, one first gives the definition for pseudo-differential operators $`P:𝒮C_c^{\mathrm{}}(𝐑^m)`$, eventually proving that a coordinate-free definition can be given and extended to smooth Riemannian manifolds .
In the presence of pseudo-differential operators, both ellipticity in the interior of $`M`$ and strong ellipticity of the boundary-value problem need a more involved formulation. In our paper, inspired by the flat-space analysis in Ref. , we make the following requirements.
4.1 Ellipticity in the interior
Let $`U`$ be an open subset with compact closure in $`M`$, and consider an open subset $`U_1`$ whose closure $`\overline{U}_1`$ is properly included into $`U`$: $`\overline{U}_1U`$. If $`p`$ is a symbol of order $`d`$ on $`U`$, it is said to be elliptic on $`U_1`$ if there exists an open set $`U_2`$ which contains $`\overline{U}_1`$ and positive constants $`C_0,C_1`$ so that
$$|p(x,\xi )|^1C_1(1+|\xi |)^d,$$
$`(4.14)`$
for $`|\xi |C_0`$ and $`xU_2`$, where $`|\xi |\sqrt{g^{ab}(x)\xi _a\xi _b}`$. The corresponding operator $`P`$ is then elliptic.
4.2 Strong ellipticity in the absence of boundaries
Let us assume that the symbol under consideration is polyhomogeneous, in that it admits an asymptotic expansion of the form
$$p(x,\xi )\underset{l=0}{\overset{\mathrm{}}{}}p_{dl}(x,\xi ),$$
$`(4.15)`$
where each term $`p_{dl}`$ has the homogeneity property
$$p_{dl}(x,t\xi )=t^{dl}p_{dl}(x,\xi )\mathrm{if}t1\mathrm{and}|\xi |1.$$
$`(4.16)`$
The leading symbol is then, by definition,
$$p^0(x,\xi )p_d(x,\xi ).$$
$`(4.17)`$
Strong ellipticity in the absence of boundaries is formulated in terms of the leading symbol, and it requires that
$$\mathrm{Re}p^0(x,\xi )c(x)|\xi |^d,$$
$`(4.18)`$
where $`xM`$ and $`|\xi |1`$, $`c`$ being a positive function on $`M`$. It can then be proved that the Gärding inequality holds, according to which, for any $`\epsilon >0`$,
$$\mathrm{Re}(Pu,u)bu_{\frac{d}{2}}^2b_1u_{\frac{d}{2}\epsilon }^2\mathrm{for}uH^{\frac{d}{2}}(M),$$
$`(4.19)`$
with $`b>0`$, where $`H^s(M)`$ is the standard notation for Sobolev spaces, for all $`s`$ .
4.3 Strong ellipticity in the presence of boundaries
The homogeneity property (4.16) only holds for $`t1`$ and $`|\xi |1`$. Consider now the case $`l=0`$, for which one obtains the leading symbol which plays the key role in the definition of ellipticity. If $`p^0(x,\xi )p_d(x,\xi )\sigma _L(P;x,\xi )`$ is not a polynomial (which corresponds to the genuinely pseudo-differential case) while being a homogeneous function of $`\xi `$, it is irregular at $`\xi =0`$. When $`|\xi |1`$, the only control over the leading symbol is provided by estimates of the form
$$\begin{array}{ccc}& \left|(i)^{_{k=1}^m(\alpha _k+\beta _k)}\left(\frac{}{x_1}\right)^{\alpha _1}\mathrm{}\left(\frac{}{x_m}\right)^{\alpha _m}\left(\frac{}{\xi _1}\right)^{\beta _1}\mathrm{}\left(\frac{}{\xi _m}\right)^{\beta _m}p^0(x,\xi )\right|\hfill & \\ & c(x)\xi ^{d|\beta |}.\hfill & (4.20)\hfill \end{array}$$
We therefore come to appreciate the problematic aspect of symbols of pseudo-differential operators . The singularity at $`\xi =0`$ can be dealt with either by modifying the leading symbol for small $`\xi `$ to be a $`C^{\mathrm{}}`$ function (at the price of loosing the homogeneity there), or by keeping the strict homogeneity and dealing with the singularity at $`\xi =0`$ .
On the other hand, we are interested in a definition of strong ellipticity of pseudo-differential boundary-value problems that reduces to Eqs. (4.4)–(4.6) when both $`P`$ and the boundary operator reduce to the form considered in section 2. For this purpose, and bearing in mind the occurrence of singularities in the leading symbols of $`P`$ and of the boundary operator, we make the following requirements.
Let $`(P+G)`$ be a pseudo-differential operator subject to boundary conditions described by the pseudo-differential boundary operator $``$ (the consideration of $`(P+G)`$ rather than only $`P`$ is necessary to achieve self-adjointness, as is described in detail in Refs. and ). The pseudo-differential boundary-value problem $`((P+G),)`$ is strongly elliptic with respect to $`𝒞𝐑_+`$ if: (I) The inequalities (4.14) and (4.18) hold; (II) There exists a unique solution of the equation
$$[\sigma _L((P+G);\left\{\widehat{x}^k\right\},r=0,\left\{\zeta _j\right\},\omega i\frac{}{r})\lambda ]\phi (r,\widehat{x},\zeta ;\lambda )=0,$$
$`(4.4^{})`$
subject to the boundary conditions
$$\sigma _L()(\left\{\widehat{x}^k\right\},\left\{\zeta _j\right\})\psi (\phi )=\psi ^{}(\phi )$$
$`(4.5^{})`$
and to the asymptotic condition (4.6). It should be stressed that, unlike the case of differential operators, Eq. (4.4’) is not an ordinary differential equation in general, because $`(P+G)`$ is pseudo-differential. (III) The strictly homogeneous symbols associated to $`(P+G)`$ and $``$ have limits for $`|\zeta |0`$ in the respective leading symbol norms, with the limiting symbol restricted to the boundary which avoids the values $`\lambda 𝒞𝐑_+`$ for all $`\left\{\widehat{x}\right\}`$.
Condition (III) requires a last effort for a proper understanding. Given a pseudo-differential operator of order $`d`$ with leading symbol $`p^0(x,\xi )`$, the associated strictly homogeneous symbol is defined by
$$p^h(x,\xi )|\xi |^dp^0(x,\frac{\xi }{|\xi |})\mathrm{for}\xi 0.$$
$`(4.21)`$
This extends to a continuous function vanishing at $`\xi =0`$ when $`d>0`$. In the presence of boundaries, the boundary-value problem $`((P+G),)`$ has a strictly homogeneous symbol on the boundary equal to (some indices are omitted for simplicity)
$$\left(\begin{array}{c}p^h(\left\{\widehat{x}\right\},r=0,\left\{\zeta \right\},i\frac{}{r})+g^h(\left\{\widehat{x}\right\},\left\{\zeta \right\},i\frac{}{r})\lambda \\ b^h(\left\{\widehat{x}\right\},\left\{\zeta \right\},i\frac{}{r})\end{array}\right),$$
where $`p^h,g^h`$ and $`b^h`$ are the strictly homogeneous symbols of $`P,G`$ and $``$ respectively, obtained from the corresponding leading symbols $`p^0,g^0`$ and $`b^0`$ via equations analogous to (4.21), after taking into account the split (4.3), and upon replacing $`\omega `$ by $`i\frac{}{r}`$. The limiting symbol restricted to the boundary (also called limiting $`\lambda `$-dependent boundary symbol operator) and mentioned in condition III reads therefore
$$\begin{array}{ccc}& a^h(\left\{\widehat{x}\right\},r=0,\zeta =0,i\frac{}{r})\hfill & \\ & =\left(\begin{array}{c}p^h(\left\{\widehat{x}\right\},r=0,\zeta =0,i\frac{}{r})+g^h(\left\{\widehat{x}\right\},\zeta =0,i\frac{}{r})\lambda \\ b^h(\left\{\widehat{x}\right\},\zeta =0,i\frac{}{r})\end{array}\right),\hfill & (4.22)\hfill \end{array}$$
where the singularity at $`\xi =0`$ of the leading symbol in absence of boundaries is replaced by the singularity at $`\zeta =0`$ of the leading symbols of $`P,G`$ and $``$ when a boundary occurs.
5. Application of the strong ellipticity criterion
Let us now see how the previous conditions on the leading symbol of $`(P+G)`$ and on the graded leading symbol of the boundary operator can be used. The equation (4.4’) is solved by a function $`\phi `$ depending on $`r,\left\{\widehat{x}^k\right\},\left\{\zeta _j\right\}`$ and, parametrically, on the eigenvalues $`\lambda `$. For simplicity, we write $`\phi =\phi (r,\widehat{x},\zeta ;\lambda )`$, omitting indices. Since the leading symbol is no longer a polynomial when $`(P+G)`$ is genuinely pseudo-differential, we cannot make any further specification on $`\phi `$ at this stage, apart from requiring that it should reduce to (here $`|\zeta |^2\zeta _i\zeta ^i`$)
$$\chi (\widehat{x},\zeta )e^{r\sqrt{|\zeta |^2\lambda }}$$
when $`(P+G)`$ reduces to a Laplacian (and hence $`\mathrm{\Lambda }`$ reduces to (2.3)).
The equation (4.5’) involves the graded leading symbol of $``$ and restrictions to the boundary of the field and its covariant derivative along the normal direction. Such a restriction is obtained by setting to zero the geodesic distance $`r`$, and hence we write, in general form (here we denote again by $`\mathrm{\Lambda }`$ the full matrix element $`_{21}`$ in the boundary operator (3.9)),
$$\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \sigma _L(\mathrm{\Lambda })& I\mathrm{\Pi }\end{array}\right)\left(\begin{array}{c}\phi (0,\widehat{x},\zeta ;\lambda )\\ \phi ^{}(0,\widehat{x},\zeta ;\lambda )\end{array}\right)=\left(\begin{array}{c}\mathrm{\Pi }\rho (0,\widehat{x},\zeta ;\lambda )\\ (I\mathrm{\Pi })\rho ^{}(0,\widehat{x},\zeta ;\lambda )\end{array}\right),$$
$`(5.1)`$
where $`\rho `$ differs from $`\phi `$, because Eq. (4.5’) is written for $`\psi (\phi )`$ and $`\psi ^{}(\phi )\psi (\phi )`$. Now Eq. (5.1) leads to
$$\mathrm{\Pi }\phi (0,\widehat{x},\zeta ;\lambda )=\mathrm{\Pi }\rho (0,\widehat{x},\zeta ;\lambda ),$$
$`(5.2)`$
$$\sigma _L(\mathrm{\Lambda })\phi (0,\widehat{x},\zeta ;\lambda )+(I\mathrm{\Pi })\phi ^{}(0,\widehat{x},\zeta ;\lambda )=(I\mathrm{\Pi })\rho ^{}(0,\widehat{x},\zeta ;\lambda ),$$
$`(5.3)`$
and we require that, for $`\phi `$ solving Eq. (4.4’) and the asymptotic decay (4.6), with $`\lambda 𝒞𝐑_+`$, Eqs. (5.2) and (5.3) can be always solved with given values of $`\rho (0,\widehat{x},\zeta ;\lambda )`$ and $`\rho ^{}(0,\widehat{x},\zeta ;\lambda )`$, whenever $`(\zeta ,\lambda )(0,0)`$. The idea is now to relate, if possible, $`\phi ^{}(0,\widehat{x},\zeta ;\lambda )`$ to $`\phi (0,\widehat{x},\zeta ;\lambda )`$ in such a way that Eq. (5.2) can be used to simplify Eq. (5.3). For this purpose, we consider the function $`f`$ such that
$$\frac{\phi ^{}(0,\widehat{x},\zeta ;\lambda )}{\phi (0,\widehat{x},\zeta ;\lambda )}=\frac{\rho ^{}(0,\widehat{x},\zeta ;\lambda )}{\rho (0,\widehat{x},\zeta ;\lambda )}=f(\widehat{x},\zeta ;\lambda ),$$
$`(5.4)`$
$$\mathrm{\Pi }(\widehat{x})f(\widehat{x},\zeta ;\lambda )=f(\widehat{x},\zeta ;\lambda )\mathrm{\Pi }(\widehat{x}).$$
$`(5.5)`$
If both (5.4) and (5.5) hold, Eq. (5.3) reduces indeed to
$$\begin{array}{ccc}& \sigma _L(\mathrm{\Lambda })\phi (0,\widehat{x},\zeta ;\lambda )+f(\widehat{x},\zeta ;\lambda )(\phi (0,\widehat{x},\zeta ;\lambda )\rho (0,\widehat{x},\zeta ;\lambda ))\hfill & \\ & =f(\widehat{x},\zeta ;\lambda )\mathrm{\Pi }(\phi (0,\widehat{x},\zeta ;\lambda )\rho (0,\widehat{x},\zeta ;\lambda )),\hfill & (5.6a)\hfill \end{array}$$
and hence, by virtue of (5.2),
$$[\sigma _L(\mathrm{\Lambda })+f(\widehat{x},\zeta ;\lambda )]\phi (0,\widehat{x},\zeta ;\lambda )=\rho ^{}(0,\widehat{x},\zeta ;\lambda ).$$
$`(5.6b)`$
Thus, the strong ellipticity condition with respect to $`𝒞𝐑_+`$ implies in this case the invertibility of $`[\sigma _L(\mathrm{\Lambda })+f(\widehat{x},\zeta ;\lambda )]`$, i.e.
$$\mathrm{det}[\sigma _L(\mathrm{\Lambda })+f(\widehat{x},\zeta ;\lambda )]0\lambda 𝒞𝐑_+.$$
$`(5.7)`$
Moreover, by virtue of the identity
$$[f(\widehat{x},\zeta ;\lambda )+\sigma _L(\mathrm{\Lambda })][f(\widehat{x},\zeta ;\lambda )\sigma _L(\mathrm{\Lambda })]=[f^2(\widehat{x},\zeta ;\lambda )\sigma _L^2(\mathrm{\Lambda })],$$
$`(5.8)`$
the condition (5.7) is equivalent to
$$\mathrm{det}[f^2(\widehat{x},\zeta ;\lambda )\sigma _L^2(\mathrm{\Lambda })]0\lambda 𝒞𝐑_+.$$
$`(5.9)`$
Since $`f(\widehat{x},\zeta ;\lambda )`$ is, in general, complex-valued, one can always express it in the form
$$f(\widehat{x},\zeta ;\lambda )=\mathrm{Re}f(\widehat{x},\zeta ;\lambda )+i\mathrm{Im}f(\widehat{x},\zeta ;\lambda ),$$
$`(5.10)`$
so that (5.9) reads eventually
$$\mathrm{det}[\mathrm{Re}^2f(\widehat{x},\zeta ;\lambda )\mathrm{Im}^2f(\widehat{x},\zeta ;\lambda )\sigma _L^2(\mathrm{\Lambda })+2i\mathrm{Re}f(\widehat{x},\zeta ;\lambda )\mathrm{Im}f(\widehat{x},\zeta ;\lambda )]0.$$
$`(5.11)`$
In particular, when
$$\mathrm{Im}f(\widehat{x},\zeta ;\lambda )=0,$$
$`(5.12)`$
condition (5.11) reduces to
$$\mathrm{det}[\mathrm{Re}^2f(\widehat{x},\zeta ;\lambda )\sigma _L^2(\mathrm{\Lambda })]0.$$
$`(5.13)`$
A sufficient condition for strong ellipticity with respect to the cone $`𝒞𝐑_+`$ is therefore the negative-definiteness of $`\sigma _L^2(\mathrm{\Lambda })`$:
$$\sigma _L^2(\mathrm{\Lambda })<0,$$
$`(5.14)`$
so that
$$\mathrm{Re}^2f(\widehat{x},\zeta ;\lambda )\sigma _L^2(\mathrm{\Lambda })>0,$$
$`(5.15)`$
and hence (5.13) is fulfilled.
In the derivation of the sufficient conditions (5.11) and (5.14), the assumption (5.5) plays a crucial role. In general, however, $`\mathrm{\Pi }`$ and $`f`$ have a non-vanishing commutator, and hence a $`C(\widehat{x},\zeta ;\lambda )`$ exists such that
$$\mathrm{\Pi }(\widehat{x})f(\widehat{x},\zeta ;\lambda )f(\widehat{x},\zeta ;\lambda )\mathrm{\Pi }(\widehat{x})=C(\widehat{x},\zeta ;\lambda ).$$
$`(5.16)`$
The occurrence of $`C`$ is a peculiar feature of the fully pseudo-differential framework. Equation (5.3) is then equivalent to (now we write explicitly also the independent variables in the leading symbol of $`\mathrm{\Lambda }`$)
$$\begin{array}{ccc}& [(\sigma _L(\mathrm{\Lambda })C)(\widehat{x},\zeta ;\lambda )+f(\widehat{x},\zeta ;\lambda )]\phi (0,\widehat{x},\zeta ;\lambda )\hfill & \\ & =\rho ^{}(0,\widehat{x},\zeta ;\lambda )C(\widehat{x},\zeta ;\lambda )\rho (0,\widehat{x},\zeta ;\lambda ).\hfill & (5.17)\hfill \end{array}$$
On defining
$$\gamma (\widehat{x},\zeta ;\lambda )[\sigma _L(\mathrm{\Lambda })C](\widehat{x},\zeta ;\lambda ),$$
$`(5.18)`$
we therefore obtain strong ellipticity conditions formally analogous to (5.7) or (5.11) or (5.13), with $`\sigma _L(\mathrm{\Lambda })`$ replaced by $`\gamma (\widehat{x},\zeta ;\lambda )`$ therein, i.e.
$$\mathrm{det}[\gamma (\widehat{x},\zeta ;\lambda )+f(\widehat{x},\zeta ;\lambda )]0\lambda 𝒞𝐑_+,$$
$`(5.19)`$
which is satisfied if
$$\mathrm{det}[\mathrm{Re}^2f(\widehat{x},\zeta ;\lambda )\mathrm{Im}^2f(\widehat{x},\zeta ;\lambda )\gamma ^2(\widehat{x},\zeta ;\lambda )+2i\mathrm{Re}f(\widehat{x},\zeta ;\lambda )\mathrm{Im}f(\widehat{x},\zeta ;\lambda )]0.$$
$`(5.20)`$
6. Further applications
In the case of more mathematical interest where the operator in the interior of $`M`$ remains a Laplacian, while the boundary operator has a pseudo-differential sector, the analysis is much simpler. For example, two cases can be considered. (i) If $`\stackrel{~}{\mathrm{\Lambda }}`$ is a pseudo-differential operator of order $`1`$, the leading symbol of the boundary operator (3.9) can be cast in the form (cf. (4.7))
$$\sigma _L()=\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ i(T+\stackrel{~}{T})& I\mathrm{\Pi }\end{array}\right),$$
$`(6.1)`$
where $`T\mathrm{\Gamma }^j\zeta _j`$ and $`\stackrel{~}{T}`$ results from the occurrence of $`\stackrel{~}{\mathrm{\Lambda }}`$. The sufficient condition for finding solutions of Eq. (4.5’) for all $`\psi ^{}`$ reads now
$$(T+\stackrel{~}{T})^2+|\zeta |^2I>0\zeta 0,$$
$`(6.2)`$
because one can simply replace $`T`$ with $`T+\stackrel{~}{T}`$ in the analysis of Ref. , if Eq. (6.1) holds. Thus, if $`\stackrel{~}{\mathrm{\Lambda }}`$ is chosen in such a way that
$$(T+\stackrel{~}{T})^2(\left\{\widehat{x}\right\},\left\{\zeta \right\})>0\zeta 0,$$
$`(6.3)`$
Eq. (4.5’) can always be solved with arbitrary $`\psi ^{}(\phi )`$. The condition (6.3) can be made explicit after re-writing the DeWitt supermetric (3.6) in the more general form
$$E^{abcd}\frac{1}{2}(g^{ac}g^{bd}+g^{ad}g^{bc})+\alpha g^{ab}g^{cd}.$$
$`(6.4)`$
Thus, on defining (with $`e_a^i`$ being a local tangent frame on $`M`$)
$$\zeta _ae_a^j\zeta _j,$$
$`(6.5)`$
and introducing the nilpotent matrices
$$(p_1)_{ab}^{cd}N_aN_b\zeta ^{(c}N^{d)},$$
$`(6.6)`$
$$(p_2)_{ab}^{cd}N_{(a}\zeta _{b)}N^cN^d,$$
$`(6.7)`$
the work in Ref. finds the useful formula
$$T=\frac{1}{(1+\alpha )}p_1+p_2,$$
$`(6.8)`$
and this should be inserted into (6.3) to restrict the kernel of $`\stackrel{~}{\mathrm{\Lambda }}`$, whose leading symbol is equal to $`i\stackrel{~}{T}`$. The resulting restriction on $`\alpha `$ should be made compatible with the values of $`\alpha `$ for which the ellipticity condition (4.14) is fulfilled in the interior of $`M`$. From this point of view, one has definitely more choice than in the case of the local boundary operator (2.2) for an operator of Laplace type on metric perturbations, because the values of $`\alpha `$ for which the condition
$$T^2+|\zeta |^2I>0\zeta 0$$
$`(6.9)`$
holds (cf. (6.2)) are incompatible with the occurrence of an operator of Laplace type on metric perturbations . (ii) If $`\stackrel{~}{\mathrm{\Lambda }}`$ is a pseudo-differential operator of order $`d>1`$ (but not necessarily integer), the leading symbol of the boundary operator (3.9) can be expressed in the form
$$\sigma _L()=\left(\begin{array}{cc}\mathrm{\Pi }& 0\\ \widehat{T}& I\mathrm{\Pi }\end{array}\right).$$
$`(6.10)`$
The sufficient condition for finding solutions of Eq. (4.5’) reads instead
$$\widehat{T}^2+|\zeta |^2I>0\zeta 0.$$
$`(6.11)`$
It is therefore sufficient to choose $`\stackrel{~}{\mathrm{\Lambda }}`$ in such a way that
$$\widehat{T}^2<0\zeta 0.$$
$`(6.12)`$
7. Concluding remarks
The mathematical literature and, in particular, the work by Grubb , had already considered boundary conditions of the form (3.8), where $`(\mathrm{\Lambda }+\stackrel{~}{\mathrm{\Lambda }})`$ is allowed to be a pseudo-differential operator, but for elliptic differential operators. In physics, however, the requirement of gauge invariance of the boundary conditions for quantum gravity leads to an operator on metric perturbations (see (4.8)–(4.11)) which is itself pseudo-differential, since (3.8) is obtained from the vanishing of the gauge-averaging functional at the boundary (see (3.4)). Our physical problem remains therefore original with respect to the mathematical investigations . Our main contributions are as follows. (1) The projector condition for the boundary operator in Euclidean quantum gravity at one loop has been derived in the form (3.15). (2) A careful definition of strong ellipticity of pseudo-differential boundary-value problems in Euclidean quantum gravity has been proposed in section 4, with detailed physical applications in section 5 (see (5.7), (5.11), (5.13), (5.19) and (5.20)), and further mathematical examples in section 6. In other words, we have provided a complete characterization of the properties of the symbol of the boundary operator for which a set of boundary conditions completely invariant under infinitesimal diffeomorphisms are compatible with a strongly elliptic one-loop quantum theory. The analysis of section 5 is detailed but general, and hence has the merit (as far as we can see) of including all pseudo-differential boundary operators for which the sufficient conditions derived therein can be imposed.
It would be now very interesting to prove that, by virtue of the pseudo-differential nature of $``$ in (3.9), the quantum state of the universe in one-loop semiclassical theory can be made of surface-state type . This would describe a wave function of the universe with exponential decay away from the boundary, which might provide a novel description of quantum physics at the Planck length. It therefore seems that by insisting on path-integral quantization, strong ellipticity of the Euclidean theory and invariance principles, new deep perspectives are in sight. These are in turn closer to what we may hope to test, i.e. the one-loop semiclassical approximation in quantum gravity . In the seventies, such calculations could provide a guiding principle for selecting couplings of matter fields to gravity in a unified field theory . Now they can lead instead to a deeper understanding of the interplay between non-local formulations \[19–21\], elliptic theory , gauge-invariant quantization and a quantum theory of the very early universe .
Acknowledgment
This work has been partially supported by PRIN97 “Sintesi”. Correspondence and conversations with Gerd Brubb have been very helpful. The author is indebted to Ivan Avramidi for previous collaboration, and to Pietro Santorelli for encouragement.
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warning/0001/cond-mat0001071.html | ar5iv | text | # Generalized symmetric nonextensive thermostatistics and 𝑞-modified structures
## Abstract
We formulate a convenient generalization of the $`q`$-expectation value, based on the analogy of the symmetric quantum groups and $`q`$-calculus, and show that the $`qq^1`$ symmetric nonextensive entropy preserves all of the mathematical structure of thermodynamics just as in the case of non-symmetric Tsallis statistics. Basic properties and analogies with quantum groups are discussed.
PACS: 05.20.-y; 05.70.-a; 05.30.-d
In the last few years there has been much interest in nonextensive classical and quantum physics. The nonextensive statistical mechanics proposed by Tsallis , has been the source of inspiration for many investigations in systems which represent multifractal properties, long-range interactions and/or long-range memory effects . On the other hand, quantum groups and the derived $`q`$-deformed algebraic structure such as $`q`$-oscillators, based on the deformation of standard oscillator commutation-anticommutation relations, have created considerable interest in mathematical physics and in several applications . The investigations described above are only two apparently unrelated areas in nonextensive physics. Although a complete understanding of the connection between nonextensive statistics and $`q`$-deformed structure is still lacking, many papers are devoted to the study of a deep connection between these two non-extensive formalisms .
Tsallis statistics is a $`q`$-nonsymmetric formalism i.e., not invariant under $`qq^1`$. Recently Abe has employed a connection between Tsallis entropy and the non-symmetric Jackson derivative. Because the requirement of invariance under $`qq^1`$ is very important in quantum groups , the above connection allows him to extend the Tsallis entropy to the $`q`$-symmetric one by means of a symmetric Jackson derivative. However, Abe has not extended the definition of the expectation value of an observable to the symmetric case and thus unable to formulate the thermostatistics which will preserve the Legendre transformation of standard thermodynamics in contrast to the Tsallis statistics which does. We would like to point out that Ref. introduces a two-parameter modification for the entropy and for the expectation value of an observable but does not also produce a consistent formulation of thermostatistics.
The purpose of this letter is to show how the $`qq^1`$ symmetric generalization of the Tsallis entropy together with a natural generalization of the $`q`$-expectation value produces a thermostatistics that preserves the mathematical structure of standard thermodynamics and show that this property is a direct consequence of the generalization from the non-symmetric $`q`$-calculus to the symmetric one.
Before investigating the symmetric nonextensive thermostatistics, let us briefly review the fundamental properties of the Tsallis thermostatistics, which is based upon the following two postulates .
* A nonextensive generalization of the Boltzmann-Gibbs entropy (Boltzmann constant is set equal to unity)
$$S_q=\frac{1}{q1}\left(1\underset{i=1}{\overset{W}{}}p_i^q\right),\mathrm{with}\underset{i=1}{\overset{W}{}}p_i=1,$$
(1)
where $`p_i`$ is the probability of a given microstate among $`W`$ different ones and $`q`$ is a fixed real parameter. The new entropy has the usual properties of positivity, equiprobability, and reduces to the conventional Boltzmann–Gibbs entropy $`S=_ip_i\mathrm{ln}p_i`$ in the limit $`q1`$.
* A generalized definition of internal energy
$$U_q=\underset{i}{}ϵ_ip_i^q$$
(2)
and, accordingly, a generalization of the $`q`$-expectation value of an observable $`A`$ which can be expressed as $`A_q_iA_ip_i^q`$. In the limit $`q1`$, $`A_1`$ corresponds to the standard mean value. This postulate plays a central role in the derivation of the equilibrium distribution and leads to the correct thermodynamic relations.
The deformation parameter $`q`$ measures the degree of nonextensivity of the theory. In fact, if we have two independent systems $`A`$ e $`B`$, such that the probability of $`A+B`$ is factorized into $`p_{A+B}(u_A,u_B)=p_A(u_A)p_B(u_B)`$, the global entropy is not simply the sum of their entropies and it is easy to verify that
$$S_q(A+B)=S_q(A)+S_q(B)+(1q)S_q(A)S_q(B).$$
(3)
Another important property is that $`S_q`$ is consistent with Laplace’s maximum ignorance principle, i.e., if $`p_i=1/W,i`$ and $`W1`$ (equiprobability) one has the following extremum value
$$S_q=\mathrm{ln}_qW,$$
(4)
where we have defined the generalized logarithmic function $`\mathrm{ln}_qx`$ as
$$\mathrm{ln}_qx=\frac{x^{1q}1}{1q}.$$
(5)
In the limit $`q1`$, $`\mathrm{ln}_qx\mathrm{ln}x`$ and Eq.(4) reproduces Boltzmann’s celebrated formula $`S=\mathrm{ln}W`$. Let us stress the importance of this result for the purpose of the following discussion, because it defines the generalized logarithmic function in nonextensive statistics and plays a crucial role in the determination of the $`q`$-expectation value of an observable, consistent with the thermodynamic relations as we shall see in Eq.(8) below.
Working with the canonical ensemble, the probability distribution can be obtained by extremizing the entropy $`S_q`$ under fixed internal energy $`U_q`$ constraint and norm constraint ($`_ip_i=1`$). The outcome of this optimization procedure gives the result
$`p_i={\displaystyle \frac{\left[1(1q)\beta ϵ_i\right]^{\frac{1}{1q}}}{Z_q}},`$ (6)
where $`Z_q`$ is the partition function given by
$`Z_q`$ $`=`$ $`{\displaystyle \underset{i}{}}\left[1(1q)\beta ϵ_i\right]^{\frac{1}{1q}}.`$ (7)
Using the generalized expression (5) of the logarithmic function and Eqs.(2) and (7), it has been shown that
$$U_q=\frac{}{\beta }\mathrm{ln}_qZ_q.$$
(8)
On the basis of the above relation, the entire mathematical structure of the connection between standard statistical mechanics and thermodynamics is preserved by the generalization of the Tsallis entropy, the definition of the internal energy and replacing $`\mathrm{ln}Z`$ by $`\mathrm{ln}_qZ_q`$.
Recently Abe has observed the connection between Tsallis entropy and Jackson derivatives which can be expressed as
$$S_q=_x^{(q)}\underset{i}{}p_i^x|_{x=1}\underset{i}{}\frac{p_i^qp_i}{q1},$$
(9)
where
$$_x^{(q)}f(x)=\frac{f(qx)f(x)}{x(q1)},$$
(10)
is the Jackson derivative , which in the limit $`q1`$, becomes the ordinary differential. The above connection is not just a coincidence but in fact it has been shown that the Jackson derivative can be identified with the generators of fractal and multifractal sets with discrete dilatation symmetries and thus it is strictly related to Tsallis statistics.
The Jackson derivative in Eq.(10) is intimately connected with $`q`$-deformed structures in $`q`$-oscillator theory, signified by the $`q`$-basic number
$$[x]_q=\frac{q^x1}{q1}.$$
(11)
It has been shown that the pseudo-additivity property of the Tsallis entropy, displayed in Eq.(3), is also valid for the above $`q`$-basic number .
We now develop the $`q`$-symmetric theory of the Tsallis thermostatistics. In $`q`$-deformed structures, when one constructs the theory which is invariant under $`qq^1`$, the Jackson derivative has to be generalized to the form
$$𝒟_x^{(q)}f(x)=\frac{f(qx)f(q^1x)}{x(qq^1)},$$
(12)
and correspondingly, it is possible to introduce the symmetric Tsallis entropy given by
$$S_q^S=𝒟_x^{(q)}\underset{i}{}p_i^x|_{x=1}=\underset{i}{}\frac{p_i^qp_i^{q^1}}{qq^1}.$$
(13)
The above expression for the symmetric Tsallis entropy satisfies a generalized pseudo-additivity property formally similar to the symmetric $`q`$-basic number defined by
$$[x]_q^S=\frac{q^xq^x}{qq^1}.$$
(14)
It is easy to show that the $`q`$-symmetric Tsallis entropy can be written in terms of the nonsymmetric Tsallis entropy in the compact form
$$S_q^S=c_1S_q+c_2S_{q^1},$$
(15)
where $`c_1`$ and $`c_2`$ are two coefficients (always positive) which depend only on $`q`$ and $`q^1`$,
$$c_1=\frac{q1}{qq^1},c_2=\frac{1q^1}{qq^1},c_1+c_2=1.$$
(16)
We wish to stress that the above expression is a consequence of the connection between the $`q`$-symmetric and the nonsymmetric $`q`$-structures. In fact we observe that the Jackson symmetric derivative in Eq.(12) can be expressed in terms of the nonsymmetric ones
$$𝒟_x^{(q)}=c_1_x^{(q)}+c_2_x^{(q^1)}.$$
(17)
An analogous relation also holds for the $`q`$-basic number in quantum groups
$$[x]_q^S=c_1[x]_q+c_2[x]_{q^1}.$$
(18)
Accordingly the above relations provide us with a recipe to obtain the symmetric generalization from the nonsymmetric structures.
As in the case of Tsallis entropy, the equiprobability distribution ($`p_i=1/W,i`$) can be derived by employing a new symmetric definition of the logarithmic function
$$S_q^S=\mathrm{ln}_q^SW=c_1\mathrm{ln}_qW+c_2\mathrm{ln}_{q^1}W.$$
(19)
The above result is very important for the correct construction of a generalized symmetric thermostatistics. In fact this allows us to obtain the symmetric logarithm of the partition function (thermodynamic potential) in terms of a linear combination of the Tsallis one
$$\mathrm{ln}_q^SZ_q=c_1\mathrm{ln}_qZ_q+c_2\mathrm{ln}_{q^1}Z_{q^1}.$$
(20)
Because Tsallis statistics satisfies the condition in Eq.(8), it is also verified immediately in symmetric nonextensive statistics that
$$U_q^S=\frac{}{\beta }\mathrm{ln}_q^SZ_q,$$
(21)
if we choose the symmetric internal energy to be in the form
$$U_q^S=c_1U_q+c_2U_{q^1}\underset{i}{}ϵ_i\frac{(q1)p_i^q(q^11)p_i^{q^1}}{qq^1}.$$
(22)
The definition in Eq.(22) of the internal energy implies a generalized symmetric $`q`$-expectation value of a physical observable $`A`$
$$A_q^S=c_1A_q+c_2A_{q^1}\underset{i}{}A_i\frac{(q1)p_i^q(q^11)p_i^{q^1}}{qq^1},$$
(23)
where we note that the second relation on the right hand side is true only if $`A`$ does not depend on $`q`$ (such as the case of energy or particle number).
The above results are very important because we have a new generalized definition of the internal energy, which together with the symmetric Tsallis entropy, preserves all the thermodynamic relations (Legendre transformations). This follows directly from the $`qq^1`$ invariance of the $`q`$-deformed algebra, thus offering a closer connection between nonextensive Tsallis statistics and $`q`$-deformed structures.
Following the standard procedure, the probability distribution can be obtained by extremizing the entropy $`S_q^S`$ under fixed internal energy $`U_q^S`$ constraint and the norm constraint ($`_ip_i=1`$). The result can be written as a linear combination
$$p_i^S=c_1\frac{\left[1(1q)\beta ϵ_i\right]^{\frac{1}{1q}}}{Z_q}+c_2\frac{\left[1(1q^1)\beta ϵ_i\right]^{\frac{1}{1q^1}}}{Z_{q^1}}.$$
(24)
In the limit $`q1`$, $`p_i^S`$ reduces to the standard Maxwell-Boltzmann distribution. We note that the extremization procedure only establishes that the solution for the distribution function is a linear combination of Tsallis distribution evaluated at $`q`$ and $`q^1`$. Eq.(24) is a reasonable choice and follows the prescription of the $`q`$-calculus. In fact in the $`q`$-oscillator theory the statistical distribution function can be written as the same linear combination of the non-symmetric ones
$$f_q^S=c_1f_q+c_2f_{q^1},$$
(25)
where $`f_q`$ and $`f_{q^1}`$ are the distribution functions in non-symmetric $`q`$-boson oscillators
$$f_q=\frac{1}{e^{\beta \omega }q}.$$
(26)
In Fig. 1, we show the plot of the normalized probability function (24) against $`\beta ϵ`$ for different values of $`q`$. Let us note that the above distribution has no cut-off as in Tsallis’s distribution for $`q<1`$ and the high energy tail of the distribution is always enhanced compared to the Maxwell-Boltzmann distribution since $`p_i^S`$ has the following power law behavior at high energy, $`p_i^S=aE^{1/(1q)}+bE^{1/(1q^1)}`$.
In light of the above discussion, making a Legendre transform of the function $`\mathrm{ln}_q^SZ_q`$ it is easy to verify the validity of the relation
$$S_q^S=\beta U_q^S+\mathrm{ln}_q^SZ_q,$$
(27)
which implies the standard thermodynamic relation
$$\frac{S_q^S}{U_q^S}=\frac{1}{T}$$
(28)
and the $`q`$-deformed symmetric free energy is given by
$$F_q^S=\frac{1}{\beta }\mathrm{ln}_q^SZ_q=U_q^STS_q^S.$$
(29)
All the above equations reduce to the standard thermodynamic relations in the limit $`q1`$.
Finally, we note that Tsallis recently introduced a normalization procedure for the $`q`$-expectation value of an observable in order to remove some anomalies, such as non-additivity of the generalized internal energy and non-invariance of the probability distribution under the choice of origin of the energy spectrum. In the framework of the new generalization, all the results of the present investigation remain unaltered if we implement the normalization according to
$$\stackrel{~}{U}_q^S=c_1\stackrel{~}{U}_q+c_2\stackrel{~}{U}_{q^1},$$
(30)
where $`\stackrel{~}{U}_q`$ is the normalized Tsallis internal energy
$$\stackrel{~}{U}_q=\frac{_i\stackrel{~}{p}_i^qϵ_i}{_i\stackrel{~}{p}_i^q},$$
(31)
and $`\stackrel{~}{p}_i`$ is the modified Tsallis distribution in the normalized $`q`$-expectation value given by
$`\stackrel{~}{p}_i={\displaystyle \frac{\left[1(1q)\beta (ϵ_i\stackrel{~}{U}_q)/_j\stackrel{~}{p}_j^q\right]^{\frac{1}{1q}}}{\stackrel{~}{Z}_q}},`$ (32)
with
$`\stackrel{~}{Z}_q={\displaystyle \underset{i}{}}\left[1(1q)\beta (ϵ_i\stackrel{~}{U}_q)/{\displaystyle \underset{j}{}}\stackrel{~}{p}_j^q\right]^{\frac{1}{1q}}.`$ (33)
In summary, we have shown that it is possible to extend Tsallis thermostatistics to the $`qq^1`$ symmetric generalization which preserves the Legendre transformation of standard thermodynamics. This is achieved by introducing a $`q`$-symmetric expectation value which follows directly from the extension of the nonsymmetric $`q`$-deformed theory to the symmetric one. We thus establish a closer connection between nonextensive Tsallis statistics and $`q`$-deformed structures.
In conclusion, the relevance of the $`qq^1`$ symmetry is well-known in $`q`$-oscillators from the mathematical structure as well as in applications . Similarly we expect the symmetric Tsallis thermostatistics to be useful in many future investigations.
Acknowledgments
We are grateful to P. Quarati and C. Tsallis for reading this manuscript and for useful suggestions. One of us (A.L.) would like to thank the Physics Department of Southern Illinois University for warm hospitality where this work was done. |
warning/0001/cond-mat0001333.html | ar5iv | text | # The interacting system of electrons, positrons and photons in high external electric and arbitrary magnetic fields
## 1 Introduction
In earlier my publications made the theoretical investigation for the interacting system of electrons and phonons in semiconductors, semimetals and gaseous plasmas at high external electric and magnetic fields also under the conditions of the propagation of strong electromagnetic waves. In \[1-4\] the connected system of kinetic equations of interacting electrons, holes and phonons at high electric and magnetic fields was solved, with taking into account arbitrary heating and the mutual drag of carriersand phonons. For the phonons the solution of non-stationary kinetic equation was found and was shown that the non-equilibrium and non-stationary distribution function of phonons has the stationary limit, when the drift velocity of connected by the mutual drag system of carriers and phonons is less than the velocity of sound $`(u<s)`$. For the drift velocities $`u>s`$ the distribution function of phonons becomes is exponentially grown by the time. It accords to the generation of intrinsic phonons and amplification of phonons introduced to the system externally (the stream of phonons).
It was shown that the mutual drag leads to the re-normalization the mass of carriers. As a result of the mutual drag electrons and holes are ”dressed” by phonons and formed the ”quasi-particles” which has the electrons or holes charge $`(\pm e)`$ and the phonons mass $`m=T_i/s^2`$ (where $`T_i`$ \- is the temperature of the coupled by the mutual drag system of carriers and phonons) \[1-5\]. In weak electric fields the mass of phonons are $`m_o=T/s^2`$. In the case of propagation of strong electromagnetic wave in semiconductors and semimetals at external magnetic field this lead to the cyclotron resonance on such ”quasi-particles” with the frequencies $`\omega _H`$$`=(3/4)eH/\left[(T_i/s^2)c\right]`$ . It was shown in , that the effect of connection carriers and phonons as a result of their mutual drag is a common phenomenon. Really, as it follows from under the conditions of longitudinal propagation of strong electromagnetic wave in semimetals with equal electron and hole concentrations the frequencies of Alfven and magneto-sound waves are identically equal to each other and have the form:
$$\omega _A=\omega _{ms}=kH/\left(4mN\right)^{1/2}=kv_A,$$
where $`m_i=4(T_i/s^2)/3`$ also is so-called ” mass of phonons”, $`T_i=T_i(+0)`$ \- the effective temperature of the connected by the mutual drag system of ”electron + phonons” or ”hole + phonons”. Thus we may conclude that the mutual drag of carriers and phonons at external fields lead to the formation of compound particles (”quasi-particles”) - ”carriers dressed by phonons” with the joint drift velocity u.
In the first investigation the considerations was made for the degenerate semiconductors and semimetals under the conditions of the propagation of weak electromagnetic waves. In particular in was considered the propagation of weak electromagnetic waves on degenerate semiconductors with one type of carriers and was predicted the cyclotron resonance on the connected system of electron + phonon. But soon it was shown in that this resonance does not observable because of the presence of impurities with the concentrations $`N_i=n`$ (where $`n`$ is the concentration of electrons). Also in was shown that in pure semimetals it is necessary to take into account the presence of two type of carriers which drifted on the opposite directions at $`H=0`$. The remarks about non-observably in applies also to the uniform and non-uniform low frequency cyclotron resonance in degenerate semiconductors and semimetals with one type of carriers under the conditions of propagation strong electromagnetic waves considered in and . These questions was discussed in and was shown that in intrinsic semiconductors and semimetals realized both uniform and non-uniform cyclotron resonance on quasi-particle ”hole + phonon”. In also was shown that the such type of resonance realized on non-degenerate impurity and intrinsic semiconductors and discussed the question about their observably.
## 2 General
In present paper the connected system of kinetic equations for interacting system of electrons, positrons and photons in external high electric $`\stackrel{}{E}`$ and arbitrary magnetic $`\stackrel{}{H}`$ fields are solved. The non-equilibrium distribution functions of electrons, positrons and photons are founded by the taking into account of the arbitrary heating and mutual drag.
For the photons the general solution of non-stationary and non-uniform Boltzmann equations was founded. The cases of week $`(\omega _H^\pm \tau _c)1`$ , classically high $`(\omega _H^\pm \tau _c)1`$, and quantizing $`(\mathrm{}\omega _H^\pm T,T_c)`$ high magnetic fields are considered. Here $`T`$\- is the initial temperature of equilibrium system (at the $`t=0`$, before the external fields is applied), $`T_c(\stackrel{}{E},\stackrel{}{H})`$ -is the temperature of heated electrons and positrons, $`\omega _H^\pm `$ \- is their cyclotron frequencies, $`\tau _c^1=\nu _c`$ the relaxation frequencies of electrons and positrons( the photon production by the annihilation electron-positron’s pair or by the scattering of electrons and positrons by photons). $`c=e,p`$ means electrons and positrons, accordingly.
In the absence of or in a weak magnetic fields the inter-carrier collisions frequencies $`\nu _{ee}`$ and $`\nu _{pp}`$ assumed much more than others, and that is why the isotropic parts of the distribution functions of carriers assumed to be equilibrium one, with effective temperatures of carriers $`T_c=T_{e,p}(E,H)`$. This approximation corresponds to the case of high concentration of carriers $`n>n_k`$.
If the external field is a field of strong electromagnetic wave, then it is necessary to fulfil of the condition:
$$\omega \nu _ϵ^\pm $$
(1)
$`\nu _ϵ^\pm `$\- is the collusion frequencies of carriers for energy transfer to scatters, $`\omega `$ \- is the frequency of electromagnetic wave.
Under the conditions (1) the isotropic parts of carriers distribution functions at zero approximation by $`\nu _ϵ^\pm /\omega `$ do not depend on time directly. By the direct solution of quantum kinetic equations in common case for the arbitrary spherical symmetric dispersion low of carriers, it was shown that at quantizing and classically high magnetic fields the stationary distribution functions of carriers, satisfying the boundary conditions, $`F^\pm (ϵ)_ϵ\mathrm{}=0`$ has the form (for the arbitrary quantities of their concentrations):
$$F^\pm (ϵ)=\left\{c^1\mathrm{exp}\left(𝑑ϵ^\mathrm{`}/T_c(ϵ^\mathrm{`},t)+1\right)\right\}^1$$
(2)
There $`T_c(ϵ,t)=A(ϵ,t)/B(ϵ)`$ \- the temperature of carriers which have occupied the energetic level $`ϵ`$,
$$A(ϵ)=(2\pi /\mathrm{})\underset{\alpha \beta q}{}C_q^2I_{\alpha \beta }^2\left(\mathrm{}\omega _q^{}\right)^2N(q,t)\delta \left(ϵ_\beta ϵ_\alpha \mathrm{}\omega _q\right)\delta \left(ϵ_\alpha ϵ\right)$$
$$B(ϵ)=(2\pi /\mathrm{})\underset{\alpha \beta q}{}C_q^2I_{\alpha \beta }^2\mathrm{}\omega _q\delta \left(ϵ_\beta ϵ_\alpha \mathrm{}\omega _q\right)\delta \left(ϵ_\alpha ϵ\right)\text{}\omega _H^{}=\mathrm{}\omega _q\stackrel{}{V}\stackrel{}{q}$$
$`C_q`$\- being the constant of interaction and $`I_{\alpha \beta }`$ \- is a matrix elements for the transition from state $`\alpha `$ to $`\beta `$ and back (reverse).
For the arbitrary degree of quantization we have:
$$F^\pm (ϵ)=\{1+\mathrm{exp}\left[(ϵ\zeta (E,H)/T_c]\right\}^1$$
$$T_c=T_i\{1+(\frac{V^\pm }{u})1]2(f11)\}\}$$
(3)
$$\phi _1=\left[1\frac{u^2}{c^2}\right]^{1/2}$$
At the classical region of strong magnetic field we have:
$$T_c=T_i\left\{1+\frac{1}{3}\left(\frac{V^\pm }{c}\right)^2+\left[1\frac{V^\pm }{u}\right]\left(\phi _21\right)\right\},\text{ }\phi _2=\frac{c}{2u}\mathrm{ln}\left|\frac{c+u}{cu}\right|\text{ (3’)}$$
Here $`V^\pm =cE/H`$ is the Hall’s drift velocity of carriers. It was shown that for all considered cases the solution of non-stationary kinetic equation for the photons is:
$$N(\stackrel{}{q},t)=\{N(\stackrel{}{q},\stackrel{}{r}\stackrel{}{u_o}t,0)+\beta \underset{0}{\overset{t}{}}N(q,\tau \mathrm{`})exp(\underset{0}{\overset{\tau }{}}\gamma _q\left(\tau \mathrm{`}\right)d\tau \mathrm{`})\}\times $$
$$\times exp\left[\underset{0}{\overset{t}{}}\gamma _q\left(\tau \right)𝑑\tau \right]$$
(4)
Where $`N(\stackrel{}{q},\stackrel{}{r}\stackrel{}{u_o}t,0)`$ the distribution function of photons in the absence of electric and magnetic fields (at $`t=0`$), which in the case of space uniformity is equilibrium Plank’s function at the temperature $`T`$. The increasing increment of photons is
$$\gamma _q=\beta \left[\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}1\right]$$
(5)
$`\beta =\left(\beta _e+\beta _p+\beta _{ph}+\beta \right)`$is a total collisions frequency of photons with electrons $`(e)`$, positrons $`(p)`$ (including the photon decay to electron-positron pair), photons $`(ph)`$ and boundaries $`(b)`$ of region occupied by the system, if such one exists.
$$\stackrel{}{u}(t)=\underset{\pm }{}\stackrel{}{u^\pm }(t)=\frac{\beta _e}{\beta }\stackrel{}{V^{}}(t)+\frac{\beta _p}{\beta }\stackrel{}{V^+}(t)$$
(6)
$`\stackrel{}{V^\pm }(t)`$ is average drift velocity of carriers, $`\stackrel{}{u^{}}(t)`$ \- the drift velocity of connected by the mutual drag system of ”electron + positrons” and $`\stackrel{}{u^+}(t)`$ is the same for the ”positron + photons”.
In the common case, when the heating of carriers was realized by the field of strong electromagnetic wave $`\stackrel{}{E}=\stackrel{}{E_o}e^{i\omega t}+\stackrel{}{E_o^{}}e^{+i\omega t}`$, $`\stackrel{}{V^\pm }(t)=\stackrel{}{V^\pm }\mathrm{cos}\omega t`$ we have:
$$N(\stackrel{}{q},t)=\{N(q,o)+\beta \underset{0}{\overset{t}{}}d\tau N(q,\tau )\mathrm{exp}\left[\beta (\tau \frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}\frac{\mathrm{sin}\omega \tau }{\omega })\right]\}\times $$
$$\times exp\left\{\beta \left[t\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}\frac{\mathrm{sin}\omega t}{\omega }\right]\right\}$$
(7)
In the case of constant external electric field $`(\omega 0)`$ we have:
$$N(\stackrel{}{q},t)=\left\{\frac{N(q,o)N(q,T_i)}{1\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q}\right\}exp\left\{\beta \left[\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}1\right]t\right\}+\frac{N(q,T_i)}{1\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q}$$
(8)
Here $`T_i=\left(\beta _c/\beta \right)T_c+\left(\beta _{ph}/\beta \right)T_{ph}+\left(\beta _b/\beta \right)T_b`$ is the temperature of the coupled by the mutual drag system of heated complexes of carriers and photons.
We have still considered the case when the initial state of photons (at $`t=0`$) was assumed to be equilibrium state without distinguished direction. If the part of initial distribution of photons has directional drift (the photons stream), then the kinetic equation for the photons has the form:
$$\frac{N(q,\stackrel{}{r},t)}{t}+\frac{N(q,\stackrel{}{r},t)}{\stackrel{}{r}}\frac{d\stackrel{}{r}}{dt}=\beta \left\{N(q,t)\left(1\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}\right)\right\}+N_i(q,T_i)$$
(9)
By the single substitution $`t\mathrm{`}=t`$ and $`\stackrel{}{r\mathrm{`}}=\stackrel{}{r}\stackrel{}{u_o}t`$ the solution (10) may be transformed into the form:
$$N(\stackrel{}{q},\stackrel{}{r},t)=\{N(\stackrel{}{q},\stackrel{}{r}\stackrel{}{u_o}t,0)+\beta N_i(q,T_i)\underset{0}{\overset{t}{}}exp\left[\beta \underset{0}{\overset{\tau }{}}(1\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q})d\tau \mathrm{`}\right]d\tau \}\times $$
$$\times exp\left[\beta \underset{0}{\overset{\tau }{}}\left(1\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}\right)𝑑\tau \right]\text{ (9’)}$$
Thus it is seemed that the solution of the uniform equation for the photons (8) and non-uniform equation (9’) have the same form corresponding to the different initial conditions. Therefore, if the initial distribution function of photons has the form of the drifted Plank distribution function $`N(q,\stackrel{}{r}\stackrel{}{u_o}t,0)`$ and if the external fields is uniform, then this non-uniformity has to be served with time and the drift at external field is to be added to them. Thus the equation (9’) allows to consider processes of absorption or amplification of photons, introduced to the system from outside (initial stream of photons), and the generation of own photons of system in external fields. In principle it is the most common form of the initial distribution function, which is taking the chance for examination of the affirmation of the special theory of relativity about equivalency of all inertial frames of reference. Really by the transition to the frame of reference drifting jointly with photons, as a result, we have the Plank’s equilibrium distribution function at the temperature $`T`$ in this frame of reference. In the other words, in the absence of external fields $`(E=H=0)`$ for the initial non-uniform system of photons (9) from kinetic equation we receive the uniform one, by the transition from one frame of reference to an other but it is not means that the two frames of reference is equivalent.
In fact, by the transition from the frame of reference drifting jointly with the photons by the velocity$`\stackrel{}{u_o}=c\stackrel{}{q}/q`$ to the frame of reference which in the rest which is equivalent to the transition from one internal frames of reference $`(u=u_o)`$ to the other $`(u=0)`$ we receive $`N(\stackrel{}{q},\stackrel{}{r}\stackrel{}{u_o}t,0)=N(q,0)=N_o(q,T)`$, for all moments of time, that is the solution (9’) transform to (8), but it is not means that this two frame of reference is equivalent. On the other words the demand of equivalency of the laws of physics in that two inertial frame of reference is equivalent to the demand of the equivalency the equilibrium Plank’s distribution function to the drifted Plank’s distribution function or to the demand of equivalency the laws of physics in the uniform and non-uniform cases (or spaces).
As one can see from (4), (8) and (9’) the general solution of non-stationary equation of photons have the stationary limit in the region of drift velocities $`(\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q)<1`$
$$\underset{t\mathrm{}}{lim}N(\stackrel{}{q},t)=N(\stackrel{}{q})=N(q,Ti)\left(1\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}\right)$$
(10)
As it seems from equations (4), (8) and (9’) for the drift velocities $`(\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q)>1`$ the distribution function of photons becomes exponentially grown by the time. It is accords to the generation of intrinsic photons by the increment of grow $`\gamma _q`$ and amplification of photons introduced to the system from outside (stream of photons) by the coefficient of amplification
$$\mathrm{\Gamma }_q=\frac{\gamma _q}{c}=\frac{\beta }{c}\left(\frac{\stackrel{}{u}\stackrel{}{q}}{\mathrm{}\omega _q}1\right)=\frac{\beta }{c}\left[\frac{u}{c}\mathrm{cos}a1\right]$$
(11)
$`\alpha `$$`=\left(\stackrel{}{u}\mathrm{^}\stackrel{}{q}\right)`$ \- the angle between the drift velocity of connected by the mutual drag system ”electron + photons” (”the dressed electron”) or ” positron + photons” (”the dressed positron”) and momentum of photon. The expression for the electrical current of electrons and positrons at classically high external magnetic fields has the form:
In the case of the propagation of strong electromagnetic wave and at the presence of external magnetic field the current of electrons and positrons has the form:
$`\stackrel{}{j_\pm }=ne\stackrel{}{V^\pm },`$ $`\stackrel{}{V^\pm }=\stackrel{}{V^\pm }(ϵ)`$ \- is the averaged drift velocity of carriers. Here
$$\stackrel{}{V^\pm }\left(ϵ\right)=\frac{e\mathrm{\Omega }_\pm (ϵ)}{m_c}\frac{\stackrel{}{E_{}}(\omega _H^\pm /\mathrm{\Omega }_\pm \left(ϵ\right))\left[\stackrel{}{h}\stackrel{}{E}\right](\omega _H^\pm /\mathrm{\Omega }_\pm \left(ϵ\right))^2\stackrel{}{h}\left(\stackrel{}{h}\stackrel{}{E}\right)}{\mathrm{\Omega }_\pm ^2\left(\omega _H^\pm \right)^2},\text{ }\stackrel{}{h}=\stackrel{}{H}/H\text{ }$$
(12)
In the case of $`\omega 0,`$ $`\stackrel{}{E}\stackrel{}{H}`$ (or $`H=0`$):
$$\stackrel{}{V^\pm }\left(ϵ\right)=\frac{e\stackrel{}{E}\beta _c}{m_c\nu _{ph}(ϵ,u)\beta _{ph,b}}=\frac{e\stackrel{}{E}}{m(T_i,u)\beta _{ph,b}}$$
(13)
Here $`m=m_c\nu _{ph}(ϵ,u)/\beta _c`$-the mass of connected by the mutual drag system of carriers and photons and $`\beta _{ph,b}=\beta _{ph}+\beta _b`$.
Really the interacting system of electrons, positrons and photons at the external high electromagnetic and the classically high or the quantizing magnetic fields under the conditions of their heating, mutual drags and at the stationary conditions $`(\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q)<1`$ has the cyclotron resonance with the frequencies
$$\omega _H^\pm =\frac{eH}{\left(T_i/c^2\right)c}=\frac{eH}{m(T_i)c}$$
(14)
As it seems from the equation (12) the resonance is taking place on frequencies of electromagnetic wave less than the collision frequencies of photons with carriers. The width of the resonance lines defined by the expression $`\gamma `$$`=(3/2)\left[\omega ^2/\beta _c+\beta _{ph}+\beta _b\right]`$.
In the other words, because of the mutual drag, the electrons and positrons turn into the compound particles (into the coupled system of ”electron + photons” and ”positrons + photons” i.e. so called ”dressed” by the photons ”quasi-electron” or ”quasi-positron” with the effective mass $`m(T_i)`$. In fact, we receive the quasi-particle with the electron’s or positrons charge and the photons mass:
$$m(T_i)=m(E,H)=T_i/c^2\text{ or }E=T_i=m(T_i)c^2$$
(15)
Since $`T_i`$ and $`T`$ means the average kinetic energy, then from the equation (10), (11) and (14) we receive that the, so-called ”velocity of light” in vacuum $`\mathrm{"}c\mathrm{"}`$ is the average velocity of photons in the equilibrium or stationary state with the temperatures $`T`$. The ”dressing”of electrons and positrons in quantum electrodynamics was connected with the virtual absorption and emission of photons by the electrons and positrons, which is occupied the given level of energy, i.e. with the finiteness of the lifetime of electrons and positrons or the natural width of the given energetic level. For the interacting system of electrons, holes and phonons in semiconductors, semimetals and gaseous plasmas the analogous problem was solved in \[1- 5\].
As it seems from (4), (8), (9’) and (10) at the external electric and magnetic fields under the stationary conditions the relativistic factor enters to the distribution function of photons in first order in the form $`[1(\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q)]^1`$, instead of the form $`[1(v^2/c^2)]^{1/2}`$ in relativistic electrodynamics. This is connected with the violation of $`T`$\- symmetry $`(tt)`$ of equations in electrodynamics and, in common dynamics at the external fields. In this case the uniformity of the space and as a result the law of the conservation of momentum is violated too. Since the external fields acts constantly but not instantaneously i.e. we have the motion with acceleration and the oscillatory regime is absent. The substitution $`tt`$ do not simply lead to substitution $`\stackrel{}{v}\stackrel{}{v}`$, because of the motions along and opposite of field direction are differs and do not compensate each other.
The relativistic factor of the type $`[1(v^2/c^2)]^1`$ may appears only at the absence of external field in equilibrium or stationary conditions, by the using the isotropic part of the distribution function of photons. Factually, by the separation of thestationary distribution function of photons (10) to the isotropic and anisotropy parts we have:
$$N(q)=N_s(q)+N_\alpha (q)=N(q,T_i)\left[1\frac{u^2}{c^2}\mathrm{cos}{}_{}{}^{2}a\right]^1+N(q,T_i)\frac{u}{c}\mathrm{cos}\alpha \left[1\frac{u^2}{c^2}\mathrm{cos}{}_{}{}^{2}a\right]^1$$
(16)
Since as a result of the mutual drag the carriers and photons form the connected system (complex) with the common drift velocity, then under the conditions of strong (full) mutual drag $`\alpha =0`$ or $`\pi `$ and we received:
$$N_s(q)=N(q,T_i)\left[1\frac{u^2}{c^2}\right]1;\text{ }N_a(q)=\left(\frac{u}{c}\mathrm{cos}\alpha \right)N(q,T_i)\left[1\frac{u^2}{c^2}\right]^1$$
(17)
In the absence of external electric and magnetic fields, in common case $`u=u_o=const`$ and we have:
$$N(\stackrel{}{q},\stackrel{}{u_o})=N(\stackrel{}{q},\stackrel{}{r}\stackrel{}{u_o}t,0)=\{exp\left(\frac{\mathrm{}\omega _q^{}}{T}\right)1\}^1$$
(18)
$`\mathrm{}\omega _q`$$`=\mathrm{}\omega _q\stackrel{}{u_o}\stackrel{}{q}`$. By the transition to the frame of reference drifting together with photons we can received (16) the equilibrium Plank’s distribution function with the temperature $`T`$. Let us discuss the main question now: may the presence of the relativistic factor of first or secondary order lead to any singularities in physical phenomena or quantities? As it seems from the non-stationary solution for the distribution function of photons (8) or (9’) the Lorenz-Einstein theory corresponds to uniform (equilibrium) case and must satisfy the stationary condition $`v/c<1!`$ The case of $`v=c`$ is not included to their theory and for this reason the conclusions of the Einstein theory about equality the rest mass of photons to zero and $`c`$ \- is the ultimate velocity of propagation of all types of interaction in nature do not have the real basis.
Really, from general solution of the non-stationary kinetic equation for the photons (8), by the dividing the exponent to the series near the point $`\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q=1`$ we have
$$N(q,t)=\left\{N_o(q,T)+\frac{\beta N(q,T_i)}{\gamma _q}\right\}\left\{1+\gamma _qt+\frac{1}{2}(\gamma _qt)^2+\mathrm{}\right\}\frac{\beta N(q,T_i)}{\gamma _q}=$$
$$=\left\{N_o(q,T)(1+\gamma _qt)+N(q,T_i)\beta t\right\}+\frac{1}{2}\left\{N_o(q,T)(\gamma _qt)2+\beta \gamma _qN(q,T_i)t^2\right\}$$
(19)
In the point $`\gamma _q=0`$ we have
$$N(q,0)=\underset{\gamma _q0}{lim}N(q,t)=N_o(q,t)+N(q,T_i)\beta t$$
(20)
As it is shown from this expression at the point $`\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q=1`$, i.e. at the point $`u=c`$ the distribution function of photons is non-stationary and grows linearly by the time. What about of the singularity it is abbreviated clearly!
Since the Einstein theory is a stationary one it did not applied to the non-stationary conditions, namely to the region of the drift velocities v = c or u = c. For the drift velocities $`vc`$ or $`uc`$ the theory must be non-stationary. The Einstein theory is a one mode theory and that is why must be received from the many particle (or many mode) theory by the limiting transition to the one mode case. The effect of connection charged carriers and photons as a result of their mutual drag is a common phenomenon. It is a reaction of the system on action of external fields for the conservation the stationary state of system (the analog of ”self-conservation” in biology). Thus we are found the ”dressing” mechanism of charge carriers used earlier in quantum electrodynamics.
## 3 Conclusions
As a result of analysis the phenomena of connected system of charge carriers and photons by the using of equations (4), (7) - (10) and (15) we receive the following conclusions:
1. As it seems from non-stationary distribution function of photons the Lorenz-Einstein theory corresponds to space uniform (equilibrium) case and must satisfy the stationary condition $`v<c`$ ! The case of $`v=c`$ is not included in their theory. In the point $`u=c`$ (i.e. $`\gamma _q=0`$) the distribution function of photons (18) is non-stationary do not content any singularity and grows linearly by the time.
For this reason the conclusions of Einstein theory that the rest mass of photons is equal to zero and $`c`$ \- is the ultimate velocity of propagation of all types of interactions in nature do not have the real basis.
2. Since the Einstein theory was a one mode theory and that is why must be received from the many body (mode) theory by the limiting transition to the one mode case at $`v<c`$. For the $`v>c`$ or ($`u>c`$) the theory must be nonstationary. As we say the mutual drag lead to the formation the ”quasi-particles” - the electron or positron ”dressed by photons”. The average energy in stationary state $`u<c`$ for the one mode case is:
$$ϵ=\mathrm{}\omega N_i(q,T_i)=\frac{T_iN(\omega ,T_i)}{1u^2/c^2}=$$
$$=\frac{TN(\omega ,T_i)}{1u^2/c^2}\frac{T_i}{T}=\frac{M_oc^2}{1u^2/c^2}\left(\frac{T_i}{T}\right)^2$$
(21)
The mass of the heated photons for one mode:
$$M_i=\frac{M_o}{1u^2/c^2}\left(\frac{T_i}{T}\right)^2$$
The mass of one heated photon:
$$m=\frac{M_i}{N(\omega ,T_i)}=m_o\frac{T_i}{T}\left(1\frac{u^2}{c^2}\right)^1\text{ (21’)}$$
At the absence of heating
$$ϵ=Mc^2=N_o(\omega ,T)\frac{T}{1u^2/c^2}=M_oc^2\frac{T}{1u^2/c^2}=$$
$$=m_oc^2\frac{N_o(\omega ,T)}{1u^2/c^2}\text{}m=\frac{m_o}{1u^2/c^2}\text{,}$$
(22)
$`m_o=M_o/N_o(\omega ,T)=T/c^2`$ is the rest mass of photon, i.e. the mass of photon in frame of reference which drifted jointly with photons at the temperature $`T`$, $`N`$ is the concentration of photons for one mode.
3. As it seems from (19) at high electric and magnetic fields for the drift velocities $`u<c`$ the energy (or the mass) of photons for one mode grows as a result of the mutual drag, as well as the heating of the carriers and photons.
4. As it seems from (4)-(10) at the external electric and magnetic fields the relativistic factor enters to the expressions of the distribution function and other physical quantities in first order in the form $`\left(1u^2/c^2\right)^1`$, instead the $`\left(1v^2/c^2\right)^{1/2}`$ in Einstein theory. This is together with the conclusions 1 and 2 is solve the main problem of super- luminal particles named -tachyons, because of in our theory the imaginaryty of mass of tachyons is liquidated.
5. There is the opinion that the original conception about tachyons, as individual particles such as the electrons, protons and etc. is not correct and the tachyons in such understanding is absent . Our investigations shows that the ordinary particles such as the electrons, positrons and also photons may stand a super-luminal in high external fields under the conditions of mutual drag. Also there are the opinion that the tachyons as an elementary excitations (quasi-particles) have the wide - spread in complex systems which is lose the stability and made the phase transition to the stabile state \[ibid\] . In origin the tachyons in general was considered only in amplifying mediums \[13-16\]. As it seems from our investigations factually for the drift velocities more than the velocity of light the super-luminal particles are generate or are amplify the electromagnetic waves (photons) and the super-luminal particles are placed on the regime of generation or amplification independently from on type of medium (see also ). As it will be shown in a special report in general all elementary excitations including so- called ” elementary particles ” are the quasi-particles.
6. In the point of $`u=c`$ the angle $`\alpha `$ between the $`\stackrel{}{u}`$ and $`\stackrel{}{q}`$ is equal to $`\pi `$$`/2`$ and we have the condition when the electromagnetic wave becomes a free and is emitted. Thus the point of $`u=c`$ is the point of transition of system from absorption regime to the regime of emission of electromagnetic waves (photons).
7. It is shown that the relativistic expression for the deceleration of time really taking place for the relaxation time of carriers on photons, for the life-time of carriers and also for the period of electromagnetic oscillations. Factually, in dynamics and electrodynamics the time is enter as a parameter, but not as a free coordinate and for this reason the Einstein relation is impossible to apply to the time.
Since $`\tau ^1\mathrm{~}N_i(q,T_i)`$ we have $`\tau _i\tau _o\left(1u^2/c^2\right)\left(T/T_i\right)`$; $`l_i=u\tau _i=l_o\left(1u^2/c^2\right)\left(T/T_i\right)`$. If $`T_i=T`$ we have $`\tau _i\tau _o\left(1u^2/c^2\right)`$; $`l=l_o\left(1u^2/c^2\right)`$.
8. It is shown that the so-called ”velocity of light” in vacuum $`c`$, as well as the velocity of sound for phonons, is an averaged velocity of photons at the ground state.
9. As it seems from (4), (7), (8) and (10) at the point $`u=c`$ the distribution function of photons is an isotropic one and the anisotropy part of photons distribution function at this point is equal to zero. It means that the ground and the stationary states of electromagnetic field is spherically symmetric (this question will be discussed in a special report) .At the point $`u=c`$ the stimulated absorption and emission are equal to each other and there are only spontaneous emission of photons.
10. It is shown that the demand for invariance of the Maxwell equations or the laws of physics in all inertial frames of reference is not correct. It is equivalent to the demand of invariance the lows of physics in the uniform and non-uniform spaces or to the demand of equivalency the lows of physics in cases of the presence and absence of external field (force).
11. It is shown that the presence of the second inertial frame of reference moving with the constant drift velocity relatively to first one may be a result of the presence of space non-uniformity or the non-uniform external field.
Really in the uniform space because of the equivalency of all points of space it is impossible to get simultaneously two, or more inertial frame of reference with the different constant drift velocity relatively to each other. Because it will lead to the violation of the space uniformity, i.e. to the change of the distance $`r_{12}`$ between the initial points of that frames of reference ($`r_{12}const.`$). In the uniform space for the conservation of the space uniformity during the motion it is necessary a motion of all points of the space with the same constant velocity (in the absence of the external field), or with the constant acceleration (in the presence of the uniform force). For the both cases the frames of reference, connected with the different points of the space, do not have the motion relatively to each other without the violation of the space uniformity or the uniformity of the external force. Since all points of the spaces in both cases are placed on the same conditions and are equivalent. In the second case for the reason that the force in the Newton’s second law is external one the system (or the space) must be opened. To choose two frames of reference moving with the constant drift velocity relatively to each other, it is necessary to have the space, which consists of two half-spaces. All points of first half-space are in rest or have the motion with the constant velocity v and all points of second half-space are move with the constant acceleration. The first of them is inertial, but the second accelerates and for this reason the demand of equivalency of laws of physics in that two frames of reference is equivalent to the demand of the equivalency of first and second Newton’s laws. This demand is absurd, of course!
12. As it seems from (1) - (4) the demand of the equivalency of laws of physics in the all inertial frames of reference for the photons is equivalent to the demand about equivalency the Plank’ s equilibrium distribution function to the drifted Plank’s distribution function.
13. It is shown that at the external electric fields $`E<H`$ the drift velocity of carriers $`u<c`$ and the energy received from external field is accumulated by the connection of carriers with photons, as a result of the mutual drag (as a result of the construction of structure by the mechanism of ” dressing of carriers with photons ”) and stationary state is conserved. In this region of the drift velocities the absorption is more than the emission. Under the conditions $`E>H`$ i.e. at the drift velocities $`u>c`$ the generation and amplification is dominate and there are the exponential grown the number of photons by the time. The violation of the stationary state is begins from the point $`u=c`$ and from this point is begins the transition of the carriers” dressed by photons ” to the following stationary state by the reactive emission of photons. In the region of drift velocities $`u>c`$ the anisotropy part of the photons distribution function is much more than the isotropic one. Thus in our investigation the mechanism of the transition of particles to the following stationary state is obtained.
14. It is shown that for the drifted velocities $`u>c`$ so - called energy (or mass) for one mode are:
$$ϵ=\left(\frac{T_i}{T}\right)^2\frac{TN(q,T)}{u/c1}\left\{\left[\mathrm{exp}\left(\gamma _qt\right)1\right]+\frac{T}{T_i}\mathrm{exp}\left(\gamma _qt\right)\right\}=$$
$$=\frac{M_oc^2}{u/c1}\left(\frac{T_i}{T}\right)^2\left\{\left[\mathrm{exp}\left(\gamma _qt\right)1\right]+\frac{T}{T_i}\mathrm{exp}\left(\gamma _qt\right)\right\}$$
(23)
or the mass for one mode:
$$M=\frac{M_o}{1u/c}\left\{\left[\mathrm{exp}\left(\gamma _qt\right)1\right]+\frac{m_o}{m_i}\mathrm{exp}\left(\gamma _qt\right)\right\}\text{ (23’)}$$
In the absence of heating
$$M=\frac{M_o}{u/c1}\left[2\mathrm{exp}\left(\gamma _qt\right)1\right]\frac{2M_o}{u/c1}\mathrm{exp}\left(\gamma _qt\right)$$
(24)
As it seems from this equations the mass of photons for one mode in this region of drift velocities is $`u>c`$ grows by the time exponentially.
The case considered by Einstein may be correspond to the case when from two choosing frame of reference the first was connected with the photons and drifted together with them and the second was connected with the charged carriers (electrons or positrons) drifted with constant velocity relatively to first one. Here the photons (electromagnetic field) is plays the role of media and the charged carriers drifted relatively to that media. Thus the system must consists of three subsystems. Factually, in electrodynamics the space (or the system) is consists from three half - space (or subsystems): the half -system of negatively charged carriers, half-space of positively charged carriers and the half-space (or space) of photons. For this reason the electrodynamics space is non-uniform initially, because of the point of space where is placed of negatively charge carrier do not equivalent to the point of space where is placed the positively charged carrier and both do not equivalent to the point where is placed the photon. The presence of the two type of charge lead to the presence of so-called Lorentz force. The condition of stationary of state is the equality of this force to zero $`F=0`$! This condition corresponds to the annihilation of charge carriers with production of photons, i.e. production of free electromagnetic field without charges. It means that the space of photons (i.e. free electromagnetic field) can decay to the two half-spaces: the spaces of the negative and positive charges and also two half-spaces of negative and positive charges can product the space of photons (the electromagnetic space or media).
As it seems from our investigations considered by Lorentz and Einstein case corresponds to the case of the presence of weak external field when the heating of carriers and photons are absent and there are only their mutual drag. Also they was consideredthe case of one type of charged carriers. In the presence of mutual drag of the electrons and photons the distribution function of photons $`N(q)`$ has the form of the displaced Plank’s function with the constant drift velocity and as a result with the renormalized frequency of emission $`\omega _{em}=\omega _q^{}=\omega _q\stackrel{}{u}\stackrel{}{q}/\mathrm{}`$. For the drift velocities $`u<c`$ is decrease with the increasing of drift velocity $`u`$ because of the stationary function of photons has the form (18) or (10).
Thus in the second frame of reference charge carrier is emitted or absorbed photon with the frequency $`\omega _{em}^{}=\omega _{abs.}[1\stackrel{}{u}\stackrel{}{q}/\mathrm{}\omega _q]`$. In the one mode case the observed frequency $`\omega _{abs}=\stackrel{\mathrm{\_}}{\omega }=T_i/\mathrm{}`$. In the other word
$$\omega _{abs.}=\omega _{em.}\left(1\frac{u}{c}\mathrm{cos}\alpha \right)^1\text{ or }\lambda _{abs.}=\lambda _{em.}\left(1\frac{u}{c}\mathrm{cos}\alpha \right)$$
For the case when the both frames of reference assumed to be inertial one, i.e. to move along the one line (along the $`x`$-axis) cos$`\alpha =1`$ and we have
$$\omega _{abs.}=\omega _{em.}\left(1\frac{u}{c}\right)^1\text{ or }\lambda _{abs.}=\lambda _{em.}\left(1\frac{u}{c}\right)$$
As it seems from this equation the Doppler effect is also a result of the mutual drag of carriers and photons. Factually the source of emission (charge carrier) drifts relatively to frame of reference connected with the photon (observer) with the drift velocity $`\stackrel{}{u}`$. By the increasing of drift velocity $`\stackrel{}{u}`$ the distance between the source and the detector are increased too and as a result the observed frequency of photons is increased or the wavelength is decreased. In the opposite case the wavelength is increased. In the region of $`u<c`$ the source (the charge carrier) moves slowly than the detector (photon) and as a result the distance between them is decreased and the wavelength of observed photons (light) is increased. Thus the mutual drag of electrons and photons in the region $`u<c`$ lead to decreasing the frequency of emission or absorption (see ”Low frequency cyclotron resonance for the phonons ” ). It means that the frequency of observed light (photons) is increased.
15. As it seems from the present consideration the so-called relativistic phenomena and the Doppler effect is a same ones and are a result of the mutual drag of interacting system of carriers and photons at external field. The case was considered by Lorentzand Einstein corresponds to the case, when the charge carriers are drifted at the electromagnetic field under the conditions of the mutual drag of carriers and photons in the absence of their heating by the field.
Distinction between the results of the theory of relativity and the Doppler Effect is connected with them, that the Doppler Effect is dealing with the total stationary distribution function but no only it’s isotropic path as in theory of relativity. In the other words the Doppler Effect is received as a result of taking into account the violation of the $`T`$-symmetry at the external field.
## 4 Reference
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7. Gasymov T. M. 1976. Dokl. Acad. Nauk Azerb. SSR, 32, No 6, p.p. 15 - 17.
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12. Andreev A.U., Kirjnis D.A. 1996, Uspechi Phys. Nauk , 166, No 10, p.p. 1135-1140. \[U1\]
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warning/0001/hep-ph0001103.html | ar5iv | text | # Higgs Boson Production at the LHC with Soft Gluon Effects
## I Introduction
One of the fundamental questions of the Standard Model (SM) of elementary particle physics is the dynamics of the electroweak symmetry breaking. Within the SM, the Higgs mechanism postulates the existence of a scalar field, the elementary excitation of which is called the Higgs boson. Four experimental collaborations at the LEP II collider search for the Higgs boson in the $`e^+e^{}Z^0H`$ process up to 202 GeV center of mass energy. DELPHI and L3 set a preliminary exclusion limit of $`m_H>96`$ GeV on the Higgs mass, followed closely by the limit of $`m_H>91`$ GeV set by ALEPH and OPAL . According to recent preliminary information, the combined lower limit is close to $`m_H\stackrel{>}{}106`$ GeV . Global fits to electroweak observables appear to prefer a low mass Higgs particle, with a mean value close to 100 GeV, and less than 250 GeV within 95 percent of confidence .
Among other aims the main goal of the next proton-proton accelerator, the 14 TeV center of mass energy Large Hadron Collider (LHC) at CERN, is to establish the existence of the Higgs boson and to measure its basic properties. At the LHC a light SM Higgs boson will be mainly produced through the partonic subprocess $`gg`$ (via top quark loop) $`HX`$ . It can be detected, after about 1.5 years of running with a statistical significance of at least 4, in its $`H\gamma \gamma `$ decay mode, if its mass is in the 100-150 GeV range . If the Higgs mass is higher than about 150 GeV, then its $`HZ^{0()}Z^0`$ decay mode is the cleanest and most significant . In this letter our focus is on the Higgs boson production, and in our numerical illustration we choose $`m_H=150`$ GeV.
According to earlier studies, a statistical significance on the order of 5-10 can be reached for the inclusive $`H\gamma \gamma `$ and for the $`H+\text{jet}\gamma \gamma +\text{jet}`$ signals, although actual values depend on luminosity and background estimates. In Ref. it was found that in order to optimize the significance it is necessary to impose a 30 GeV cut on the transverse momentum of the jet, or equivalently (at next-to-leading order precision), on the transverse momentum ($`Q_T`$) of the photon pair. With this cut in place, extraction of the signal in the Higgs + jet mode requires the precise knowledge of both the signal and background distributions in the medium to high $`Q_T`$ region.
To reliably predict the $`Q_T`$ distribution of Higgs bosons at the LHC, especially for the low to medium $`Q_T`$ region where the bulk of the rate is, the effects of the multiple soft–gluon emission have to be included. One approach to achieve this is parton showering . Although the universality of this method makes it a very powerful tool, present drawbacks of this ansatz are the lack of the proper normalization (which takes into account the full fixed order QCD corrections), the lack of exact matrix elements even in the high $`Q_T`$ region, and the lack of uniqueness of the prediction (”tunability”). There is ongoing work to correct these problems .
A more reliable prediction of the Higgs $`Q_T`$ can be obtained utilizing the Collins-Soper-Sterman (CSS) resummation formalism , which takes into account the effects of the multiple soft–gluon emission while reproducing the rate, systematically including the higher order corrections. It is possible to smoothly match the CSS result to the fixed order one in the medium to high $`Q_T`$ region, thus obtaining the best prediction in the full $`Q_T`$ region . Compared to fixed orders, the resummed result depends on a few extra parameters. These parameters are new renormalization scales ($`C_i`$, only two of which are independent) , and a few universal, non-perturbative parameters ($`g_i`$), which are extracted from present experiments and then used to predict the results of future ones . In this letter, we use this formalism to calculate the total cross sections and $`Q_T`$ distributions of Higgs bosons at the LHC.
Our results here, together with the resummed calculations for the diphoton and $`Z^0`$ boson pair production , provide a consistent set of QCD calculations of the transverse momentum (and other) distribution(s) of the Higgs bosons and their backgrounds at the LHC. These results systematically include both the multiple soft–gluon effects and the finite order QCD corrections, and can be used to tune the shower MC’s which experimentalists extensively use when extracting the Higgs signal, or can be utilized independently by means of the ResBos Monte Carlo event generator .
## II Analytical Results
Within the SM , as well as in the minimal supersymmetric standard model (MSSM) with small $`\mathrm{tan}\beta `$ , the dominant production mode of neutral Higgs bosons at the LHC is gluon fusion via a heavy quark loop. The lowest order cross section of this process is formally $`𝒪(\alpha _s^2)`$ in the strong coupling. Fixed order QCD corrections to this production mechanism are known to substantially increase the rate. For a light Higgs boson the $`𝒪(\alpha _s^3)`$ to $`𝒪(\alpha _s^2)`$ $`K`$-factor is in the order of 2 (cf. Fig. 1). The full $`𝒪(\alpha _s^4)`$ calculation is not completed yet, but the real emission and the virtual contributions are separately available. Resummed calculations, taking into account the soft–gluon effect, were also performed to estimate the size of the uncalculated higher order corrections , as well as to predict the shape of the Higgs transverse momentum distribution .
In this work we use the Collins-Soper-Sterman (CSS) soft–gluon resummation formalism to calculate the QCD corrections from the multiple–soft gluon emission. Calculations similar to this were earlier performed in Refs. . Our present calculation improves these by including $`𝒪(\alpha _s^4)`$ terms in the Sudakov exponent, by applying the state of the art matching to the latest fixed order distributions, by using a QCD improved gluon-Higgs effective coupling , by utilizing an improved non-perturbative function, and by including the effect of the Higgs width.
We also utilize the approximation that the object which couples the gluons to the Higgs (the top quark in the SM) is much heavier than the Higgs itself. This approximation is not essential to our calculation and can be released by including the complete Wilson coefficients with all the relevant masses. The heavy quark approximation in the SM has been shown to be reliable within 5 percent for $`m_H<2m_t`$ , and still reasonable even in the range of $`m_H\stackrel{>}{}2m_t`$ . It has also been shown that the approximation remains valid for the $`Q_T`$ distribution in the large $`Q_T`$ region, provided that $`m_H<m_t`$ and $`Q_T<m_t`$ . In this work we assume that the approximation is valid in the whole $`Q_T`$ region. In the MSSM the heavy quark approximation is also a reliable ansatz for the case of the light Higgs boson and small $`\mathrm{tan}\beta `$ when the Yukawa coupling of the bottom quark is negligible (c.f. and references therein). Using the CSS formalism we resum large logs of the type $`\mathrm{ln}(Q/Q_T)`$ in the low $`Q_T`$ region, and we match the resummed result to the fixed order calculation which is valid for high $`Q_T`$ . We also include the $`qg`$ and $`q\overline{q}`$ subprocesses which, in combination, can constitute up to 10 percent of the total rate, depending on the Higgs mass .
The resummed differential cross section of a neutral Higgs boson, denoted by $`\varphi ^0`$ in the SM or MSSM, produced in hadronic collisions is written as
$`{\displaystyle \frac{d\sigma (h_1h_2\varphi ^0X)}{dQ^2dQ_T^2dy}}=\sigma _0{\displaystyle \frac{Q^2}{S}}{\displaystyle \frac{Q^2\mathrm{\Gamma }_\varphi /m_\varphi }{(Q^2m_\varphi ^2)^2+(Q^2\mathrm{\Gamma }_\varphi /m_\varphi )^2}}`$ (1)
$`\times \{{\displaystyle \frac{1}{(2\pi )^2}}{\displaystyle }d^2be^{i\stackrel{}{Q}_T\stackrel{}{b}}\stackrel{~}{W}_{gg}(b_{},Q,x_1,x_2,C_{1,2,3})`$ (2)
$`\times \stackrel{~}{W}_{gg}^{NP}(b,Q,x_1,x_2)+Y(Q_T,Q,x_1,x_2,C_4)\}.`$ (3)
The kinematical variables $`Q`$, $`Q_T`$, and $`y`$ are the invariant mass, transverse momentum, and rapidity of the Higgs boson, respectively, in the laboratory frame. The parton momentum fractions are defined as $`x_1=e^yM_T/\sqrt{S}`$, and $`x_2=e^yM_T/\sqrt{S}`$, with $`M_T=\sqrt{Q^2+Q_T^2}`$, and $`\sqrt{S}`$ being the center–of–mass (CM) energy of the hadrons $`h_1`$ and $`h_2`$. The lowest order cross section is
$$\sigma _0=\kappa _\varphi (Q)\frac{\sqrt{2}G_F\alpha _s^2(Q^2)}{576\pi },$$
(4)
where $`G_F`$ is the Fermi constant, and $`\kappa _\varphi `$, the QCD corrected effective coupling of the Higgs boson to gluons in the heavy top quark limit (cf. Ref.), is defined as
$`\kappa _\varphi (Q)`$ $`=`$ $`1+{\displaystyle \frac{11}{2}}{\displaystyle \frac{\alpha _s^{(5)}(m_t^2)}{\pi }}+{\displaystyle \frac{3866201N_f}{144}}\left({\displaystyle \frac{\alpha _s^{(5)}(m_t^2)}{\pi }}\right)^2`$ (6)
$`+{\displaystyle \frac{15319N_f}{332N_f}}{\displaystyle \frac{\alpha _s^{(5)}(Q^2)\alpha _s^{(5)}(m_t^2)}{\pi }}+𝒪(\alpha _s^3)`$
where $`\alpha _s^{(5)}`$ is the strong coupling constant in the $`\overline{\mathrm{MS}}`$ scheme with 5 active flavors, and $`m_t`$ denotes the pole mass of the top quark.
The renormalization group invariant kernel of the Fourier integral $`\stackrel{~}{W}_{gg}(b_{},Q,x_1,x_2,C_{1,2,3})`$, and the $`Q_T`$ regular term $`Y(Q_T,Q,x_1,x_2,C_4)`$, together with the variables $`b_{}`$ and $`C_1`$ to $`C_4`$, are given in Ref. . The definition of $`\stackrel{~}{W}_{gg}`$, contains the Sudakov exponent
$`𝒮(Q,b_{},C_1,C_2)=`$ (7)
$`{\displaystyle _{C_1^2/b_{}^2}^{C_2^2Q^2}}{\displaystyle \frac{d\overline{\mu }^2}{\overline{\mu }^2}}[A(\alpha _s(\overline{\mu }),C_1)\mathrm{ln}\left({\displaystyle \frac{C_2Q^2}{\overline{\mu }^2}}\right)+`$ (8)
$`B(\alpha _s(\overline{\mu }),C_1,C_2)].`$ (9)
In the perturbative expansion of the $`A(\alpha _s(\overline{\mu }),C_1)`$ and $`B(\alpha _s(\overline{\mu }),C_1,C_2)`$ functions we follow the notation of Ref. . In our present calculation, we include the process independent next-to-next-to-leading order coefficient
$`A^{(2)}(C_1)=`$ (10)
$`4C_A\left[\left({\displaystyle \frac{67}{36}}{\displaystyle \frac{\pi ^2}{12}}\right)N_C{\displaystyle \frac{5}{18}}N_f2\beta _1\mathrm{ln}\left({\displaystyle \frac{b_0}{C_1}}\right)\right],`$ (11)
in the expansion of the $`A(\alpha _s(\overline{\mu }),C_1)`$ function, where $`C_A=3`$ is the Casimir of the adjoint representation of $`SU(3)`$, $`N_C=3`$ is the number of colors, and $`N_f=5`$ is the number of active quark flavors. With the inclusion of $`A^{(2)}`$ the only missing next-to-next-to-leading order contribution in the Sudakov exponent is the $`B^{(2)}`$ term, which is suppressed by $`1/\mathrm{ln}\left(\frac{Q^2}{Q_T^2}\right)`$ with respect to $`A^{(2)}`$, and by $`\alpha _s`$ with respect to $`B^{(1)}`$. This is illustrated by the expansion of the asymptotic part of the cross section:
$`\underset{Q_T0}{lim}{\displaystyle \frac{d\sigma }{dQ^2dQ_T^2dy}}=`$ (12)
$`\sigma _0{\displaystyle \frac{1}{Q_T^2}}{\displaystyle \underset{i,j}{}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{m=0}{\overset{2n1}{}}}\left({\displaystyle \frac{\alpha _s(Q)}{\pi }}\right)^nC_{nm}^{(ij)}\mathrm{ln}^m\left({\displaystyle \frac{Q^2}{Q_T^2}}\right),`$ (13)
where $`i`$ and $`j`$ label incoming partons. While the $`A^{(2)}`$ coefficient contributes to the above series via
$`C_{21}^{(ij)}`$ (14)
$`{\displaystyle \frac{1}{2}}\left[\left(B^{(1)}\right)^2A^{(2)}\beta _0A^{(1)}\mathrm{ln}\left({\displaystyle \frac{\mu _R^2}{Q^2}}\right)\beta _0B^{(1)}\right]f_if_j,`$ (15)
the $`B^{(2)}`$ coefficient only occurs as
$`C_{20}^{(ij)}`$ (17)
$`\left[\zeta (3)\left(A^{(1)}\right)^2+{\displaystyle \frac{1}{2}}B^{(2)}+{\displaystyle \frac{1}{2}}\beta _0B^{(1)}\mathrm{ln}\left({\displaystyle \frac{\mu _R^2}{Q^2}}\right)\right]f_if_j,`$ (18)
where our notation coincides with that of Ref. . Hence, the contribution from the uncalculated $`B^{(2)}`$, compared to that from $`A^{(2)}`$, is expected to be smaller because of the additional log weighting the $`A^{(2)}`$ coefficient. To estimate the size of the contribution from $`B^{(2)}`$, we follow the usual practice in a perturbative calculation by varying the renormalization constants ($`C_1`$ and $`C_2`$) in the Sudakov factor by a factor of 2. The results are shown in Fig. 2.
The form of our non-perturbative function $`\stackrel{~}{W}_{gg}^{NP}`$ coincides with the one used for the $`gg\gamma \gamma `$ process in Ref.
$`\stackrel{~}{W}_{gg}^{NP}(b,Q,Q_0,x_1,x_2)=`$ (19)
$`\mathrm{exp}\left[g_1b^2{\displaystyle \frac{C_A}{C_F}}g_2b^2\mathrm{ln}\left({\displaystyle \frac{Q}{2Q_0}}\right)g_1g_3b\mathrm{ln}(100x_1x_2)\right],`$ (20)
(21)
where the Casimir of the fundamental $`SU(3)`$ representation is $`C_F=4/3`$. The values of the non-perturbative parameters $`g_i`$ are defined in Ref. . The uncertainties of the resummed distribution, stemming from the non-perturbative function, were found to be in the order of 5 percent (cf. ). In the high $`Q_T`$ region Eq. (3) is matched to the fixed order perturbative result (at the $`𝒪(\alpha _s)`$) of Ref. in the manner described in Ref. .
## III Numerical Results
The analytic results are coded in the ResBos Monte Carlo event generator , which uses the following electroweak input parameters :
$`G_F=1.16639\times 10^5\mathrm{GeV}^2,m_Z=91.187\mathrm{GeV},`$ (22)
$`m_W=80.36\mathrm{GeV}.`$ (23)
As in the background calculation , we use the canonical choice of the renormalization constants ($`C_1=C_3=2e^{\gamma _E}C_0`$ and $`C_2=C_4=1`$ ), the NLO expressions for the running electromagnetic and strong couplings $`\alpha (\mu )`$ and $`\alpha _S(\mu )`$, as well as the NLO parton distribution function set CTEQ4M (defined in the modified minimal subtraction, $`\overline{MS}`$, scheme) . We set the renormalization scale equal to the factorization scale: $`\mu _R=\mu _F=Q`$. In the choice of the non-perturbative parameters, we follow Ref. . Since we are not concerned with the decays of Higgs bosons in this work, we do not impose any kinematic cuts. We defer the more extensive study, including various decay modes and QCD backgrounds, to a future publication Ref. .
Fig. 1 displays Higgs boson production cross sections via the gluon fusion process at the LHC, calculated with various QCD corrections in the SM as the function of the Higgs mass. The ratio of the fixed order $`𝒪(\alpha _s^3)`$ (dashed) and the lowest order $`𝒪(\alpha _s^2)`$ (dotted) curves varies between 2.0 and 2.3. We note that less than 2 percent of the $`𝒪(\alpha _s^3)`$ corrections come from the $`qg`$ and $`q\overline{q}`$ initial states for Higgs masses below 200 GeV. The resummed curve is about 10 percent higher than the $`𝒪(\alpha _s^3)`$ one, as expected based on the findings that the CSS formalism preserves the fixed order rate within the error of the matching (the latter being higher order) . The resummed rate is close to the $`𝒪(\alpha _s^3)`$, because we used the $`𝒪(\alpha _s^3)`$ fixed order results to derive the Wilson coefficients which are utilized in our calculation. In Ref. a resummed calculation estimated the size of the $`𝒪(\alpha _s^4)`$ corrections, and a typical value of 1.5 of the $`𝒪(\alpha _s^4)`$ to $`𝒪(\alpha _s^3)`$ $`K`$-factor can be inferred from that work. Based on this, we also plot the $`𝒪(\alpha _s^3)`$ curve rescaled by 1.5, to illustrate the possible size of the $`𝒪(\alpha _s^4)`$ corrections and to establish the normalization of our resummed calculation among the fixed order results.
Fig. 2 compares the Higgs boson transverse momentum distributions calculated by ResBos (curves) and by PYTHIA (histograms from version 6.122). The middle solid curve is calculated using the canonical choice for the renormalization constants in the Sudakov exponent: $`C_1=C_0`$, and $`C_2=1`$. To estimate the size of the uncalculated $`B^{(2)}`$ term, we varied these renormalization constants multiplying both by 1/2 and 2. The upper solid curve shows the result for $`C_1=C_0/2,C_2=1/2`$, and the lower solid curve for $`C_1=2C_0,C_2=2`$. The band between these two curves gives the order of the uncertainty following from the exclusion of $`B^{(2)}`$. The typical size of this uncertainty, e.g. around the peak region, is in the order of $`\pm 10`$ percent. The corresponding uncertainty in the total cross section is also in the same order. The dashed PYTHIA histogram is plotted without altering its output. The normalization of PYTHIA, as that of any parton shower MCs, is the lowest order $`𝒪(\alpha _s^2)`$ for this process. The default PYTHIA histogram is also plotted after the rate is multiplied by the factor $`K=2`$ (dotted). The shape of the PYTHIA histogram agrees reasonably with the resummed curve in the low and intermediate $`Q_T`$ ($`125`$ GeV) region. For large $`Q_T`$ the PYTHIA prediction falls under the ResBos curve, since ResBos mostly uses the exact fixed order $`𝒪(\alpha _s^3)`$ matrix elements in that region (c.f. Ref. ), while PYTHIA still relies on the multi-parton radiation ansatz. PYTHIA can be tuned to agree with ResBos in the high $`Q_T`$ region (dash-dotted), by changing the maximal virtuality a parton can acquire in the course of the shower, i.e. the $`Q_{max}^2`$ parameter, from the default value to the partonic center of mass energy $`s`$. In that case, however, the low $`Q_T`$ region will have disagreement, since the normalization in PYTHIA is conserved, so events in the low $`Q_T`$ region are depleted.
## IV Conclusions
In this letter we presented Higgs boson production rates and $`Q_T`$ distributions for the LHC, including $`𝒪(\alpha _s^3)`$ fixed order QCD and multiple soft–gluonic corrections by means of the CSS resummation formalism. We showed that the resummed rate recovers the fixed order $`𝒪(\alpha _s^3)`$ rate, as expected within the CSS formalism. We investigated the uncertainty of the resummed prediction due to uncalculated terms in the Sudakov exponent. We found that the shape of the resummed prediction in the low $`Q_T`$ region is in reasonable agreement with the default result of PYTHIA.
## Acknowledgments
We thank the CTEQ Collaboration, and M. Spira for invaluable discussions. We are indebted to I. Puljak for providing us with the PYTHIA results. C. B. thanks the organizers of the les Houches SUSY/Higgs/QCD workshop for their hospitality and financial support, and the Fermilab Theory Group for their invitation and financial support. This work was supported in part by the NSF under grant PHY–9802564 and by DOE under grant DE-FG-03-94ER40833. |
warning/0001/hep-ph0001249.html | ar5iv | text | # 1 Masses of the SUSY particles, in GeV, for the mAMSB model with 𝑚₀=200GeV, 𝑚_{3/2}=35TeV, tan𝛽=3, and sgn𝜇=+ from ISAJET (left side) and from Ref. (right side) using the ISAJET sign conventions.
BNL-HET-00/1
UCD-2000-5
hep-ph/0001249
Anomaly Mediated SUSY Breaking at the LHC
Frank E. Paige<sup>a</sup> and James Wells<sup>b</sup>
<sup>a</sup>Physics Department, Brookhaven National Laboratory, Upton, NY 11973
<sup>b</sup>Physics Department, University of California, Davis, CA 95616
Anomaly Mediated SUSY Breaking models are reviewed. Possible signatures at the LHC for one case of the minimal realistic model are examined.
Introduction
The signatures for SUSY at the LHC depend very much on the SUSY masses, which presumably result from spontaneous SUSY breaking. It is not possible to break SUSY spontaneously using just the MSSM fields; instead one must do so in a hidden sector and then communicate the breaking through some interaction. In supergravity models, the communication is through gravity. In gauge mediated models it is through gauge interactions; the gravitino is then very light and can play an important role. Simple examples of both have been discussed previously. A third possibility is that the hidden sector does not have the right structure to provide masses through either mechanism; then the leading contributions come from a combination of gravity and anomalies. This is known as Anomaly Mediated SUSY Breaking (AMSB), and it predicts a different pattern of masses and signatures.
Anomaly-Mediated Supersymmetry Breaking
In the supersymmetric standard model there exist AMSB contributions to the soft mass parameters that arise via the superconformal anomaly . The effect can be understood by recognizing several important features of supersymmetric theories. First, supersymmetry breaking can be represented by a chiral superfield $`\mathrm{\Phi }=1+m_{3/2}\theta ^2`$ which also acts as a compensator for super-Weyl transformations. Treating $`\mathrm{\Phi }`$ as a spurion, one can transform a theory into a super-conformally invariant theory. Even if a theory is superconformal at the outset (i.e., no dimensionful couplings), the spurion $`\mathrm{\Phi }`$ is employed since the quantum field theory requires a regulator that implies scale dependence (Pauli-Villars mass, renormalization scale in dimensional reduction, etc.). To preserve scale invariance the renormalization scale parameter $`\mu `$ in a quantum theory then becomes $`\mu /\sqrt{\mathrm{\Phi }^{}\mathrm{\Phi }}`$. It is the dependence of the regulator on $`\mathrm{\Phi }`$ that induces supersymmetry breaking contributions to the scalars and gauginos.
The anomaly induced masses can be derived straightforwardly for the scalar masses. The Kähler kinetic terms depend on wave function renormalization as in the following superfield operator,
$$d^2\theta d^2\overline{\theta }Z_Q\left(\frac{\mu }{\sqrt{\mathrm{\Phi }^{}\mathrm{\Phi }}}\right)Q^{}Q.$$
(1)
Taylor expanding $`Z`$ around $`\mu `$ and projecting out the $`FF^{}`$ terms yields a supersymmetry breaking mass for the scalar field $`\stackrel{~}{Q}`$:
$$m_{\stackrel{~}{Q}}^2=\frac{1}{4}\frac{d^2\mathrm{ln}Z_Q}{d(\mathrm{ln}\mu )^2}m_{3/2}^2=\frac{1}{4}\left(\frac{\gamma _Q}{g}\beta _g+\frac{\gamma _Q}{y}\beta _y\right)m_{3/2}^2.$$
(2)
Similar calculations can be done for the gauginos and the $`A`$ terms:
$`M_i`$ $`=`$ $`{\displaystyle \frac{g_i^2}{2}}{\displaystyle \frac{dg_i^2}{d\mathrm{ln}\mu }}m_{3/2}={\displaystyle \frac{\beta _{g_i}}{g_i}}m_{3/2},`$ (3)
$`A_y`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{a}{}}{\displaystyle \frac{d\mathrm{ln}Z_{Q_a}}{d\mathrm{ln}\mu }}m_{3/2}={\displaystyle \frac{\beta _y}{y}}m_{3/2}`$ (4)
where the sum over $`a`$ includes all fields associated with the Yukawa coupling $`y`$ in the superpotential.
There are several important characteristics of the AMSB spectrum to note. First, the equations for the supersymmetry breaking contributions are scale invariant. That is, the value of the soft masses at any scale is obtained by simply plugging in the gauge couplings and Yukawa couplings at that scale into the above formulas. Second, the masses are related to the gravitino mass by a one loop suppression. In AMSB $`M_im_{3/2}\alpha _i/4\pi `$, whereas in SUGRA $`M_im_{3/2}`$. While the AMSB contributions are always present in a theory independent of how supersymmetry breaking is accomplished, they may be highly suppressed compared to standard hidden sector models. Therefore, for AMSB to be the primary source of scalar masses, one needs to assume or arrange that supersymmetry breaking is not directly communicated from a hidden sector. This can be accomplished, for example, by assuming supersymmetry breaking on a distant brane . Finally, the squared masses of the sleptons are negative (tachyonic) because $`\beta _g>0`$ for $`U(1)`$ and $`SU(2)`$ gauge groups. This problem rules out the simplest AMSB model based solely on eqs. 2-4.
Given the tachyonic slepton problem, it might seem most rational to view AMSB as a good idea that did not quite work out. However, there are many reasons to reflect more carefully on AMSB. As already mentioned above, AMSB contributions to scalar masses are always present if supersymmetry is broken. Soft masses in the MSSM come for free, whereas in all other successful theories of supersymmetry breaking a communication mechanism must be detailed. In particle, hidden sector models require singlets to give the gauginos an acceptable mass. In AMSB, singlets are not necessary. Also, there may be small variations on the AMSB idea that can produce a realistic spectrum and can have important phenomenological consequences. This is our motivation for writing this note.
Two realistic minimal models of AMSB: mAMSB and DAMSB
As we discussed in the introduction, the pure AMSB model gives negative squared masses for the sleptons, thus breaking electromagnetic gauge invariance, so some additional contributions must be included. The simplest assumption that solves this problem is to add at the GUT scale a single universal scalar mass $`m_0^2`$ to all the sfermions’ squared masses. We will call this model mAMSB. The description and many phenomenological implications of this model are given in Refs. . The parameters of the model after the usual radiative electroweak symmetry breaking are then
$$m_0,m_{3/2},\mathrm{tan}\beta ,sgn\mu =\pm .$$
This model has been implemented in ISAJET 7.48 ; a pre-release version of ISAJET has been used to generate the events for this analysis.
For this note the AMSB parameters were chosen to be
$$m_0=200\mathrm{GeV},m_{3/2}=35\mathrm{TeV},\mathrm{tan}\beta =3,sgn\mu =+$$
For this choice of parameters the slepton squared masses are positive at the weak scale, but they are still negative at the GUT scale. This means that charge and color might be broken (CCB) at high temperatures in the early universe. However, at these high energies there are also large finite temperature effects on the mass, which are positive (symmetry restoration occurs at higher $`T`$). In fact, a large class of SUSY models with CCB minima naturally fall into the correct SM minimum when you carefully follow the evolution of the theory from high T to today. If CCB minima are excluded at all scales, then the value of $`m_0`$ must be substantially larger, so the sleptons must be quite heavy.
The masses from ISAJET 7.48 for this point are listed in Table 1. The mass spectrum has some similarity to that for SUGRA Point 5 studied previously : the gluino and squark masses are similar, and the decays $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}`$ and $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0h`$ are allowed. Thus, many of the techniques developed for Point 5 are applicable here. But there are also important differences. In particular, the $`\stackrel{~}{\chi }_1^\pm `$ is nearly degenerate with the $`\stackrel{~}{\chi }_1^0`$, not with the $`\stackrel{~}{\chi }_2^0`$. The mass splitting between the lightest chargino and the lightest neutralino must be calculated as the difference between the lightest eigenvalues of the full one-loop neutralino and chargino mass matrices. The mass splitting is always above $`m_{\pi ^\pm }`$, thereby allowing the two-body decays $`\chi _1^\pm \chi _1^0+\pi ^\pm `$ . Decay lifetimes of $`\chi ^\pm `$ are always less than 10 cm over mAMSB parameter space, and are often less than 1 cm.
Another unique feature of the spectrum is the near degeneracy of the $`\stackrel{~}{\mathrm{}}_L`$ and $`\stackrel{~}{\mathrm{}}_R`$ sleptons. The mass splitting is
$`m_{\stackrel{~}{\mathrm{}}_L}^2m_{\stackrel{~}{\mathrm{}}_R}^20.037\left(m_Z^2\mathrm{cos}2\beta +M_2^2\mathrm{ln}{\displaystyle \frac{m_{\stackrel{~}{\mathrm{}}_R}}{m_Z}}\right).`$ (5)
There is no symmetry requiring this degeneracy, but rather it is an astonishing accident and prediction of the mAMSB model.
It is instructive to compare the masses from ISAJET with those calculated in Ref. to provide weak-scale input to ISAJET. These masses are listed in the right hand side of Table 1. Since the agreement is clearly adequate for the purposes of the present study, no attempt has been made to understand or resolve the differences. It is clear, however, that if SUSY is discovered at the LHC and if masses or combinations of masses are measured with the expected precision, then more work is needed to compare the LHC results with theoretical models in a sufficiently reliable way.
Another variation on AMSB is deflected AMSB (DAMSB). The idea is based on Ref. who demonstrated that realistic sparticle spectrums with non-tachyonic sleptons can be induced if a light modulus field $`X`$ (SM singlet) is coupled to heavy, non-singlet vector-like messenger fields $`\mathrm{\Psi }_i`$ and $`\overline{\mathrm{\Psi }}_i`$:
$$W_{\mathrm{mess}}=\lambda _\mathrm{\Psi }X\mathrm{\Psi }_i\overline{\mathrm{\Psi }}_i.$$
To ensure gauge coupling unification we identify $`\mathrm{\Psi }_i`$ and $`\overline{\mathrm{\Psi }}_i`$ as $`5+\overline{5}`$ representations of $`SU(5)`$. When the messengers are integrated out at some scale $`M_0`$, the beta functions do not match the AMSB masses, and the masses are deflected from the AMSB renormalization group trajectory. The subsequent evolution of the masses below $`M_0`$ induces positive mass squared for the sleptons, and a reasonable spectrum can result. Although there may be additional significant parameters associated with the generation of the $`\mu `$ and $`B_\mu `$ term in the model, we assume for this discussion that they do not affect the spectra of the MSSM fields. The values of $`\mu `$ and $`B_\mu `$ are then obtained by requiring that the conditions for EWSB work out properly.
The parameters of DAMSB are
$$m_{3/2},n,M_0,\mathrm{tan}\beta ,sgn\mu =\pm $$
where $`n`$ is the number of $`5+\overline{5}`$ messenger multiplets, and $`M_0`$ is the scale at which the messengers are integrated out. Practically, the spectrum is obtained by imposing the boundary conditions at $`M_0`$, and then using SUSY soft mass renormalization group equations to evolve these masses down to the weak scale. Expressions for the boundary conditions can be found in Refs. , and details on how to generate the low-energy spectrum are given in . The resulting spectrum of superpartners is substantially different from that of mAMSB. The most characteristic feature of the DAMSB spectrum is the near proximity of all superpartner masses. In Table 2 we show the spectrum of a model with $`n=5`$, $`M_0=10^{15}\mathrm{GeV}`$, and $`\mathrm{tan}\beta =4`$ as given in . The LSP is the lightest neutralino, which is a Higgsino. (Actually, the LSP is the fermionic component of the modulus $`X`$, but the decay of $`\chi _1^0`$ to it is much greater than collider time scales.) All the gauginos and squarks are between $`300\mathrm{GeV}`$ and $`500\mathrm{GeV}`$, while the sleptons and higgsinos are a bit lighter ($`150\mathrm{GeV}`$ to $`250\mathrm{GeV}`$) in this case.
In summary, we have outlined two interesting directions to pursue in modifying AMSB to make a realistic spectrum. The first direction we call mAMSB, and is constructed by adding a common scalar mass to the sfermions at the GUT scale to solve the negative squared slepton mass problem of pure AMSB. The other direction that we outlined is deflected anomaly mediation that is based on throwing the scalar masses off the pure AMSB renormalization group trajectory by integrating out heavy messenger states coupled to a modulus. The spectra of the two approaches are significantly different, and we should expect the LHC signatures to be different as well. In this note, we study the mAMSB carefully in a few observables to demonstrate how it is distinctive from other, standard approaches to supersymmetry breaking, such as mSUGRA and GMSB.
LHC studies of the example mAMSB model point
We now turn to a study of the example mAMSB spectra presented in Table 1. A sample of $`10^5`$ signal events was generated; since the total signal cross section is $`16\mathrm{nb}`$, this corresponds to an integrated LHC luminosity of $`6\mathrm{fb}^1`$. All distributions shown in this note are normalized to $`10\mathrm{fb}^1`$, corresponding to one year at low luminosity at the LHC. Events were selected by requiring
* At least four jets with $`p_T>100,50,50,50\mathrm{GeV}`$;
* $`\text{ / }E_T>\mathrm{min}(100\mathrm{GeV},0.2M_{\mathrm{eff}})`$;
* Transverse sphericity $`S_T>0.2`$;
* $`M_{\mathrm{eff}}>600\mathrm{GeV}`$;
where the “effective mass” $`M_{\mathrm{eff}}`$ is given by the scalar sum of the missing $`E_T`$ and the $`p_T`$’s of the four hardest jets,
$$M_{\mathrm{eff}}=\text{ / }E_T+p_{T,1}+p_{T,2}+p_{T,3}+p_{T,4}.$$
Standard model backgrounds from gluon and light quark jets, $`t\overline{t}`$, $`W+\mathrm{jets}`$, $`Z+\mathrm{jets}`$, and $`WW`$ have also been generated, generally with much less equivalent luminosity. The $`M_{\mathrm{eff}}`$ distributions for the signal and the sum of all backgrounds with all except the last cut are shown in Figure 1. The ISAJET IDENT codes for the SUSY events contributing to this plot are also shown. It is clear from this plot that the Standard Model backgrounds are small with these cuts, as would be expected from previous studies .
The mass distribution for $`\mathrm{}^+\mathrm{}^{}`$ pairs with the same and opposite flavor is shown in Figure 2. The opposite-flavor distribution is small, and there is a clear endpoint in the same-flavor distribution at
$$M_{\mathrm{}\mathrm{}}^{\mathrm{max}}=\sqrt{\frac{(M_{\stackrel{~}{\chi }_2^0}^2M_\stackrel{~}{\mathrm{}}^2)(M_\stackrel{~}{\mathrm{}}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_\stackrel{~}{\mathrm{}}^2}}=213.6,215.3\mathrm{GeV}$$
corresponding to the endpoints for the decays $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}_{L,R}^\pm \mathrm{}^{}\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{}`$. This is similar to what is seen in SUGRA Point 5, but in that case only one slepton contributes. It is clear from the $`e^+e^{}+\mu ^+\mu ^{}e^\pm \mu ^{}`$ dilepton distribution with finer bins shown in the same figure that the endpoints for $`\stackrel{~}{\mathrm{}}_R`$ and $`\stackrel{~}{\mathrm{}}_L`$ cannot be resolved with the expected ATLAS dilepton mass resolution. More work is needed to see if the presence of two different endpoints could be inferred from the shape of the edge of the dilepton distribution.
Since the main source for $`\stackrel{~}{\chi }_2^0`$ is $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0q`$, information on the squark masses can be obtained by combining the leptons from $`\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}`$ decays with one of the two hardest jets in the event, since the hardest jets are generally products of the squark decays. Figure 3 shows the distribution for the smaller of the two $`\mathrm{}^+\mathrm{}^{}j`$ masses formed with the two leptons and each of the two hardest jets in the event. The dashed curve in this figure shows the same distribution for $`M_{\mathrm{}\mathrm{}}>175\mathrm{GeV}`$, for which the backgrounds are smaller. Both distributions should have endpoints at the kinematic limit for $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0\stackrel{~}{\mathrm{}}\mathrm{}\stackrel{~}{\chi }_1^0\mathrm{}\mathrm{}`$,
$$\left[\frac{(M_{\stackrel{~}{q}_R}^2M_{\stackrel{~}{\chi }_2^0}^2)(M_{\stackrel{~}{\chi }_2^0}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_{\stackrel{~}{\chi }_2^0}^2}\right]^{1/2}=652.9\mathrm{GeV}.$$
Figure 3 also shows the $`\mathrm{}^\pm j`$ mass distribution formed with each of the two leptons combined with the jet that gives the smaller of the two $`\mathrm{}\mathrm{}j`$ masses. This should have a 3-body endpoint at
$$\left[\frac{(M_{\stackrel{~}{q}_R}^2M_{\stackrel{~}{\chi }_2^0}^2)(M_{\stackrel{~}{\chi }_2^0}^2M_\stackrel{~}{\mathrm{}}^2)}{M_{\stackrel{~}{\chi }_2^0}^2}\right]^{1/2}=605.4\mathrm{GeV}.$$
The branching ratio for $`\stackrel{~}{b}_1\stackrel{~}{\chi }_2^0b`$ is very small, so the same distributions with $`b`$-tagged jets contain only a handful of events and cannot be used to determine the $`\stackrel{~}{b}_1`$ mass.
The decay chain $`\stackrel{~}{q}_R\stackrel{~}{\chi }_2^0q\stackrel{~}{\mathrm{}}_{L,R}^\pm \mathrm{}^{}q\stackrel{~}{\chi }_1^0\mathrm{}^+\mathrm{}^{}q`$ also implies a lower limit on the $`\mathrm{}\mathrm{}q`$ mass for a given limit on $`z=\mathrm{cos}\theta ^{}`$ or equivalently on the $`\mathrm{}\mathrm{}`$ mass. For $`z>0`$ (or equivalently $`M_{\mathrm{}\mathrm{}}>M_{\mathrm{}\mathrm{}}^{\mathrm{max}}/\sqrt{2}`$) this lower limit is
$$\begin{array}{ccc}(M_\mathrm{}\mathrm{}q^{\mathrm{min}})^2\hfill & =& \frac{1}{4M_2^2M_e^2}\times \hfill \\ & & [M_1^2M_2^4+3M_1^2M_2^2M_e^2M_2^4M_e^2M_2^2M_e^4M_1^2M_2^2M_q^2\hfill \\ & & M_1^2M_e^2M_q^2+3M_2^2M_e^2M_q^2M_e^4M_q^2+(M_2^2M_q^2)\times \hfill \\ & & \sqrt{(M_1^4+M_e^4)(M_2^2+M_e^2)^2+2M_1^2M_e^2(M_2^46M_2^2M_e^2+M_e^4)}]\hfill \\ M_\mathrm{}\mathrm{}q^{\mathrm{min}}\hfill & =& 376.6\mathrm{GeV}\hfill \end{array}$$
where $`M_q`$, $`M_2`$, $`M_e`$, and $`M_1`$ are the (average) squark, $`\stackrel{~}{\chi }_2^0`$, (average) slepton, and $`\stackrel{~}{\chi }_1^0`$ masses. To determine this lower edge, the larger of the two $`\mathrm{}\mathrm{}j`$ masses formed from two opposite-sign leptons and one of the two hardest jets is plotted in Figure 4. An endpoint at about the right value can clearly be seen.
The $`\mathrm{}^+\mathrm{}^{}`$, $`\mathrm{}^+\mathrm{}^{}q`$, $`\mathrm{}^\pm q`$, and lower $`\mathrm{}^+\mathrm{}^{}q`$ edges provide four constraints on the four masses involved. Since the cross sections are similar to those for SUGRA Point 5, we take the errors at high luminosity to be negligible on the $`\mathrm{}^+\mathrm{}^{}`$ edge, 1% on the $`\mathrm{}^+\mathrm{}^{}q`$ and $`\mathrm{}^\pm q`$ upper edges, and 2% on the $`\mathrm{}^+\mathrm{}^{}q`$ lower edge. Random masses were generated within $`\pm 50\%`$ of their nominal values, and the $`\chi ^2`$ for the four measurements with these errors were used to determine the probability for each set of masses. The resulting distribution for the $`\stackrel{~}{\chi }_1^0`$ mass, also shown in Figure 4, has a width of $`\pm 11.7\%`$, about the same as for Point 5; the errors for the other masses are also comparable. Of course, the masses being measured in this case are different: for example the squark mass is the average of the $`\stackrel{~}{q}_R`$ rather than the $`\stackrel{~}{q}_L`$ masses.
The leptons from $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\mathrm{}^\pm \nu `$ are very soft. This implies that the rate for events with one or three leptons or for two leptons with opposite flavor are all suppressed. Figure 5 shows as a solid histogram the multiplicity of leptons with $`p_T>10\mathrm{GeV}`$ and $`|\eta |<2.5`$ for the AMSB signal with a veto on hadronic $`\tau `$ decays. The same figure shows the distribution for a model with the same weak-scale mass parameters except that the gaugino masses $`M_1`$ and $`M_2`$ are interchanged. This model has a wino $`\stackrel{~}{\chi }_1^\pm `$ approximately degenerate with the $`\stackrel{~}{\chi }_2^0`$ rather than with the $`\stackrel{~}{\chi }_1^0`$. Clearly the AMSB model has a much smaller rate for single leptons and a somewhat smaller rate for three leptons; these rates can be used to distinguish AMSB and SUGRA-like models.
While the decay $`\stackrel{~}{\chi }_2^0\stackrel{~}{\chi }_1^0h`$ is kinematically allowed, the branching ratio is only about 0.03%. Other sources of $`h`$ in SUSY events are also quite small, so in contrast to SUGRA Point 5 there is no strong $`hb\overline{b}`$ signal. However, there is a fairly large branching ratio for $`\stackrel{~}{g}\stackrel{~}{b}\overline{b},\stackrel{~}{t}\overline{t}`$ with $`\stackrel{~}{b}\stackrel{~}{\chi }_1^0b`$, $`\stackrel{~}{t}\chi _1^+b`$, giving two hard $`b`$ jets and hence structure in the $`M_{bb}`$ distribution. For this analysis $`b`$ jets were tagged by assuming that any $`B`$ hadron with $`p_{T,B}>10\mathrm{GeV}`$ and $`|\eta _B|<2`$ is tagged with an efficiency $`ϵ_B=60\%`$; the jet with the smallest
$$R=\sqrt{(\mathrm{\Delta }\eta )^2+(\mathrm{\Delta }\varphi )^2}$$
was then taken to be $`b`$ jets. The two hardest jets generally come from the squarks. To reconstruct $`\stackrel{~}{g}\stackrel{~}{b}\overline{b}`$ one of the two hardest jets, tagged as a $`b`$, was combined with any remaining jet, also tagged as a $`b`$. In addition to the standard multijet and $`\text{ / }E_T`$ cuts, a cut $`M_{\mathrm{eff}}>1200\mathrm{GeV}`$ was made to reduce the Standard Model background. The resulting distributions for the $`b`$ jet multiplicity and for the smallest $`bb`$ dijet mass are shown in Figure 6. The dijet mass should have an endpoint at the kinematic limit for $`\stackrel{~}{g}\stackrel{~}{b}_1\overline{b}\stackrel{~}{\chi }_1^0b\overline{b}`$,
$$M_{bb}^{\mathrm{max}}=\sqrt{\frac{(M_{\stackrel{~}{g}}^2M_{\stackrel{~}{b}}^2)(M_{\stackrel{~}{b}}^2M_{\stackrel{~}{\chi }_1^0}^2)}{M_{\stackrel{~}{b}}^2}}=418.7\mathrm{GeV}.$$
While the figure is roughly consistent with this, the endpoint is not very sharp; more work is needed to assign an error and to understand the high mass tail. There should also be a $`b\overline{t}`$ endpoint resulting from $`\stackrel{~}{g}\stackrel{~}{t}\overline{t}`$, $`\stackrel{~}{t}\stackrel{~}{\chi }_1^+b`$, with $`M_{\stackrel{~}{\chi }_1^+}M_{\stackrel{~}{\chi }_1^0}`$ and essentially invisible. Of course $`m_t`$ has to be kept in the formula. This would be an apparent strong flavor violation in gluino decays and so quite characteristic of these models. Reconstructing the top is more complicated, so this has not yet been studied.
The splitting between the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{\chi }_1^0`$ is very small in AMSB models. ISAJET gives a splitting of $`0.189\mathrm{GeV}`$ for this point and $`c\tau =2.8\mathrm{cm}`$, with the dominant decay being the two-body mode $`\stackrel{~}{\chi }_1^\pm \stackrel{~}{\chi }_1^0\pi `$ via a virtual $`W`$. Ref. gives a somewhat smaller value of $`\mu `$ and so a smaller splitting. The lifetime is of course quite sensitive to the exact splitting. Since the pion or electron is soft and so difficult to reconstruct, it seems better to look for the tail of long-lived winos. The signature is an isolated stiff track in a fraction of the events that ends in the tracking volume and produces no signal in the calorimeter or muon system. Figure 7 shows the radial track length $`R_T`$ distribution in units of $`c\tau `$ for winos with $`|\eta |<1`$ and the (generated) momentum distribution for those with $`R_T>10c\tau `$. Note that the ATLAS detector has three layers of pixels with very low occupancy at radii of 4, 11, and 14 cm and four double layers of silicon strips between 30 and 50 cm. It seems likely that the background for tracks that end after the pixel layers would be small.
It is instructive to compare this signature to that for GMSB models with an NLSP slepton. Both models predict long-lived charged particles with $`\beta <1`$. In the GMSB models, two NLSP sleptons occur in every SUSY event, and they decay into a hard $`e`$’s, $`\mu `$’s, or $`\tau `$’s plus nearly massless $`\stackrel{~}{G}`$’s. In the AMSB models, only a fraction of the SUSY events contain long-lived charged tracks, and these decay into a soft pion or electron plus an invisible particle. A detailed tracking simulation should be done for both cases.
Acknowledgement: This work was supported in part by the U.S. Department of Energy under Contract DE-AC02-98CH10886. We also acknowledge the support of the Les Houches Physics Center, where part of this work was done. |
warning/0001/math0001006.html | ar5iv | text | # Summation and transformation formulas for elliptic hypergeometric series
## 1. Introduction
In the preface to their book “Basic Hypergeometric Series” , Gasper and Rahman refer to the enchanting nature of the theory of $`q`$-series or basic hypergeometric series as the highly infectious “$`q`$-disease”. Indeed, from the time of Heine, about a century and a half ago, till the present day, many researchers have been pursuing the task of finding $`q`$-analogues of classical results in the theory of special functions, orthogonal polynomials and hypergeometric series. It thus seems somewhat surprising that what appears to be the next natural line of research, replacing “$`q`$-analogue” by “elliptic analogue”, has so-far found very few practitioners.
The elliptic (or “modular”) analogues of hypergeometric series were introduced by Frenkel and Turaev in their study of elliptic $`6j`$-symbols. These $`6j`$-symbols, which correspond to certain elliptic solution of the Yang–Baxter equation found by Baxter and Date et al. , can be expressed in terms of elliptic generalizations of terminating, balanced, very-well-poised $`{}_{10}{}^{}\varphi _{9}^{}`$ series. Moreover, the tetrahedral symmetry of the elliptic $`6j`$ symbols implies an elliptic analogue of the famous Bailey transformation for $`{}_{10}{}^{}\varphi _{9}^{}`$ series. So far, the only follow up on the work of Frenkel and Turaev appears to be the paper by Spiridonov and Zhedanov, who presented several contiguous relations for the elliptic analogue of $`{}_{10}{}^{}\varphi _{9}^{}`$ series and who studied an elliptic version of Wilson’s family of biorthogonal rational functions. As an independent development towards elliptic analogues, we should also mention the work by Ruijsenaars and Felder and Varchenko we studied an elliptic variant of the $`q`$-gamma function.
The aim of this paper is to prove several further results for elliptic hypergeometric series. After an introduction to basic and elliptic hypergeometric series in section 2, we use section 3 to derive an elliptic matrix inverse. This matrix inverse, which generalizes a well-known result from the theory of basic series, is used repeatedly in section 4 to derive an extensive list of summation and transformation formulas for terminating, balanced, very-well-poised, elliptic hypergeometric series. The “$`q`$ limits” of most of these identities correspond to known results by Gasper and Rahman, Gessel and Stanton, and Chu. In section 5 we establish an elliptic, multivariable extension of Jackson’s $`{}_{8}{}^{}\varphi _{7}^{}`$ sum associated with the C<sub>n</sub> root system, generalizing the basic case due to Schlosser. Our proof involves an elliptic extension of a general determinant evaluation by Krattenthaler. We conclude the paper with a conjectured C<sub>n</sub> Bailey transformation for elliptic hypergeometric series.
## 2. Basic hypergeometric series and their elliptic analogues
Assume $`|q|<1`$ and define the $`q`$-shifted factorial for all integers $`n`$ by
$$(a;q)_{\mathrm{}}=\underset{k=0}{\overset{\mathrm{}}{}}(1aq^k)\text{and}(a;q)_n=\frac{(a;q)_{\mathrm{}}}{(aq^n;q)_{\mathrm{}}}.$$
Specifically,
$$(a;q)_n=\{\begin{array}{cc}_{k=0}^{n1}(1aq^k)\hfill & n>0\hfill \\ 1\hfill & n=0\hfill \\ 1/_{k=0}^{n1}(1aq^{n+k})=1/(aq^n;q)_n\hfill & n<0\text{.}\hfill \end{array}$$
With the usual condensed notation
$$(a_1,\mathrm{},a_m;q)_n=(a_1;q)_n\mathrm{}(a_m;q)_n$$
we can define an $`{}_{r+1}{}^{}\varphi _{r}^{}`$ basic hypergeometric series as
$${}_{r+1}{}^{}\varphi _{r}^{}[\genfrac{}{}{0.0pt}{}{a_1,a_2,\mathrm{},a_{r+1}}{b_1,\mathrm{},b_r};q,z]=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(a_1,a_2,\mathrm{},a_{r+1};q)_k}{(q,b_1,\mathrm{},b_r;q)_k}z^k.$$
Here it is assumed that the $`b_i`$ are such that none of the terms in the denominator of the right-hand side vanishes. When one of the $`a_i`$ is of the form $`q^n`$ ($`n`$ a nonnegative integer) the infinite sum over $`k`$ can be replaced by a sum from $`0`$ to $`n`$. In this case the series is said to be terminating. A $`{}_{r+1}{}^{}\varphi _{r}^{}`$ series is called balanced if $`b_1\mathrm{}b_r=qa_1\mathrm{}a_{r+1}`$ and $`z=q`$. A $`{}_{r+1}{}^{}\varphi _{r}^{}`$ series is said to be very-well-poised if $`a_1q=a_2b_1=\mathrm{}=a_{r+1}b_r`$ and $`a_2=a_3=qa_1^{1/2}`$. In this paper we exclusively deal with balanced, very-well poised series (or rather, their elliptic analogues) and departing from the standard notation of Gasper and Rahman’s book we use the abbreviation
$`{}_{r+1}{}^{}W_{r}^{}(a_1;a_4,\mathrm{},a_{r+1};q)`$ $`={}_{r+1}{}^{}\varphi _{r}^{}[{\displaystyle \genfrac{}{}{0.0pt}{}{a_1,qa_1^{1/2},qa_1^{1/2},a_4,\mathrm{},a_{r+1}}{a_1^{1/2},a_1^{1/2},qa_1/a_4,\mathrm{},qa_1/a_{r+1}}};q,q]`$
$`={\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1a_1q^{2k}}{1a_1}}{\displaystyle \frac{(a_1,a_4,\mathrm{},a_{r+1};q)_kq^k}{(q,a_1q/a_4,\mathrm{},a_1q/a_{r+1};q)_k}},`$
where we always assume the parameters in the argument of $`{}_{r+1}{}^{}W_{r}^{}`$ to obey the relation $`(a_4\mathrm{}a_{r+1})^2=a_1^{r3}q^{r5}`$.
One of the deepest results in the theory of basic hypergeometric series is Bailey’s transformation , \[17, Eq. (III.28)\]
(2.1)
$$\begin{array}{c}{}_{10}{}^{}W_{9}^{}(a;b,c,d,e,f,g,q^n;q)\hfill \\ \hfill =\frac{(aq,aq/ef,\lambda q/e,\lambda q/f;q)_n}{(aq/e,aq/f,\lambda q/ef,\lambda q;q)_n}{}_{10}{}^{}W_{9}^{}(\lambda ;\lambda b/a,\lambda c/a,\lambda d/a,e,f,g,q^n;q),\end{array}$$
where
$$bcdefg=a^3q^{n+2}\text{and}\lambda =a^2q/bcd.$$
This identity contains many well-known transformation and summation theorems for basic series as special cases. For example, setting $`cd=aq`$ (so that $`\lambda b/a=1`$) and then replacing $`e,f,g`$ by $`c,d,e`$ gives Jackson’s $`q`$-analogue of Dougall’s $`{}_{7}{}^{}F_{6}^{}`$ sum , \[17, Eq. (II.22)\]
(2.2)
$${}_{8}{}^{}W_{7}^{}(a;b,c,d,e,q^n;q)=\frac{(aq,aq/bc,aq/bd,aq/cd;q)_n}{(aq/b,aq/c,aq/d,aq/bcd;q)_n},$$
where
$$bcde=a^2q^{n+1}.$$
To introduce the elliptic analogues of basic hypergeometric series we need the elliptic function
(2.3)
$$E(x)=E(x;p)=(x;p)_{\mathrm{}}(p/x;p)_{\mathrm{}},$$
for $`|p|<1`$. Some elementary properties of $`E`$ are
(2.4)
$$E(x)=xE(1/x)=E(p/x)$$
and the (quasi)periodicity
(2.5)
$$E(x)=(x)^kp^{\left(\genfrac{}{}{0pt}{}{k}{2}\right)}E(xp^k),$$
which follows by iterating (2.4).
Using definition (2.3) one can define an elliptic analogue of the $`q`$-shifted factorial by
(2.6)
$$(a;q,p)_n=\{\begin{array}{cc}_{k=0}^{n1}E(aq^k)\hfill & n>0\hfill \\ 1\hfill & n=0\hfill \\ 1/_{k=0}^{n1}E(aq^{n+k})=1/(aq^n;q,p)_n\hfill & n<0\text{.}\hfill \end{array}$$
Note that $`E(x;0)=1x`$ and hence $`(a;q,0)_n=(a;q)_n`$. Again we use condensed notation, setting
$$(a_1,\mathrm{},a_m;q,p)_n=(a_1;q,p)_n\mathrm{}(a_m;q,p)_n.$$
Many of the relations satisfied by the $`q`$-shifted factorials (see (I.7)–(I.30) of ) trivially generalize to the elliptic case. Here we only list those identities needed later. The proofs merely require manipulation of the definition of $`(a;q,p)_n`$;
(2.7a) $`(aq^n;q,p)_n`$ $`=(q/a;q,p)_n(a/q)^nq^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)}`$
(2.7b) $`(aq^n;q,p)_k`$ $`=(q/a;q,p)_n(a;q,p)_kq^{nk}/(q^{1k}/a;q,p)_n`$
(2.7c) $`(aq^n;q,p)_k`$ $`=(aq^k;q,p)_n(a;q,p)_k/(a;q,p)_n=(aq;q,p)_{n+k}/(a;q,p)_n`$
(2.7d) $`(a;q,p)_{nk}`$ $`=(a;q,p)_n(q^{1n}/a)^kq^{\left(\genfrac{}{}{0pt}{}{k}{2}\right)}/(q^{1n}/a;q,p)_k`$
(2.7e) $`(a;q,p)_{kn}`$ $`=(a,aq,\mathrm{},aq^{k1};q^k,p)_n.`$
Finally we will need the identity
(2.8)
$$(a;q,p)_n=(a)^{nk}p^{n\left(\genfrac{}{}{0pt}{}{k}{2}\right)}q^{k\left(\genfrac{}{}{0pt}{}{n}{2}\right)}(ap^k;q,p)_n,$$
which follows from (2.5) and (2.6).
After these preliminaries we come to Frenkel and Turaev’s definition of balanced, very-well-poised, elliptic (or modular) hypergeometric series ,
(2.9)
$${}_{r+1}{}^{}\omega _{r}^{}(a_1;a_4,\mathrm{},a_{r+1};q,p)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{E(a_1q^{2k})}{E(a_1)}\frac{(a_1,a_4,\mathrm{},a_{r+1};q;p)_kq^k}{(q,a_1q/a_4,\mathrm{},a_1q/a_{r+1};q,p)_k},$$
where $`(a_4\mathrm{}a_{r+1})^2=a_1^{r3}q^{r5}`$. Following we will stay clear of any convergence problems by demanding terminating series, i.e., one of the $`a_i`$ $`(i=4,\mathrm{},r+1)`$ is of the form $`q^n`$ with $`n`$ a nonnegative integer. Remark that by $`E(x;p)E(x;p)=E(x^2;p^2)`$ the above ratio of two elliptic $`E`$-functions can be written as
$$\frac{(qa_1^{1/2},qa_1^{1/2};q,p^{1/2})_k}{(a_1^{1/2},a_1^{1/2};q,p^{1/2})_k}.$$
Hence in the $`p0`$ limit we recover the usual definition of a balanced, very-well-poised, basic hypergeometric series.
An important result of Frenkel and Turaev is the elliptic analogue of Bailey’s transformation (2.1).
###### Theorem 2.1.
Let $`bcdefg=a^3q^{n+2}`$ and $`\lambda =a^2q/bcd.`$ Then
(2.10)
$$\begin{array}{c}{}_{10}{}^{}\omega _{9}^{}(a;b,c,d,e,f,g,q^n;q,p)\hfill \\ \hfill =\frac{(aq,aq/ef,\lambda q/e,\lambda q/f;q,p)_n}{(aq/e,aq/f,\lambda q/ef,\lambda q;q,p)_n}{}_{10}{}^{}\omega _{9}^{}(\lambda ;\lambda b/a,\lambda c/a,\lambda d/a,e,f,g,q^n;q,p).\end{array}$$
Of course we can again specialize $`cd=aq`$ to arrive at an elliptic Jackson sum.
###### Corollary 2.2.
For $`a^2q^{n+1}=bcde`$ there holds
(2.11)
$${}_{8}{}^{}\omega _{7}^{}(a;b,c,d,e,q^n;q,p)=\frac{(aq,aq/bc,aq/bd,aq/cd;q,p)_n}{(aq/b,aq/c,aq/d,aq/bcd;q,p)_n}.$$
## 3. A matrix inverse
Before deriving new summation and transformation formulas for elliptic hypergeometric series we need to prepare several, mostly elementary, results for the elliptic function $`E`$ of equation (2.3). This will result in a matrix inverse that will be at the heart of all results of the subsequent section.
The obvious starting point is the well-known addition formula
(3.1)
$$\begin{array}{c}E(ux)E(u/x)E(vy)E(v/y)E(uy)E(u/y)E(vx)E(v/x)\hfill \\ \hfill =\frac{v}{x}E(xy)E(x/y)E(uv)E(u/v).\end{array}$$
By iterating this equation one readily derives the following lemma, which for $`p=0`$ reduces to a result of Macdonald, first published by Bhatnagar and Milne \[4, Thm. 2.27\].
###### Lemma 3.1.
For $`n`$ a nonnegative integer and $`a_j,b_j,c_j,d_j`$ ($`j=0,\mathrm{},n`$) indeterminates there holds
(3.2)
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}b_k/c_kE(a_kb_k)E(a_k/b_k)E(c_kd_k)E(c_k/d_k)\hfill \\ \hfill \times \underset{j=0}{\overset{k1}{}}E(a_jc_j)E(a_j/c_j)E(b_jd_j)E(b_j/d_j)\underset{j=k+1}{\overset{n}{}}E(a_jd_j)E(a_j/d_j)E(b_jc_j)E(b_j/c_j)\\ \hfill =\underset{j=0}{\overset{n}{}}E(a_jc_j)E(a_j/c_j)E(b_jd_j)E(b_j/d_j)\underset{j=0}{\overset{n}{}}E(a_jd_j)E(a_j/d_j)E(b_jc_j)E(b_j/c_j).\end{array}$$
###### Proof.
We carry out induction on $`n`$. For $`n=0`$ the lemma is nothing but (3.1) with $`u=a_0`$, $`v=b_0`$, $`x=c_0`$ and $`y=d_0`$. Now write (3.2) as $`L_n=R_n`$ and assume this to hold for $`nm1`$. With the abbreviations
$`f_j`$ $`=E(a_jb_j)E(a_j/b_j)E(c_jd_j)E(c_j/d_j)`$
$`g_j`$ $`=E(a_jc_j)E(a_j/c_j)E(b_jd_j)E(b_j/d_j)`$
$`h_j`$ $`=E(a_jd_j)E(a_j/d_j)E(b_jc_j)E(b_j/c_j)`$
we then have
$`L_m`$ $`=h_mL_{m1}+{\displaystyle \frac{b_m}{c_m}}f_m{\displaystyle \underset{j=0}{\overset{m1}{}}}g_j`$
$`=h_m{\displaystyle \underset{j=0}{\overset{m1}{}}}g_j{\displaystyle \underset{j=0}{\overset{m}{}}}h_j+{\displaystyle \frac{b_m}{c_m}}f_m{\displaystyle \underset{j=0}{\overset{m1}{}}}g_j`$
$`={\displaystyle \underset{j=0}{\overset{m}{}}}g_j{\displaystyle \underset{j=0}{\overset{m}{}}}h_j=R_m,`$
where in the second line we have used the induction hypothesis and in the third line the addition formula (3.1) in the form $`h_m+b_mf_m/c_m=g_m`$. ∎
Making the substitutions
$$a_j(abd^2)^{1/2},b_j(ab/c^2)^{1/2},c_j(ab)^{1/2}r^j,d_j(a/b)^{1/2}q^j$$
and using the definition of the elliptic analogue of the $`q`$-shifted factorial (2.6) and the relations (2.7) it follows that
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^kr^k)E(bq^kr^k)}{E(a)E(b)}\frac{(a/c,c/b;q,p)_k(abd,1/d;r,p)_kq^k}{(cr,abr/c;r,p)_k(q/bd,adq;q,p)_k}\hfill \\ \hfill =\frac{E(c)E(ab/c)E(ad)E(bd)}{E(a)E(b)E(cd)E(abd/c)}\left(1\frac{(a/c,bq^n/c;q,p)_{n+1}(abd,dr^n;r,p)_{n+1}}{(bdq^n,ad;q,p)_{n+1}(r^n/c,ab/c;r,p)_{n+1}}\right).\end{array}$$
For $`p=0`$ this corresponds to \[18, Eq. (2.7)\] of Gasper and Rahman. Important will be the specialization obtained by choosing $`d=r^n`$,
(3.3)
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^kr^k)E(bq^kr^k)}{E(a)E(b)}\frac{(a/c,c/b,q,p)_k(abr^n,r^n;r,p)_kq^k}{(cr,abr/c;r,p)_k(qr^n/b,aqr^n;q,p)_k}\hfill \\ \hfill =\frac{E(c)E(ab/c)E(ar^n)E(br^n)}{E(a)E(b)E(cr^n)E(abr^n/c)}.\end{array}$$
We will use this identity in the next section to prove Theorem 4.1. Now it is needed to obtain the following pair of infinite-dimensional, lower-triangular matrices, that are inverses of each other
$`f_{n,k}`$ $`={\displaystyle \frac{(aq^kr^k,q^kr^k/b;q,p)_{nk}}{(r,abr^{2k+1};r,p)_{nk}}}`$
$`f_{n,k}^1`$ $`=(1)^{nk}q^{\left(\genfrac{}{}{0pt}{}{nk}{2}\right)}{\displaystyle \frac{E(aq^kr^k)E(q^kr^k/b)}{E(aq^nr^n)E(q^nr^n/b)}}{\displaystyle \frac{(aq^{k+1}r^n,q^{k+1}r^n/b;q,p)_{nk}}{(r,abr^{n+k};r,p)_{nk}}}.`$
For $`p=0`$ this is \[8, Eqs. (4.4) and (4.5)\], \[16, Eqs. (3.2) and (3.3)\] and \[24, Eq. (4.3)\]. To derive it from (3.3) we follow \[17, Sec. 3.6\] and set $`c=1`$ followed by the replacements $`nnl`$, $`kkl`$, $`aaq^lr^l`$ and $`bbq^lr^l`$. By (2.7) one then finds the desired orthogonality relation
(3.4)
$$\underset{k=l}{\overset{n}{}}f_{n,k}^1f_{k,l}=\delta _{n,l}$$
with $`f`$ and $`f^1`$ as given above. Finally replacing $`rq^r`$ and using (2.7) yields the new inverse pair
(3.5a) $`f_{n,k}`$ $`={\displaystyle \frac{E(abq^{2rk})}{E(ab)}}{\displaystyle \frac{(aq^n;q,p)_{rk}}{(bq^{1n};q,p)_{rk}}}{\displaystyle \frac{(ab,q^{rn};q^r,p)_k}{(q^r,abq^{rn+r};q^r,p)_k}}q^{rk}`$
(3.5b) $`f_{n,k}^1`$ $`={\displaystyle \frac{(b;q,p)_{rn}}{(aq;q,p)_{rn}}}{\displaystyle \frac{E(aq^{(r+1)k})E(bq^{(r1)k})}{E(a)E(b)}}`$
$`\times {\displaystyle \frac{(a,1/b;q,p)_k}{(q^r,abq^r;q^r,p)_k}}{\displaystyle \frac{(abq^{rn},q^{rn};q^r,p)_k}{(q^{1rn}/b,aq^{rn+1};q,p)_k}}q^k.`$
This last pair of inverse matrices will be used repeatedly in the next section. We note that it also follows by the (simultaneous) substitutions $`aab`$, $`b_iaq^i`$ and $`c_iq^{ri}`$ in the following elliptic analogue of a result due to Krattenthaler \[24, Eq (1.5)\].
###### Lemma 3.2.
Let $`a`$ and $`b_i,c_i`$ ($`i`$) be indeterminates (such that $`c_ic_j`$ for $`ij`$ and $`ac_ic_j1`$ for $`i,j)`$. Then (3.4) holds with
$$f_{n,k}=\frac{_{j=k}^{n1}E(c_kb_j)E(ac_k/b_j)}{_{j=k+1}^nc_jE(ac_kc_j)E(c_k/c_j)}$$
and
$$f_{n,k}^1=\frac{E(c_kb_k)E(ac_k/b_k)}{E(c_nb_n)E(ac_n/b_n)}\frac{_{j=k+1}^nE(c_nb_j)E(ac_n/b_j)}{_{j=k}^{n1}c_jE(ac_nc_j)E(c_n/c_j)}.$$
###### Proof.
Since for $`n=l`$ (3.4) clearly holds we may assume $`n>l`$ in the following. Now let $`n>0`$ in (3.2) and make the replacements $`nnl`$, $`kkl`$ and $`a_ja^{1/2}c_l`$, $`b_ja^{1/2}c_{n+l}`$, $`c_jb_{j+l}/a^{1/2}`$, $`d_ja^{1/2}c_{j+l}`$. Noting that, in particular, $`a_0=d_0`$ and $`b_n=d_n`$ so that the right-hand side of (3.2) vanishes, and after performing a few trivialities, one finds (3.4) (with $`n>l`$) with $`f`$ and $`f^1`$ given by Lemma 3.2. ∎
## 4. Summation and transformation formulas
Our approach to elliptic hypergeometric summation and transformation formulas is a standard one in the context of basic hypergeometric series, see e.g., . Given a pair of infinite-dimensional, lower-triangular matrices $`f`$ and $`f^1`$, i.e., a pair of matrices such that (3.4) holds, the following two statements are equivalent
(4.1)
$$\underset{k=0}{\overset{n}{}}f_{n,k}a_k=b_n$$
and
(4.2)
$$\underset{k=0}{\overset{n}{}}f_{n,k}^1b_k=a_n.$$
Our first example of how this may be usefully applied arises by noting that, thanks to (3.3), equation (4.2) holds for the pair of matrices given in (3.5) with
$`a_n`$ $`={\displaystyle \frac{(bq;q,p)_{rn}(c,ab/c;q^r,p)_n}{(a;q,p)_{rn}(abq^r/c,cq^r;q^r,p)_n}}`$
$`b_n`$ $`={\displaystyle \frac{(a/c,c/b,q,p)_n(q^r,abq^r;q^r,p)_n}{(cq^r,abq^r/c;q^r,p)_n(a,1/b;q,p)_n}}.`$
Hence also (4.1) holds leading to
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(abq^{2rk})}{E(ab)}\frac{(bq,aq^n;q,p)_{rk}}{(a,bq^{1n};q,p)_{rk}}\frac{(ab,c,ab/c,q^{rn};q^r,p)_k}{(q^r,abq^r/c,cq^r,abq^{rn+r};q^r,p)_k}q^{rk}\hfill \\ \hfill =\frac{(a/c,c/b,q,p)_n(q^r,abq^r;q^r,p)_n}{(cq^r,abq^r/c;q^r,p)_n(a,1/b;q,p)_n}.\end{array}$$
By (2.7) we may alternatively write this in hypergeometric notation as follows.
###### Theorem 4.1.
For $`r`$ a positive integer there holds
$$\begin{array}{c}{}_{2r+6}{}^{}\omega _{2r+5}^{}(ab;c,ab/c,bq,bq^2,\mathrm{},bq^r,aq^n,aq^{n+1},\mathrm{},aq^{n+r1},q^{rn};q^r,p)\hfill \\ \hfill =\frac{(a/c,c/b,q,p)_n(q^r,abq^r;q^r,p)_n}{(cq^r,abq^r/c;q^r,p)_n(a,1/b;q,p)_n}.\end{array}$$
When $`r=1`$ this corresponds to the specialization $`bc=a`$ of the elliptic Jackson sum (2.11).
With Theorem 4.1 at hand we are prepared for the proof of the following quadratic transformation.
###### Theorem 4.2.
Let $`bcd=aq`$ and $`ef=a^2q^{2n+1}`$. When $`g=a/b`$ or $`g=a/e`$, there holds
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{3k})}{E(a)}\frac{(b,c,d;q,p)_k}{(aq^2/b,aq^2/c,aq^2/d;q^2,p)_k}\frac{(e,f,q^{2n};q^2,p)_k}{(aq/e,aq/f,aq^{2n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq^2,a^2q^2/bce,a^2q^2/bdeg,agq^2/cd;q^2,p)_n}{(a^2q^2/beg,agq^2/c,aq^2/d,a^2q^2/bcde;q^2,p)_n}\\ \hfill \times {}_{10}{}^{}\omega _{9}^{}(ag/c;a/c,gq^2/c,beg/a,d,f,g,q^{2n};q^2,p).\end{array}$$
For $`p=0`$ and $`g=a/b`$ this identity reduces to a transformation of Gasper and Rahman \[18, Eq. (5.14)\]. We should also remark that the left-hand side does not depend on $`g`$ so that the two different cases actually correspond to a transformation of the right-hand side. Indeed, the equality of the $`g=a/b`$ and $`g=a/e`$ instances of the right-hand side is an immediate consequence of the $`{}_{10}{}^{}\omega _{9}^{}`$ transformation (2.10).
###### Proof.
Given the previous remark we only need to prove the $`g=a/b`$ case of the theorem. Starting point is again the pair of inverse matrices (3.5) in which we set $`r=2`$. The crux of the proof is the observation that equation (4.1) holds, with
(4.3)
$$a_n=\frac{(b/d,bdq;q^2,p)_n}{(adq^2,aq/d;q^2,p)_n}{}_{10}{}^{}\omega _{9}^{}(ad;ad/c,c,dq,dq^2,aq/b,abq^{2n},q^{2n};q^2,p)$$
and
(4.4)
$$b_n=\frac{(q^2,abq^2,aq/b;q^2,p)_n(a/c,c/d,dq;q,p)_n}{(a,1/b,bq;q,p)_n(cq^2,adq^2/c,aq/d;q^2,p)_n}.$$
Assuming this is true, we immediately recognize (4.2) as Theorem 4.2 (with $`g=a/b`$) under the simultaneous replacements $`ba/c,cc/d,ddq,eaq/b`$ and $`fabq^{2n}`$.
Of course, it remains to show (4.1) with the above pair of $`a_n`$ and $`b_n`$. Writing $`a_n=_ja_{n,j}`$ in accordance with the definition of $`{}_{10}{}^{}\omega _{9}^{}`$, we start with the trivialities
(4.5)
$$b_n=\underset{k=0}{\overset{n}{}}f_{n,k}a_k=\underset{k=0}{\overset{n}{}}\underset{j=0}{\overset{k}{}}f_{n,k}a_{k,j}=\underset{j=0}{\overset{n}{}}\underset{k=j}{\overset{n}{}}f_{n,k}a_{k,j}=\underset{j=0}{\overset{n}{}}\underset{k=0}{\overset{nj}{}}f_{n,k+j}a_{k+j,j}.$$
Using the explicit expressions for $`f_{n,k}`$ and $`a_{n,k}`$ as well as the relations in (2.7) this becomes
$$\begin{array}{c}b_n=\underset{j=0}{\overset{n}{}}\frac{(dq,aq^n;q,p)_{2j}(abq^2;q^2,p)_{2j}(ad,c,ad/c,aq/b,q^{2n};q^2,p)_j(bq^2/d)^j}{(a,bq^{1n};q,p)_{2j}(ad;q^2,p)_{2j}(q^2,adq^2/c,cq^2,aq/d,abq^{2n+2};q^2,p)_j}\hfill \\ \hfill \times {}_{8}{}^{}\omega _{7}^{}(abq^{4j};b/d,bdq^{2j+1},aq^{n+2j},aq^{n+2j+1},q^{2n+2j};q^2,p).\end{array}$$
By the elliptic Jackson sum (2.11) we can sum the $`{}_{8}{}^{}\omega _{7}^{}`$, and after the usual simplifications we find
$$\begin{array}{c}b_n=\frac{(abq^2,aq/b;q^2,p)_n}{(adq^2,aq/d;q^2,p)_n}\frac{(1/d,dq;q,p)_n}{(1/b,bq;q,p)_n}\hfill \\ \hfill \times {}_{10}{}^{}\omega _{9}^{}(ad;c,ad/c,dq,dq^2,aq^n,aq^{n+1},q^{2n};q^2,p).\end{array}$$
By the $`r=2`$ case of Theorem 4.1 the $`{}_{10}{}^{}\omega _{9}^{}`$ can be summed yielding the expression for $`b_n`$ given in (4.4). ∎
A result very similar to that of Theorem 4.2 is the following cubic transformation, which, for $`f=a/b`$, provides an elliptic analogue of \[18, Eq. (3.6)\] by Gasper and Rahman.
###### Theorem 4.3.
Let $`bcd=aq`$ and $`de=a^2q^{3n+1}`$. Then for $`f=a/b`$ or $`f=a/e`$ there holds
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q,p)_k}{(aq^3/b,aq^3/c;q^3,p)_k}\frac{(d;q,p)_{2k}}{(aq/d;q,p)_{2k}}\frac{(e,q^{3n};q^3,p)_k}{(aq/e,aq^{3n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq^3,a^2q^3/bce,a^2q^3/bdef,afq^3/cd;q^3,p)_n}{(a^2q^3/bef,aq^3f/c,aq^3/d,a^2q^3/bcde;q^3,p)_n}\\ \hfill \times {}_{10}{}^{}\omega _{9}^{}(af/c;a/c,fq^3/c,bef/a,d,dq,f,q^{3n};q^3,p).\end{array}$$
Again we note that the two different cases correspond to the $`{}_{10}{}^{}\omega _{9}^{}`$ transformation (2.10) applied to the right-hand side.
###### Proof.
By the above remark we only need a proof for $`f=a/b`$. The claim is now that if we choose $`r=3`$ in the pair of matrices (3.5) then (4.1) holds, with
(4.6)
$$a_n=\frac{(b^2/a;q^3,p)_n}{(a^2q^3/b;q^3,p)_n}{}_{10}{}^{}\omega _{9}^{}(a^2/b;ac/b,a/c,aq/b,aq^2/b,aq^3/b,abq^{3n},q^{3n};q^3,p)$$
and
(4.7)
$$b_n=\frac{(q^3,abq^3;q^3,p)_n(b/c,c;q,p)_n(aq/b;q,p)_{2n}}{(a,1/b;q,p)_n(acq^3/b,aq^3/c;q^3,p)_n(bq;q,p)_{2n}}.$$
Clearly, if this is true we are done with the proof since with these $`a_n`$ and $`b_n`$ equation (4.2) corresponds to the $`f=b/a`$ case of Theorem 4.3 with the replacements $`bb/c,daq/b`$ and $`eabq^{3n}`$.
To establish (4.1) with the above $`a_n`$ and $`b_n`$ we follow the proof of Theorem 4.2. That is, we again write $`a_n=_ja_{n,j}`$ and use (4.5). Inserting the expressions for $`f_{n,k}`$ and $`a_{n,k}`$ this yields
$$\begin{array}{c}b_n=\underset{j=0}{\overset{n}{}}\frac{(aq/b,aq^n;q,p)_{3j}(abq^3;q^3,p)_{2j}(a^2/b,ac/b,a/c,q^{3n};q^3,p)_j(b^2q^3/a)^j}{(a,bq^{1n};q,p)_{3j}(a^2/b;q^3,p)_{2j}(q^3,aq^3/c,acq^3/b,abq^{3n+3};q^3,p)_j}\hfill \\ \hfill \times {}_{8}{}^{}\omega _{7}^{}(abq^{6j};b^2/a,aq^{n+3j},aq^{n+3j+1},aq^{n+3j+2},q^{3n+3j};q^3,p).\end{array}$$
The $`{}_{8}{}^{}\omega _{7}^{}`$ can be summed by (2.11), and after some manipulations involving (2.7) we arrive at
$$\begin{array}{c}b_n=\frac{(abq^3;q^3,p)_n}{(a^2q^3/b;q^3,p)_n}\frac{(b/a;q,p)_n}{(1/b;q,p)_n}\frac{(aq/b;q,p)_{2n}}{(bq;q,p)_{2n}}\hfill \\ \hfill \times {}_{12}{}^{}\omega _{11}^{}(a^2/b;ac/b,a/c,aq/b,aq^2/b,aq^3/b,aq^n,aq^{n+1},aq^{n+2},q^{3n};q^3,p).\end{array}$$
According to Theorem 4.1 with $`r=3`$ the $`{}_{12}{}^{}\omega _{11}^{}`$ can be summed to yield (4.7) as claimed. ∎
Theorems 4.2 and 4.3 imply several other quadratic and cubic summation and transformation formulas.
The most obvious ones arise when we demand that $`g=1`$ in Theorem 4.2 or $`f=1`$ in Theorem 4.3.
###### Corollary 4.4.
Let $`bcd=aq`$ and $`ef=a^2q^{2n+1}`$. When $`b=a`$ or $`e=a`$ there holds
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{3k})}{E(a)}\frac{(b,c,d;q,p)_k}{(aq^2/b,aq^2/c,aq^2/d;q^2,p)_k}\frac{(e,f,q^{2n};q^2,p)_k}{(aq/e,aq/f,aq^{2n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq^2,a^2q^2/bce,a^2q^2/bde,aq^2/cd;q^2,p)_n}{(a^2q^2/be,aq^2/c,aq^2/d,a^2q^2/bcde;q^2,p)_n}.\end{array}$$
For $`p=0`$ and $`b=a`$ this is a summation of Gessel and Stanton \[19, Eq. (1.4)\], and for $`p=0`$, $`e=a`$ it corresponds to \[33, Eq. (1.9), $`bq^{2n}`$\] by Rahman and \[10, Eq. (5.1d)\] by Chu.
###### Corollary 4.5.
Let $`bcd=aq`$ and $`de=a^2q^{3n+1}`$. When $`b=a`$ or $`e=a`$ there holds
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q,p)_k}{(aq^3/b,aq^3/c;q^3,p)_k}\frac{(d;q,p)_{2k}}{(aq/d;q,p)_{2k}}\frac{(e,q^{3n};q^3,p)_k}{(aq/e,aq^{3n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq^3,a^2q^3/bce,a^2q^3/bde,aq^3/cd;q^3,p)_n}{(a^2q^3/be,aq^3/c,aq^3/d,a^2q^3/bcde;q^3,p)_n}.\end{array}$$
When $`b=a`$ this is the elliptic analogue of \[16, Eq. (5.22), $`cq^{3n}`$\] of Gasper. Theorem 4.3 also leads to a summation formula if we choose $`d=a`$. Indeed, $`{}_{10}{}^{}\omega _{9}^{}(af/c;a/c,fq^3/c,bef/a,a,aq,f,q^{3n};q^3,p)`$ with $`f=a/b`$, $`bc=q`$ and $`e=aq^{3n+1}`$ becomes $`{}_{7}{}^{}\omega _{6}^{}(a^2/q;a/b,ab/q,aq^{3n+1},q^{3n};q^3,p)`$ which, by (2.11), evaluates to
$$\frac{(q^3,a^2q^2,bq,q^2/b;q^3,p)_n}{(aq,q^2/a,aq^3/b,abq^2;q^3,p)_n}.$$
We therefore conclude the following result, which for $`p=0`$ corresponds to \[18, Eq. (3.7)\].
###### Corollary 4.6.
For $`bc=q`$ and $`e=aq^{3n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q,p)_k}{(aq^3/b,aq^3/c;q^3,p)_k}\frac{(a;q,p)_{2k}}{(q;q,p)_{2k}}\frac{(e,q^{3n};q^3,p)_k}{(aq/e,aq^{3n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq^2,aq^3,bq,cq;q^3,p)_n}{(q,q^2,aq^3/b,aq^3/c;q^3,p)_n}.\end{array}$$
The next three results, stated as separate theorems, are somewhat less trivial as their proof deviates from the standard polynomial argument applicable in the $`p=0`$ case.
###### Theorem 4.7.
For $`bcd=a^2q`$ and $`ef=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{3k})}{E(a)}\frac{(b,c,d;q^2,p)_k}{(aq/b,aq/c,aq/d;q,p)_k}\frac{(e,f,q^n;q,p)_k}{(aq^2/e,aq^2/f,aq^{n+2};q^2,p)_k}q^k\hfill \\ \hfill \frac{(aq,aq/bc;q,p)_n(aq^{1n}/b,aq^{1n}/c;q^2,p)_n}{(aq/b,aq/c;q,p)_n(aq^{1n},aq^{1n}/bc;q^2,p)_n}\\ \hfill \times {}_{10}{}^{}\omega _{9}^{}(a^2/ef;b,c,d,a/e,a/f,q^{1n},q^n;q^2,p).\end{array}$$
For $`p=0`$ this is \[18, Eq. (5.15)\] (with corrected misprint).
###### Theorem 4.8.
For $`bc=a^2q^{n+1}`$ and $`de=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n/2}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q^3,p)_k}{(aq/b,aq/c;q,p)_k}\frac{(q^n;q,p)_{2k}}{(aq^{n+1};q,p)_{2k}}\frac{(d,e;q,p)_k}{(aq^3/d,aq^3/e;q^3,p)_k}q^k\hfill \\ \hfill =\frac{(aq;q,p)_n(aq^{2n}/b;q^3,p)_n}{(aq/b;q,p)_n(aq^{2n};q^3,p)_n}{}_{10}{}^{}\omega _{9}^{}(a^2/de;b,c,a/d,a/e,q^{2n},q^{1n},q^n;q^3,p).\end{array}$$
When $`p=0`$ this can be recognized as \[18, Eq. (3.19), $`cbq^{n1}`$\].
###### Theorem 4.9.
For $`bcd=a^2q`$ and $`de=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q^3,p)_k}{(aq/b,aq/c;q,p)_k}\frac{(d;q,p)_{2k}}{(aq/d;q,p)_{2k}}\frac{(e,q^n;q,p)_k}{(aq^3/e,aq^{n+3};q^3,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}f_n{}_{10}{}^{}\omega _{9}^{}(a^2/deq;a/dq,a/e,b,c,d,q^{1n},q^n;q^3,p)\hfill & n2(mod3)\hfill \\ g_n{}_{10}{}^{}\omega _{9}^{}(a^2/de;a/d,a/e,b,c,dq,q^{2n},q^n;q^3,p)\hfill & n1(mod3)\hfill \\ h_n{}_{10}{}^{}\omega _{9}^{}(a^2q/de;aq/d,a/e,b,c,dq^2,q^{2n},q^{1n};q^3,p)\hfill & n0(mod3),\hfill \end{array}\end{array}$$
with
$`f_n`$ $`={\displaystyle \frac{(aq^{3\sigma },aq^{3\sigma }/bc,aq^{3\sigma }/bd,aq^{3\sigma }/cd;q^3,p)_{(n+\sigma )/3}}{(aq^{3\sigma }/b,aq^{3\sigma }/c,aq^{3\sigma }/d,aq^{3\sigma }/bcd;q^3,p)_{(n+\sigma )/3}}}`$
$`g_n`$ $`={\displaystyle \frac{(aq^{3\sigma },aq^{3\sigma }/bc,aq^{2\sigma }/bd,aq^{2\sigma }/cd;q^3,p)_{(n+\sigma )/3}}{(aq^{3\sigma }/b,aq^{3\sigma }/c,aq^{2\sigma }/d,aq^{2\sigma }/bcd;q^3,p)_{(n+\sigma )/3}}}`$
$`h_n`$ $`={\displaystyle \frac{E(aq^\sigma )E(aq^\sigma /bc)E(aq/bd)E(aq/cd)}{E(aq^\sigma /b)E(aq^\sigma /c)E(aq/d)E(aq/bcd)}}`$
$`\times {\displaystyle \frac{(aq^{3\sigma },aq^{3\sigma }/bc,aq^{1\sigma }/bd,aq^{1\sigma }/cd;q^3,p)_{(n+\sigma )/3}}{(aq^{3\sigma }/b,aq^{3\sigma }/c,aq^{1\sigma }/d,aq^{1\sigma }/bcd;q^3,p)_{(n+\sigma )/3}}}`$
where $`\sigma \{0,1,2\}`$ is fixed by $`n+\sigma 0(mod3)`$.
###### Proof.
Using identity (2.8) it readily follows that both the left- and right-hand sides of the identities in Theorems 4.74.9, viewed as functions of the variable $`b`$, satisfy the periodicity $`f(pb)=f(b)`$. If we define $`h(b)=\text{LHS}(b)/\text{RHS}(b)1`$ then $`h`$ is a meromorphic function in $`0<|b|<\mathrm{}`$ with that same periodicity and with a finite number of poles in a period annulus. Such a function is either a constant or has an equal number of zeros and poles in a period annulus (poles or zeros of order $`j`$ counted $`j`$ times). So if we can show that within a period annulus $`h(b)=0`$ for an infinite number of $`b`$ then $`h`$ must be identically zero. Without loss of generality we may assume that $`q^{m_1}p^{m_2}`$ for $`m_1`$ and $`m_2`$ positive integers. It is then enough to show that the identities in Theorems 4.74.9 hold for $`b=q^m`$ where $`m`$ runs over an infinite subset of the integers. First consider Theorem 4.7. For $`b=q^{2m}`$ ($`m`$ a nonnegative integer) it holds as can be seen by making the simultaneous replacements $`nm,be,cf`$, $`dq^n,ec`$, $`fd`$ in the $`g=a/b`$ case of Theorem 4.2 and using (2.7). Theorem 4.8 for $`b=q^{3m}`$ ($`m`$ a nonnegative integer) follows by making the simultaneous replacements $`nm,bd,ce`$, $`dq^n,ec`$ in the $`f=a/b`$ case of Theorem 4.3 and using (2.7). Finally we show that the first ($`n2(mod3)`$) of the identities of Theorem 4.9 holds for $`b=q^{3m}`$ ($`m`$ a nonnegative integer). The other two follow in similar manner. Take Theorem 4.3 with $`f=a/b`$ and make the simultaneous replacements $`nm`$, $`be`$, $`cq^n`$, $`ec`$. To the thus obtained identity apply the elliptic Bailey transformation (2.10) with $`aa^2q^n/e`$, $`baq^n`$, $`caq^{n+3}/e`$, $`ddq`$, $`ec`$, $`fd`$ and $`ga/e`$ (so that $`\lambda =a^2/deq`$). The result is the first identity of the theorem with $`b=q^{3m}`$. ∎
By appropriately specializing the transformations in the Theorems 4.74.9 we obtain several further summations. Taking $`f=a`$ in Theorem 4.7 leads to an elliptic extension of \[16, Eq. (5.15), $`bq^n`$\] and \[33, Eq. (1.9), $`cq^n`$\].
###### Corollary 4.10.
For $`bcd=a^2q`$ and $`e=q^{n+1}`$,
$`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{E(aq^{3k})}{E(a)}}`$ $`{\displaystyle \frac{(b,c,d;q^2,p)_k}{(aq/b,aq/c,aq/d;q,p)_k}}{\displaystyle \frac{(a,e,q^n;q,p)_k}{(q^2,aq^2/e,aq^{n+2};q^2,p)_k}}q^k`$
$`={\displaystyle \frac{(aq,aq/bc;q,p)_n(aq^{1n}/b,aq^{1n}/c;q^2,p)_n}{(aq/b,aq/c;q,p)_n(aq^{1n},aq^{1n}/bc;q^2,p)_n}}`$
$`={\displaystyle \frac{(aq^{2\sigma },aq^{2\sigma }/bc,aq^{2\sigma }/bd,aq^{2\sigma }/cd;q^2,p)_{(n+\sigma )/2}}{(aq^{2\sigma }/b,aq^{2\sigma }/c,aq^{2\sigma }/d,aq^{2\sigma }/bcd;q^2,p)_{(n+\sigma )/2}}},`$
where $`\sigma \{0,1\}`$ is determined by $`n+\sigma 0(mod2)`$.
To obtain a summation formula by choosing $`d=a`$ in Theorem 4.7 requires a bit more work. First observe that by this choice for $`d`$ we have $`bc=aq`$ so that
$$\frac{(aq/bc;q,p)_n}{(aq^{1n}/bc;q^2,p)_n}=\{\begin{array}{cc}\frac{(q;q^2,p)_{n/2}}{(q^n;q^2,p)_{n/2}}\hfill & \text{for }n\text{ even}\hfill \\ 0\hfill & \text{for }n\text{ odd.}\hfill \end{array}$$
Next note that $`{}_{10}{}^{}\omega _{9}^{}(a^2/ef;a,b,c,a/e,a/f,q^{1n},q^n;q^2,p)`$ with $`ef=aq^{n+1}`$ and $`bc=aq`$ reduces to $`{}_{8}{}^{}\omega _{7}^{}(a^2/ef;b,c,a/e,a/f,q^n;q^2,p)`$, which for $`n`$ even can be summed by (2.11) to give
$$\frac{(q^2,q/a,bq/e,aq^2/be;q^2,p)_{n/2}}{(q/e,aq^2/e,q^2/b,bq/a;q^2,p)_{n/2}}.$$
Combining all of the above and using (2.7) we get the following generalization of summation \[19, Eq. (6.14)\] by Gessel and Stanton.
###### Corollary 4.11.
For $`bc=aq`$ and $`de=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{3k})}{E(a)}\frac{(a,b,c;q^2,p)_k}{(q,aq/b,aq/c;q,p)_k}\frac{(d,e,q^n;q,p)_k}{(aq^2/d,aq^2/e,aq^{n+2};q^2,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(aq^2,aq^2/bc,aq^2/bd,aq^2/cd;q^2,p)_{n/2}}{(aq^2/b,aq^2/c,aq^2/d,aq^2/bcd;q^2,p)_{n/2}}\hfill & \text{for }n\text{ even}\hfill \\ 0\hfill & \text{for }n\text{ odd.}\hfill \end{array}\end{array}$$
In exactly the same manner we can derive two summations from Theorem 4.8. First when $`e=a`$ we get the elliptic analogue of \[16, Eq. (5.22), $`bq^{n+1}`$\].
###### Corollary 4.12.
For $`bc=a^2q^{n+1}`$ and $`d=q^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n/2}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q^3,p)_k}{(aq/b,aq/c;q,p)_k}\frac{(q^n;q,p)_{2k}}{(aq^{n+1};q,p)_{2k}}\frac{(a,d;q,p)_k}{(aq^3/d,q^3;q^3,p)_k}q^k\hfill \\ \hfill =\frac{(aq;q,p)_n(aq^{2n}/b;q^3,p)_n}{(aq/b;q,p)_n(aq^{2n};q^3,p)_n}.\end{array}$$
The second summation follows from $`c=a`$. Then $`b=aq^{n+1}`$ and
$$\frac{(aq^{2n}/b;q^3,p)_n}{(aq/b;q,p)_n}=\frac{(q^{12n};q^3,p)_n}{(q^n;q,p)_n}=0\text{for }n2(mod3)\text{.}$$
Further observe that $`{}_{10}{}^{}\omega _{9}^{}(a^2/de;a,b,a/d,a/e,q^{2n},q^{1n},q^n;q^3,p)`$ with $`b=aq^{n+1}`$ and $`de=aq^{n+1}`$ reduces to $`{}_{8}{}^{}\omega _{7}^{}(a^2/de;a/d,a/e,b,q^n,q^{1n};q^3,p)`$, which for $`n2(mod3)`$ can be summed to give
$$\frac{(aq^{2n},q^3,dq^{12n}/a,q^{2n}/d;q^3,p)_{n/3}}{(q^{2n}/a,q^{12n},aq^3/d,dq^{2n};q^3,p)_{n/3}}.$$
After a few simplification we arrive at the following summation.
###### Corollary 4.13.
For $`b=aq^{n+1}`$ and $`cd=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n/2}{}}\frac{E(aq^{4k})}{E(a)}\frac{(a,b;q^3,p)_k}{(q,aq/b;q,p)_k}\frac{(q^n;q,p)_{2k}}{(aq^{n+1};q,p)_{2k}}\frac{(c,d;q,p)_k}{(aq^3/c,aq^3/d;q^3,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(aq^3,aq^3/bc,aq^3/bd;q^3,p)_{n/3}}{(aq^3/c,aq^3/d,aq^3/bcd;q^3,p)_{n/3}}\hfill & n2(mod3)\hfill \\ 0\hfill & n2(mod3).\hfill \end{array}\end{array}$$
Finally we turn to Theorem 4.9. The choice $`e=a`$ immediately gives an elliptic analogue of \[16, Eq. (5.22), $`bq^n`$\].
###### Corollary 4.14.
For $`bcd=a^2q`$ and $`d=q^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q^3,p)_k}{(aq/b,aq/c;q,p)_k}\frac{(d;q,p)_{2k}}{(aq/d;q,p)_{2k}}\frac{(a,q^n;q,p)_k}{(q^3,aq^{n+3};q^3,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(aq^3,aq^3/bc,aq^3/bd,aq^3/cd;q^3,p)_{n/3}}{(aq^3/b,aq^3/c,aq^3/d,aq^3/bcd;q^3,p)_{n/3}}\hfill & n0(mod3)\hfill \\ \frac{(aq,aq/bc,aq/bd,aq/cd;q^3,p)_{(n+2)/3}}{(aq/b,aq/c,aq/d,aq/bcd;q^3,p)_{(n+2)/3}}\hfill & n1(mod3)\hfill \\ \frac{(aq^2,aq^2/bc,aq/bd,aq/cd;q^3,p)_{(n+1)/3}}{(aq^2/b,aq^2/c,aq/d,aq/bcd;q^3,p)_{(n+1)/3}}\hfill & n2(mod3).\hfill \end{array}\end{array}$$
If, on the other hand, we set $`c=a`$ in Theorem 4.9 and perform a calculation similar to the one employed in the derivation of Corollaries 4.11 and 4.13 we get the elliptic extension of \[10, Eq. (4.6d)\].
###### Corollary 4.15.
For $`bc=aq`$ and $`cd=aq^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(a,b;q^3,p)_k}{(q,aq/b;q,p)_k}\frac{(c;q,p)_{2k}}{(aq/c;q,p)_{2k}}\frac{(d,q^n;q,p)_k}{(aq^3/d,aq^{n+3};q^3,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(q,q^2,aq^3,b^2/a;q^3,p)_{n/3}}{(bq,bq^2,b/a,aq^3/b;q^3,p)_{n/3}}\hfill & n0(mod3)\hfill \\ 0\hfill & n0(mod3).\hfill \end{array}\end{array}$$
Similarly, taking $`d=a`$ in Theorem 4.9 yields the last summation of this section.
###### Corollary 4.16.
For $`bc=aq`$ and $`d=q^{n+1}`$,
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{4k})}{E(a)}\frac{(b,c;q^3,p)_k}{(aq/b,aq/c;q,p)_k}\frac{(a;q,p)_{2k}}{(q;q,p)_{2k}}\frac{(d,q^n;q,p)_k}{(aq^3/d,aq^{n+3};q^3,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(aq^3,q^2/b,q^2/c;q^3,p)_{n/3}}{(q^2/bc,aq^3/b,aq^3/c;q^3,p)_{n/3}}\hfill & n0(mod3)\hfill \\ 0\hfill & n1(mod3)\hfill \\ \frac{(aq^2,q/b,q/c;q^3,p)_{(n+2)/3}}{(q/bc,aq^2/b,aq^2/c;q^3,p)_{(n+2)/3}}\hfill & n2(mod3).\hfill \end{array}\end{array}$$
For $`p=0`$ this is \[18, Eq. (3.21), $`aq^n`$\].
Before concluding this section let us remark that there is a nice corollary to Corollary 4.4 in the form of a nontrivial determinant evaluation. For $`p=0`$ this was observed by Andrews and Stanton and the following theorem is a direct elliptic extension to their \[1, Thm. 8\]. (In fact Andrews and Stanton did not use the $`p=0`$ version of Corollary 4.4 but of the $`p=0`$, $`n`$ odd instance of Corollary 4.11 in their proof.)
###### Theorem 4.17.
For $`x,y`$ indeterminates and $`n`$ a positive integer there holds
$$\begin{array}{c}\underset{1i,jn}{det}\left(\frac{(yq^{1i}/x,q^{2i}/xy,q^{24i}/x^2;q^2,p)_{ij}}{(q^{22i}/xy,yq^{12i}/x,q^{i+1};q,p)_{ij}}\right)\hfill \\ \hfill =\underset{i=1}{\overset{n}{}}\frac{(q,x^2q^{2i2};q,p)_i}{(q,x^2q^{2i2};q^2,p)_i}\frac{(xyq^{i1},xq^i/y;q^2,p)_i}{(xyq^{i1},xq^i/y;q,p)_i}.\end{array}$$
Let us note that for $`i<j`$ we need the elliptic analogue of the $`q`$-shifted factorial (2.6) with negative subscript. Hence $`1/(q^{i+1};q,p)_{ij}=(q^{2ij+1};q,p)_{ji}=0`$ for $`2ij+10`$.
###### Proof.
To prove the theorem we establish the Gauss decomposition or “LU” factorization of the matrix $`M`$ featuring in the determinant. Let $`M_n=(M_{i,j})_{1i,jn}`$ with
$$M_{i,j}=\frac{(yq^{1i}/x,q^{2i}/xy,q^{24i}/x^2;q^2,p)_{ij}}{(q^{22i}/xy,yq^{12i}/x,q^{i+1};q,p)_{ij}},$$
and $`U_n=(U_{i,j})_{1i,jn}`$, with
$$U_{i,j}=(1)^{i+j}q^{(ij)(i+j7)/2}\frac{E(x^2q^{3i2})(q^i;q,p)_{2j2i}(q^{33j}/x^2;q,p)_{ji}}{E(x^2q^{i+2j2})(q^2,q^{44j}/x^2,q^{32j}/x^2;q^2,p)_{ji}}$$
for $`ij`$, and
$$U_{i,j}=0$$
for $`i>j`$. Next we calculate the product of the above two matrices
$$\begin{array}{c}(M_nU_n)_{i,j}=\underset{k=1}{\overset{j}{}}M_{i,k}U_{k,j}\hfill \\ \hfill =\frac{(yq^{1i}/x,q^{2i}/xy,q^{24i}/x^2;q^2,p)_{i1}(q^j,q^{1j},q^{33j}/x^2;q,p)_{j1}q^{j1}}{(q^{22i}/xy,yq^{12i}/x,q^{i+1};q,p)_{i1}(q^{44j}/x^2,q^{12j}/x^2,q^2;q^2,p)_{j1}}\\ \hfill \times \underset{k=0}{\overset{j1}{}}\frac{E(x^2q^{3k+1})}{E(x^2q)}\frac{(xq^{i+1}/y,xyq^i,q^{12i};q,p)_k(x^2q,x^2q^{2j},q^{22j};q^2,p)_kq^k}{(xyq^{2i},xq^{3i}/y,x^2q^{2i+2};q^2,p)_k(q,q^{22j},x^2q^{2j};q,p)_k}.\end{array}$$
To proceed we observe that the sum over $`k`$ can be carried out by the $`e=a`$ case of Corollary 4.4 so that the last line in the above equation may be replaced by
$$\frac{(x^2q^3,xq^{i+2}/y,xyq^{i+1},q^{22i};q^2,p)_{j1}}{(q,xyq^{2i},xq^{3i}/y,x^2q^{2i+2};q^2,p)_{j1}}.$$
We learn two things from this result. First, that the matrix $`L_n=(L_{i,j})_{1i,jn}=M_nU_n`$ is lower triangular (since $`(q^{22i};q^2,p)_{j1}=0`$ for $`j>i`$) and, second, that its diagonal entries are given by
$$L_{i,i}=\frac{(q^2,x^2q^{2i1};q,p)_{i1}(xyq^{i+1},xq^{i+2}/y;q^2,p)_{i1}}{(q^3,x^2q^{2i};q^2,p)_{i1}(xyq^i,xq^{i+1}/y;q,p)_{i1}}.$$
The calculation of $`det(M_n)`$ is now done; by $`M_nU_n=L_n`$ we get $`det(M_n)det(U_n)=det(L_n)`$, but $`det(U_n)=1`$ by the fact that $`U_n`$ is an upper-triangular matrix with $`1`$’s along the diagonal. Hence we only need to compute the determinant of $`L_n`$ which is the product of its diagonal entries, resulting in the right-hand side of the theorem. ∎
## 5. An elliptic C<sub>n</sub> Jackson sum
Building on earlier work in , Schlosser proved a multidimensional extension of Jackson’s $`{}_{8}{}^{}\varphi _{7}^{}`$ summation . Here we show that by a generalization of a determinant lemma of Krattenthaler (see Lemmas 5.2 and 5.3 below) Schlosser’s C<sub>n</sub> Jackson sum can readily be generalized to the elliptic case. This is the content of our next theorem.
###### Theorem 5.1.
For $`x_1,\mathrm{},x_n`$, $`a,b,c,d`$ and $`e`$ indeterminates and $`N`$ a nonnegative integer such that $`a^2q^{Nn+2}=bcde`$ there holds
(5.1)
$$\begin{array}{c}\underset{k_1,\mathrm{},k_n=0}{\overset{N}{}}\underset{1i<jn}{}\left(\frac{E(q^{k_ik_j}x_i/x_j)}{E(x_i/x_j)}\frac{E(ax_ix_jq^{k_i+k_j})}{E(ax_ix_jq^N)}\right)\hfill \\ \hfill \times \underset{i=1}{\overset{n}{}}\frac{E(ax_i^2q^{2k_i})}{E(ax_i^2)}\frac{(ax_i^2,bx_i,cx_i,dx_i,ex_i,q^N;q,p)_{k_i}q^{ik_i}}{(q,aqx_i/b,aqx_i/c,aqx_i/d,aqx_i/e,ax_i^2q^{N+1};q,p)_{k_i}}\\ \hfill =\underset{i=1}{\overset{n}{}}\frac{(aqx_i^2,aq^{2i}/bc,aq^{2i}/bd,aq^{2i}/cd;q,p)_N}{(aq^{2n}/bcdx_i,aqx_i/b,aqx_i/c,aqx_i/d;q,p)_N}.\end{array}$$
As remarked above, to prove this result we need the elliptic analogue of the following determinant lemma due to Krattenthaler \[23, Lemma 34\] (see also \[25, Lemma 5\]), which was crucial in the proof of the $`p=0`$ case of (5.1.
###### Lemma 5.2.
Let $`X_1,\mathrm{},X_n,A_2,\mathrm{},A_n`$ and $`C`$ be indeterminates. If, for $`j=0,\mathrm{},n1`$, $`P_j`$ is a Laurent polynomial of degree less than or equal to $`j`$ such that $`P_j(C/X)=P_j(X)`$, then
$$\begin{array}{c}\underset{1i,jn}{det}\left(P_{j1}(X_i)\underset{k=j+1}{\overset{n}{}}(1A_kX_i)(1CA_k/X_i)\right)\hfill \\ \hfill =\underset{1i<jn}{}A_jX_j(1X_i/X_j)(1C/X_iX_j)\underset{i=1}{\overset{n}{}}P_{i1}(1/A_i).\end{array}$$
Here the degree of a Laurent polynomial $`P(x)=_{i=M}^Na_ix^i`$ with $`a_N0`$ is defined to be $`N`$, and the empty product $`_{k=j+1}^n(1A_kX_i)(1CA_k/X_i)`$ for $`j=n`$ is defined to be $`1`$. For a proof of this lemma we refer to .
The needed elliptic analogue of the previous lemma can be stated as follows.
###### Lemma 5.3.
Let $`X_1,\mathrm{},X_n,A_2,\mathrm{},A_n`$ and $`C`$ be indeterminates and $`E`$ the elliptic function defined in (2.3). If, for $`j=0,\mathrm{},n1`$, $`P_j`$ is analytic in $`0<|x|<\mathrm{}`$ with periodicity $`P_j(px)=(C/x^2p)^jP_j(x)`$ and symmetry $`P_j(C/x)=P_j(x)`$, then
(5.2)
$$\begin{array}{c}\underset{1i,jn}{det}\left(P_{j1}(X_i)\underset{k=j+1}{\overset{n}{}}E(A_kX_i)E(CA_k/X_i)\right)\hfill \\ \hfill =\underset{1i<jn}{}A_jX_jE(X_i/X_j)E(C/X_iX_j)\underset{i=1}{\overset{n}{}}P_{i1}(1/A_i).\end{array}$$
###### Proof.
View both sides of (5.2) as a function of the variable $`X_i`$ $`(i=1,\mathrm{},n)`$, and write $`L(X_i)`$ ($`R(X_i)`$) for the left(right)-hand side. From the periodicity property (2.5) and the periodicity of $`P_j`$, we find that
$$F(X_i)=(pX_i^2/C)^{n1}F(pX_i),$$
where $`F=L,R`$. As a result the function $`f`$, defined as the ratio of $`L`$ over $`R`$, satisfies the periodicity $`f(X_i)=f(pX_i)`$. Since $`E(x)`$ and $`P_j(x)`$ are analytic in $`0<|x|<\mathrm{}`$, the only possible poles of $`f`$ are the zero’s of $`R`$. Since $`E(x)`$ has simple zeros at $`x=p^k`$ ($`k`$), the zeros of $`R`$ are $`X_i=p^kX_j`$ and $`X_i=p^kC/X_j`$ where $`k`$ and $`j=1,\mathrm{},i1,i+1,\mathrm{},n`$. First consider $`X_i=p^kX_j`$. When inserted into the determinant it follows from (2.5) and
(5.3)
$$P_j(x)=(x^2p^k/C)^{jk}P_j(xp^k),$$
that the $`i`$-th and $`j`$-th row become proportional (with proportionality constant $`(Cp^k/X_j^2)^{k(n1)}`$). Next, when $`X_i=p^kC/X_j`$ it follows from (2.5), (5.3) and the symmetry $`P_j(C/x)=P_j(x)`$ that the $`i`$-th and $`j`$-th row once again become proportional (with proportionality constant $`(X_j^2p^k/C)^{k(n1)}`$). We may therefore conclude that $`L`$ vanishes at the zeros of $`R`$, so that, according to Liouville’s theorem, $`f`$ must be constant.
To conclude the proof we only need to show the validity of (5.2) for some appropriately chosen values of $`X_1,\mathrm{},X_n`$. A good choice is
(5.4)
$$X_i=1/A_i,i=1,\mathrm{},n.$$
Since $`_{k=j+1}^nE(X_i/A_k)=0`$ for $`j<i`$, this leaves the determinant of an upper-triangular matrix which evaluates to
$$\underset{i=1}{\overset{n}{}}\underset{j=i+1}{\overset{n}{}}E(A_j/A_i)E(CA_iA_j)P_{i1}(1/A_i).$$
Clearly this corresponds to the right-hand side of (5.2) under the specialization (5.4), and we are done. ∎
Choosing $`A_i=Aq^{ni}`$ and $`P_i(X)=(BXq^{ni1},BCq^{ni1}/X;q,p)_i`$ in (5.2), and using (2.7) and $`_{j=1}^n(j1)(nj)=\left(\genfrac{}{}{0pt}{}{n}{3}\right)`$, we obtain the following nice corollary of Lemma 5.3.
###### Corollary 5.4.
For $`X_1,\mathrm{},X_n,A,B`$ and $`C`$ indeterminates,
$$\begin{array}{c}\underset{1i,jn}{det}\left(\frac{(AX_i,AC/X_i;q,p)_{nj}}{(BX_i,BC/X_i;q,p)_{nj}}\right)\hfill \\ \hfill =A^{\left(\genfrac{}{}{0pt}{}{n}{2}\right)}q^{\left(\genfrac{}{}{0pt}{}{n}{3}\right)}\underset{1i<jn}{}X_jE(X_i/X_j)E(C/X_iX_j)\underset{i=1}{\overset{n}{}}\frac{(B/A,ABCq^{2n2i};q,p)_{i1}}{(BX_i,BC/X_i;q,p)_{n1}}.\end{array}$$
We remark that this can be written as the following determinant identity for theta functions:
$$\begin{array}{c}\underset{1i,jn}{det}(T_{nj}(A+X_i)T_{nj}(A+CX_i)\hfill \\ \hfill \times T_{j1}(B+X_i+nj)T_{j1}(B+C+njX_i))\\ \hfill =\underset{1i<jn}{}\vartheta _1(X_iX_j)\vartheta _1(CX_iX_j)\underset{i=1}{\overset{n}{}}T_{i1}(BA)T_{i1}(A+B+C+2n2i),\end{array}$$
where $`T_n(x)=_{k=0}^{n1}\vartheta _1(x+k)`$ and $`\vartheta _1(x)`$ a standard theta function ,
$$\vartheta _1(x)=2\underset{k=0}{\overset{\mathrm{}}{}}(1)^np^{(2n+1)^2/4}\mathrm{sin}(2n+1)x=ip^{1/4}e^{ix}(p^2;p^2)_{\mathrm{}}E(e^{2ix};p^2).$$
For $`n=2`$ this is nothing but the well-known identity
$$\begin{array}{c}\vartheta _1(u+x)\vartheta _1(ux)\vartheta _1(v+y)\vartheta _1(vy)\vartheta _1(u+y)\vartheta _1(uy)\vartheta _1(v+x)\vartheta _1(vx)\hfill \\ \hfill =\vartheta _1(x+y)\vartheta _1(xy)\vartheta _1(u+v)\vartheta _1(uv),\end{array}$$
equivalent to (3.1).
###### Proof of Theorem 5.1.
By Corollary 5.4 with $`X_iq^{k_i}/x_i`$ and $`Ca`$, and $`E(x)=xE(1/x)`$ we can trade the double product
$$\underset{1i<jn}{}E(q^{k_ik_j}x_i/x_j)E(ax_ix_jq^{k_i+k_j})$$
for a determinant. If we also choose $`B=q^{2n}/c`$ and $`A=b/a`$, and use (2.7) and $`_{j=1}^n\left(\genfrac{}{}{0pt}{}{j1}{2}\right)=\left(\genfrac{}{}{0pt}{}{n}{3}\right)`$, the left-hand side of (5.1) can be rewritten as
(5.5)
$$\begin{array}{c}q^{3\left(\genfrac{}{}{0pt}{}{n}{3}\right)}\underset{1i<jn}{}\left(\frac{a^2x_j/bc^2}{E(x_j/x_i)E(ax_ix_jq^N)}\right)\underset{i=1}{\overset{n}{}}\frac{1}{(aq^{2n}/bc,bq^{n2i+2}/c;q,p)_{i1}}\hfill \\ \hfill \times \underset{1i,jn}{det}((cx_i,c/ax_i;q,p)_{j1}(bx_i,b/ax_i;q,p)_{nj}\\ \hfill \times {}_{8}{}^{}\omega _{7}^{}(ax_i^2;bx_iq^{nj},cx_iq^{j1},dx_i,ex_i,q^N;q,p)).\end{array}$$
Applying the elliptic $`{}_{8}{}^{}\omega _{7}^{}`$ sum of Theorem 2.11 and again using (2.7) as well as $`_{j=1}^n(nj)(n+j3)=4\left(\genfrac{}{}{0pt}{}{n}{3}\right)`$, we arrive at the following expression for the left-hand side of (5.1)
$$\begin{array}{c}q^{\left(\genfrac{}{}{0pt}{}{n}{3}\right)}\underset{1i<jn}{}\left(\frac{ax_jq^N/b}{E(x_j/x_i)E(ax_ix_jq^N)}\right)\underset{i=1}{\overset{n}{}}\frac{(q^{2n}/cx_i,ax_iq^{Nn+2}/c;q,p)_{n1}}{(bq^{n2i+2}/c,aq^{Nn+2}/bc;q,p)_{i1}}\hfill \\ \hfill \times \underset{i=1}{\overset{n}{}}\frac{(ax_i^2q,aq^{2i}/bc,aq^{2i}/bd,aq^{2i}/cd;q,p)_N}{(aqx_i/b,aqx_i/c,ax_iq/d,aq^{2n}/bcdx_i;q,p)_N}\\ \hfill \times \underset{1i,jn}{det}\left(\frac{(bx_i,bq^N/ax_i;q,p)_{nj}}{(q^{2n}/cx_i,ax_iq^{Nn+2}/c;q,p)_{nj}}\right).\end{array}$$
By Lemma 5.4 with $`X_i1/x_i`$, $`Abq^N/a`$, $`Bq^{2n}/c`$ and $`Caq^N`$ the first and third line are found to be reciprocal, thus resulting in the right-hand side of (5.1). ∎
## 6. Discussion
Of course the summations and transformations obtained in this paper for elliptic hypergeometric series are only a tip of the iceberg. Many more results for terminating, balanced, very-well-poised, basic hypergeometric series admit elliptic generalizations. In particular all the multivariable balanced, very-well-poised summation and transformation theorems of should admit elliptic counterparts. However, the methods of proof applied in these papers does not simply carry over the the elliptic case. In particular, the multivariable Jackson sums (from which most of the other results can be derived in ways well-tailored for elliptic generalization) are usually proved using simpler identities for series that are not both balanced and very-well poised. This is unlike the one-dimensional Jackson sum which can be proved simply by induction, without relying on other results – a method of proof that readily carries over to the elliptic case. Indeed the only higher-dimensional elliptic Jackson sum that we were able to prove so far is the one stated in Theorem 5.1.
One might expect that at least “the corresponding $`{}_{10}{}^{}\omega _{9}^{}`$ transformation” should be accessible with the techniques presented in this paper. However, all our attempts to find a C<sub>n</sub> $`{}_{10}{}^{}\omega _{9}^{}`$ transformation that implies Theorem 5.1 failed dismally. Surprisingly though, our failed attempts did suggest how to somewhat change Theorem 5.1 so that it does admit a generalization to a transformation. Since to the best of our knowledge this transformation (in the $`p0`$ limit) does not appear in the above list of references we state it here as a conjecture.
First we we need some more notation. Following Macdonald’s book we set
$$|\lambda |=\underset{i1}{}\lambda _i\text{and}n(\lambda )=\underset{i1}{}(i1)\lambda $$
for $`\lambda `$ a partition (i.e., $`\lambda =(\lambda _1,\lambda _2,\mathrm{})`$ with $`\lambda _i\lambda _{i+1}`$ and finitely many $`\lambda _i`$ nonzero). For a partition $`\lambda `$ of exactly $`n`$ parts (some of which may be zero) define
$$(a;q,p)_\lambda =\underset{i=1}{\overset{n}{}}(ax^{1j};q,p)_{\lambda _i}$$
and employ the usual condensed notation
$$(a_1,\mathrm{},a_m;q,p)_\lambda =(a_1;q,p)_\lambda \mathrm{}(a_m;q,p)_\lambda .$$
With these preliminaries we define a C<sub>n</sub> analogue of the balanced, very-well-poised, elliptic hypergeometric series (2.9) by
$$\begin{array}{c}{}_{r+1}{}^{}\mathrm{\Omega }_{r}^{}(a_1;a_4,\mathrm{},a_{r+1};q,p)\hfill \\ \hfill =\underset{\lambda _1\lambda _2\mathrm{}\lambda _n0}{}\underset{i=1}{\overset{n}{}}(\frac{E(a_1x^{2(1i)}q^{2\lambda _i})}{E(a_1x^{2(1i)})})\frac{(a_1x^{1n},a_4,\mathrm{},a_{r+1};q,p)_\lambda q^{|\lambda |}x^{2n(\lambda )}}{(qx^{n1},a_1q/a_4,\mathrm{},a_1q/a_{r+1};q,p)_\lambda }\\ \hfill \times \underset{1i<jn}{}(\frac{E(x^{ji}q^{\lambda _i\lambda _j})}{E(x^{ji})}\frac{E(a_1x^{2ij}q^{\lambda _i+\lambda _j})}{E(a_1x^{2ij})}\\ \hfill \times \frac{(a_1x^{3ij};q,p)_{\lambda _i+\lambda _j}(x^{ji+1};q,p)_{\lambda _i\lambda _j}}{(a_1qx^{1ij};q,p)_{\lambda _i+\lambda _j}(qx^{ji1};q,p)_{\lambda _i\lambda _j}}),\end{array}$$
where $`(a_4\mathrm{}a_{r+1})^2=a_1^{r3}q^{r5}x^{22n}`$. For reasons of convergence we again insist that one of the $`a_i`$ $`(i=4,\mathrm{},r+1)`$ is of the form $`q^N`$ with $`N`$ a nonnegative integer, so that the only nonvanishing contributions to the above sum come from $`\lambda _1N`$. Observe that for $`x=1`$ the double product in the summand simplifies to a multinomial coefficient, i.e., to $`_{1i<jn}(ji+1\delta _{\lambda _i,\lambda _j})/(ji)=n!/(m_0!m_1!\mathrm{}m_N!)`$, where $`m_k`$ is the number of parts of size $`k`$ in the partition $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$. Since $`lim_{x1}(a;q,p)_\lambda =_{i=1}^n(a;q,p)_{\lambda _i}`$ and
$$\underset{\lambda }{}\frac{n!}{m_0!\mathrm{}m_N!}\underset{i=1}{\overset{n}{}}a_{\lambda _i}=\underset{\begin{array}{c}0m_0,\mathrm{},m_Nn\\ m_1+\mathrm{}+m_N=n\end{array}}{}\frac{n!}{m_0!\mathrm{}m_N!}\underset{i=0}{\overset{N}{}}a_i^{m_i}=\left(\underset{i=0}{\overset{N}{}}a_i\right)^n,$$
where $`\lambda =(\lambda _1,\mathrm{},\lambda _n)=(0^{m_0}1^{m_1}\mathrm{}N^{m_N})`$), we may conclude that
$$\underset{x1}{lim}{}_{r+1}{}^{}\mathrm{\Omega }_{r}^{}(a_1;a_4,\mathrm{},a_{r+1};q,p)=\left({}_{r+1}{}^{}\omega _{r}^{}(a_1;a_4,\mathrm{},a_{r+1};q,p)\right)^n.$$
Computer assisted experiments suggest the following C<sub>n</sub> version of the $`{}_{10}{}^{}\omega _{9}^{}`$ transformation (2.10).
###### Conjecture 6.1.
Let $`bcdefgx^{n1}=a^3q^{N+2}`$ and $`\lambda =a^2q/bcd`$. Then
$$\begin{array}{c}{}_{10}{}^{}\mathrm{\Omega }_{9}^{}(a;b,c,d,e,f,g,q^N;q,p)\hfill \\ \hfill =\frac{(aq,aq/ef,\lambda q/e,\lambda q/f;q,p)_{(N^n)}}{(aq/e,aq/f,\lambda q/ef,\lambda q;q,p)_{(N^n)}}{}_{10}{}^{}\mathrm{\Omega }_{9}^{}(\lambda ;\lambda b/a,\lambda c/a,\lambda d/a,e,f,g,q^N;q,p).\end{array}$$
For $`cd=aq`$ this implies
###### Corollary 6.2.
For $`bcfgx^{n1}=a^2q^{N+1}`$ there holds
$${}_{8}{}^{}\mathrm{\Omega }_{7}^{}(a;b,c,d,e,q^N;q,p)=\frac{(aq,aq/bc,aq/bd,aq/cd;q,p)_{(N^n)}}{(aq/b,aq/c,aq/d,aq/bcd;q,p)_{(N^n)}}.$$
As remarked earlier we were unable to trace the $`p=0`$ case of the above two results in the literature, but we did find that letting $`d`$ tend to infinity after setting $`p=0`$, Corollary 6.2 reduces to a multivariable analogue of Rogers’ $`{}_{6}{}^{}\varphi _{5}^{}`$ sum due to van Diejen \[13, Thm. 3\].
Another challenging problem is to find nontrivial transformations based on the inverse pair given in (3.5) for all positive integers $`r`$. The only result for general $`r`$ obtained so far in the not-so-deep Theorem 4.1, which we were unable to generalize to a transformation. The problem with the type of transformations derived in section 4 appears to be that increasing $`r`$ has the effect of decreasing the number of available free parameters. For example, when we mimic the derivation of Theorems 4.2 and 4.3 but choose $`r=4`$ in (3.5) we no longer obtain a transformation for a $`{}_{10}{}^{}\omega _{9}^{}`$, but the less appealing quartic transformation
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(aq^{5k})}{E(a)}\frac{(b^2/aq^2;q,p)_k}{(a^2q^6/b^2;q^4,p)_k}\frac{(aq/b,aq^2/b,aq^3/b;q^2,p)_k}{(b,bq,bq^2;q^3,p)_k}\frac{(abq^{4n},q^{4n};q^4,p)_k}{(q^{14n}/b,aq^{4n+1};q,p)_k}q^k\hfill \\ \hfill =\frac{(aq;q,p)_{4n}(q^4,b^3/aq^2;q^4,p)_n}{(b;q,p)_{4n}(ab,a^2q^6/b^2;q^4,p)_n}\\ \hfill \times \underset{k=0}{\overset{n}{}}\frac{E(abq^{8k4})}{E(ab/q^4)}\frac{(ab/q^4,a^2q^2/b^2,b,b/q,b/q^2,b/q^3;q^4,p)_k}{(q^4,b^3/aq^2,a,aq,aq^2,aq^3;q^4,p)_k}q^{4k},\end{array}$$
which contains only two indeterminates. Moreover its counterpart (in the sense of $`(q^{4n};q^4,p)_k(q^n;q,p)_k)`$) no longer seems to allow for a transformation at all, admitting just
$$\begin{array}{c}\underset{k=0}{\overset{n}{}}\frac{E(a^2q^{5k})}{E(a^2)}\frac{(a^2;q^4,p)_k}{(q;q,p)_k}\frac{(a,aq,aq^2;q^3,p)_k}{(a,aq,aq^2;q^2,p)_k}\frac{(aq^{n+1},q^n;q,p)_k}{(aq^{3n},a^2q^{n+4};q^4,p)_k}q^k\hfill \\ \hfill =\{\begin{array}{cc}\frac{(q,q^2,q^3,a^2q^4;q^4,p)_{n/4}}{(aq^2,aq^3,aq^4,q/a;q^4,p)_{n/4}}\hfill & n0(mod4)\hfill \\ 0\hfill & n0(mod0),\hfill \end{array}\end{array}$$
which generalizes the quadratic and cubic summations of Corollaries 4.11 and 4.15.
### Acknowledgements
This work is supported by a fellowship of the Royal Netherlands Academy of Arts and Sciences. |
warning/0001/nucl-th0001016.html | ar5iv | text | # Strongly damped nuclear collisions: zero or first sound ?
## I Introduction
The phenomenon of strong dissipation in the collective motions of heated nuclear systems is a challenging problem for the nuclear transport theory. Following the analogy with other Fermi liquids, like the liquid <sup>3</sup>He , one expects also in nuclear systems two types of collective modes.
At small temperatures the mean field dominated modes should exist, from the possible violatation of the local thermal equilibrium due to the strong Pauli blocking which inhibits two-body collisions (the Landau zero sounds). This is actually the main nature of isoscalar giant resonances, where indeed the collective energy is essentially given by the amount needed to deform the Fermi sphere in momentum space . Isovector giant resonances are also mean field modes, this time of plasmon type since we can have both a distortion and a shift of the neutron (proton) Fermi spheres.
At high temperatures the usual hydrodynamical collective modes (first sounds) should propagate . In an infinite Fermi liquid with strong repulsion (Landau parameter $`F_01`$) the two sounds are clearly distinguishable, since the damping rate has a maximum as a function of the temperature . However, in nuclear matter the transition between the two sounds is expected to be smeared-out since the Landau parameter $`F_0`$ is small in absolute value .
The presence of this transition is still an open problem in nuclear dynamics. Only the isovector giant dipole resonances can be experimentally studied with sufficient accuracy in heated nuclei. The point is that for this kind of two-component mode the transition has some special features that make it difficult to observe a clear signature .
Aim of this work is to study the transition in the collective nuclear dynamics looking at the evolution of the damping mechanism of large amplitude quadrupole oscillations in fusion processes, in a microscopic kinetic model. Similar attempts have been recently performed in fission dynamics studies . We will also follow the temperature dependence of the attenuation of the collective mode.
In the ref.s a reduced friction coefficient
$$\beta =\frac{1}{E_{kin}^{coll}}\left(\frac{dE}{dt}\right)_{diss},$$
(1)
where $`E_{kin}^{coll}`$ is the collective kinetic energy of the dinuclear system (DNS) and $`(dE/dt)_{diss}`$ is the dissipation rate of the total collective energy $`E=E_{kin}^{coll}+E_{pot}`$, has been extracted from the measurements of the prescission neutron multiplicity in fast fission reactions. It was shown that in the temperature region $`T=2÷3`$ MeV the reduced friction coefficient $`\beta `$ is very large ($`\beta 10÷10010^{21}`$ s<sup>-1</sup>), result not explained with a one-body dissipation mechanism only.
In the present work we have performed some kinetic transport studies of the nuclear dynamics. Transport equations describe a self-consistent mean field dynamics coupled to two-body collisions and so we expect to see in a natural way the transition between the two sound propagations. We use the Boltzmann-Nordheim-Vlasov (BNV) procedure to simulate the phase space dynamics, which has been quite successful in predicting mean properties of heavy ion collisions at medium energies. In particular we study central nucleus-nucleus collisions leading to fusion at beam energies $`E_{lab}<21`$ AMeV in order to extract the coefficient $`\beta `$ as a function of the temperature from the damping of the collective quadrupole oscillations of the formed di-nuclear system (DNS). Density dependent effective interactions of Skyrme type ($`SKM`$ ) and an averaged free nucleon-nucleon cross section, $`\sigma =40mb`$ () have been used.
The structure of the work is as follows. In Sect. 2 the procedure of the extraction of the reduced damping coefficient $`\beta `$ from the time dependence of the quadrupole moment given by the BNV model is described. In Sect. 3 we present a novel method for determination of the temperature based on the energy conservation and on the BNV evolution of the potential energy. Sect. 4 contains our results on the $`\beta (T)`$ and their interpretation in terms of two-body and one-body dissipation mechanisms. Summary and conclusions are given in Sect. 5.
## II Extraction of the reduced friction coefficient from nuclear kinetic equations
According to Eq. (1), the reduced friction coefficient $`\beta `$ can be, in principle, calculated directly, once the phase-space distribution function $`f(𝐫,𝐩,t)`$ is known from the BNV output, since:
$`E_{kin}^{coll}(t)`$ $`=`$ $`{\displaystyle 𝑑𝐫\frac{mv(𝐫,t)^2}{2}\rho (𝐫,t)},`$ (2)
$`𝐯(𝐫,t)`$ $`=`$ $`{\displaystyle \frac{1}{\rho }}{\displaystyle 𝑑𝐩f(𝐫,𝐩,t)\frac{𝐩}{m}},`$ (3)
$`\rho (𝐫,t)`$ $`=`$ $`{\displaystyle 𝑑𝐩f(𝐫,𝐩,t)},`$ (4)
$`E_{pot}(t)`$ $`=`$ $`E_{pot}^{int}(t)+{\displaystyle \frac{3}{5}}{\displaystyle 𝑑𝐫ϵ_F(\rho )\rho },`$ (5)
$`E_{pot}^{int}(t)`$ $`=`$ $`{\displaystyle 𝑑𝐫ϵ_{m.f.}(\rho )}+E_{coul}(t),`$ (6)
where $`ϵ_F(\rho )=\mathrm{}^2/(2m)(3\pi ^2\rho /2)^{2/3}`$ is the Fermi energy, $`ϵ_{m.f.}(\rho )`$ is the nuclear mean field interaction energy density and $`E_{coul}`$ is the Coulomb energy. From Eq.(2) we see that the beam energy is not giving contribution to the collective kinetic energy.
However, in practice, the calculation of the collective kinetic energy $`E_{kin}`$ is quite ambiguous in the test particle technique due to the strong dependence on the width of the gaussians representing the test particles .
We calculate then the coefficient $`\beta `$ from the time evolution of the quadrupole moment of the DNS:
$$Q_{zz}(t)=𝑑𝐫q_{zz}(𝐫)\rho (𝐫,t),q_{zz}(𝐫)=2z^2x^2y^2.$$
(7)
For a damped periodical motion (see Appendix)
$$Q_{zz}(t)\mathrm{exp}(i\omega t),\omega =\omega _R+i\omega _I,\omega _I<0,$$
(8)
the coefficient $`\beta `$ is proportional to the imaginary part of the frequency $`\omega `$:
$$\beta =4\omega _I.$$
(9)
The curves in Fig. 1 show the time evolution of $`Q_{zz}`$ in the case of central collisions of <sup>64</sup>Ni + <sup>238</sup>U at beam energies $`E_{lab}=6.53÷20.53`$ AMeV. The quadrupole moment quickly approaches a minimum at $`t=100÷200`$ fm/c. Afterwards the $`Q_{zz}`$ starts to grow again, and after a time interval $`\mathrm{\Delta }t=50÷100`$ fm/c it saturates. This saturation of the $`Q_{zz}`$ is caused by an almost complete lost of the collective energy of the DNS. One can, therefore, qualitatively estimate the $`\beta `$ coefficient just assuming that $`|\omega _I|\omega _R=2\pi /t_{osc}`$, where $`t_{osc}100`$ fm/c is the period of oscillations. That gives $`\beta =4|\omega _I|7010^{21}`$ s. This value is in agreement with the results obtained in Refs. . However, we would like to stress that in Refs. the damping coefficient was extracted from the outgoing part of the trajectory of the DNS from compact mononucleus shape to scission during a time interval $`30000`$ fm/c. Our analysis is concentrated on the part of the trajectory in vicinity of the mononucleus shape and the corresponding time scale $`300`$ fm/c is much shorter.
In order to get the values of the coefficient $`\beta `$ we fit a part of the curve $`Q_{zz}(t)`$ to the function
$$Q_{zz}^{fit}(t)=𝒜+\mathrm{sin}(\omega _Rt+\varphi _0)\mathrm{exp}(|\omega _I|t).$$
(10)
The upper time limit of the fitting region is chosen at the second minimum of the quadrupole moment. The lower time limit is given by the earliest time when the same value of the $`Q_{zz}`$ as in the second minimum is reached. This definition of the time limits, from one hand, corresponds approximately to the selection of a full oscillation in vicinity of the compact mononucleus shape. On the other hand, the stage of strong dissipation from the first minimum to the first maximum is completely in the fitting region. The best fit functions of Eq. (10) are shown by full dots in Fig.1. The fitting parameters $`𝒜,,\omega _R,\varphi _0,\omega _I`$ and the corresponding coefficient $`\beta =4|\omega _I|`$ for central collisions of <sup>64</sup>Ni + <sup>238</sup>U at various beam energies are collected in the Table. We see that the damping increases with the increasing collision energy until $`E_{lab}=10.53`$ AMeV. This fact can be understood already from Fig. 1, since with increasing beam energy a larger $`\mathrm{\Delta }Q_{zz}=Q_{zz}^{saturation}Q_{zz}^{minimum}`$ is damped during a shorter time interval. But at higher energies $`E_{lab}>10.53`$ AMeV the damping starts to decrease, as we can see from the presence of quadrupole vibrations at later times $`t>200`$ fm/c (Fig. 1). Eventually at $`E_{lab}>14.53`$ AMeV a kind of saturation seems to be reached.
## III Determination of the temperature from BNV simulations
A direct way to extract the temperature is to fit the local momentum distribution given by the BNV model to a $`T0`$ Fermi distribution. However, in the case of low-energy nuclear collsions studied in present work, this direct method is not appropriate just because its accuracy of $`1`$ MeV is not enough. Therefore we will follow a procedure based on the conservation of the energy and on the time evolution of the potential energy given by the BNV. The thermal excitation energy of the DNS is (c.f. ):
$$E_{therm}^{}=E_{kin}^{c.m.}E_{coul}E_{rot}\mathrm{\Delta }E_{pot},$$
(11)
where $`E_{kin}^{c.m.}=E_{lab}A_1A_2/(A_1+A_2)`$ is the center-of-mass kinetic energy, $`E_{coul}=e^2Z_1Z_2/(R_1+R_2+3.5)`$ is the Coulomb energy, $`R_i=1.2A_i^{1/3}`$ (fm) $`i=1,2`$ are the nuclear radii, $`E_{rot}=\mathrm{}^2L^2/(2\mathrm{\Theta })`$ is the rotational energy with $`L`$ being the angular momentum (in $`\mathrm{}`$ units) and $`\mathrm{\Theta }=[\frac{2}{5}(A_1R_1^2+A_2R_2^2)+A_1A_2/(A_1+A_2)(R_1+R_2)^2]m_{nuc}`$ being the momentum of inertia ($`m_{nuc}`$ is the nucleon mass). The last term $`\mathrm{\Delta }E_{pot}`$ in the r.h.s. of Eq.(11) is the difference between the values of the potential energy just before and after the overlapping of the density profiles. The origin of this term is mostly from the sharp decrease of the nuclear surface energy when the two nuclei touch each other. In Fig. 2 we show the time dependence of the potential energy per nucleon in the central collision of <sup>64</sup>Ni + <sup>238</sup>U at 10.53 AMeV. After some small increase until $`t20`$ fm/c due to the Coulomb contribution the total potential energy quickly drops by about 1 MeV/nucleon reaching the minimum at $`t60`$ fm/c. We define $`\mathrm{\Delta }E_{pot}`$ as
$$\mathrm{\Delta }E_{pot}=E_{pot}^{min}E_{pot}^{max},$$
(12)
where $`E_{pot}^{min}`$ and $`E_{pot}^{max}`$ are the first minimum and the first maximum of the potential energy.
In the Table we report the values of the potential energy "jumps" $`\mathrm{\Delta }E_{pot}`$, the excitation energies and corresponding temperatures $`T=\sqrt{E_{therm}^{}/a},a=A/8\text{MeV}^1`$ for the <sup>64</sup>Ni + <sup>238</sup>U central collisions at various beam energies<sup>*</sup><sup>*</sup>* In the case of the central collisions $`E_{rot}=0`$ in Eq. (11) . Our temperatures are higher than those obtained in Ref. . In particular, for the Ni+U collision at 6.53 AMeV we have $`T=3.1`$ MeV and the authors of Ref. report $`T=2.4`$ MeV. This difference is mainly explained by the fact, that the temperatures in Ref. are obtained from the outgoing stage of the reaction on much larger time scales as compared to our study. A further contribute to some temperature overshooting is coming from pre-equilibrium emissions, neglected in the energy balance Eq.(11). This effect, although expected to increase with beam energy, can be still considered quite small in this energy range.
## IV Temperature dependence of the dissipation
Fig. 3 shows the temperature dependence of the reduced friction coefficient $`\beta `$. The BNV results are shown by solid line with dots. The coefficient $`\beta `$ first increases with temperature reaching the maximum at $`T4.5`$ MeV and then at higher temperatures it decreases. This signal is indeed quite robust and cannot be related to some overestimation of the temperatures at higher beam energies as discussed at the end of the previous section.
One can interpret the results of the full transport calculations within the approach of Ref. , based on the analytical solution of the linearized Landau-Vlasov equation
$$\left(\frac{}{t}+\frac{𝐩}{m^{}}\frac{}{𝐫}\right)\delta f(𝐫,𝐩,t)\frac{\delta U}{𝐫}\frac{f_{\mathrm{eq}}(𝐩)}{𝐩}=I_{coll}[\delta f],$$
(13)
where $`\delta f`$ and $`f_0`$ are the perturbation and the equilibrium value of the phase space distribution function,
$$\delta U(𝐫,𝐩,t)=\frac{1}{N(0)}\frac{gd𝐩^{}}{(2\pi \mathrm{})^3}(F_0+F_1\widehat{p}\widehat{p}^{})\delta f(𝐫,𝐩^{},t)$$
(14)
is the mean field perturbation ($`\widehat{p}𝐩/p`$, $`\widehat{p}^{}𝐩^{}/p^{}`$), $`g=4`$ is the spin-isospin degeneracy of a nucleon, $`N(0)=gm^{}p_F/(2\pi ^2\mathrm{}^3)`$ is the level density at zero temperature, $`F_0`$ and $`F_1`$ are the Landau parameters and $`m^{}=m/(1+F_1/3)`$ is the effective mass.
The collision integral in the r.h.s. of Eq. (13) is taken in the relaxation time approximation:
$$I_{coll}[\delta f]\frac{1}{\tau }\delta f_{|l2},$$
(15)
where
$$\delta f_{|l2}(𝐩)\underset{l2}{}\underset{m=l}{\overset{l}{}}Y_{lm}(\widehat{p})𝑑\mathrm{\Omega }_{\widehat{p}^{}}Y_{lm}^{}(\widehat{p}^{})\delta f(𝐩^{})_{|p^{}=p}$$
(16)
is the part of the perturbation containing the quadrupole and higher multipolarity distortions of the Fermi surface. The effective relaxation time $`\tau `$ includes two- and one-body dissipation contributions:
$$\tau ^1=\tau _{2body}^1+\tau _{1body}^1.$$
(17)
The relaxation time $`\tau _{2body}`$ was calculated in Ref. for various choices of a nucleon-nucleon scattering cross section:
$$\tau _{2body}^1=T^2/\kappa $$
(18)
with $`\kappa 1900`$ MeV<sup>2</sup>fm/c for the isotropic energy independent isospin-averaged nucleon-nucleon scattering cross section $`\sigma _{NN}=40`$ mb. For the one-body relaxation time we have used the wall-and-window formula (c.f. Ref. ):
$$\tau _{1body}^1=\frac{\overline{v}}{2R\xi },$$
(19)
where $`R=1.2A^{1/3}`$ ($`A=A_1+A_2`$) is the radius of a mononucleus composed of the two colliding nuclei,
$$\overline{v}=\frac{3}{4}v_F\left[1+\frac{\pi ^2}{6}\left(\frac{T}{ϵ_F}\right)^2\right]$$
(20)
is an average velocity of nucleons, and $`\xi `$ is a numerical factor which depends on the multipolarity and on the isospin of a collective mode. We have chosen a value $`\xi =1.85`$, which corresponds to the isoscalar quadrupole mode in the scaling wall model (see Ref. and refs. therein).
The solution of Eq. (13) inside a nucleus with uniform nonperturbed density can be found as a superposition of plane waves. In this case, as it was shown in Ref. , the intrinsic width of a giant multipole resonance $`\mathrm{\Gamma }=2\mathrm{}/\tau _{rel}`$, where $`\tau _{rel}`$ is the relaxation time of the distribution function $`\delta f`$ ($`\delta f\mathrm{exp}(i\omega t)`$, $`\omega =\omega _R+i\omega _I`$, $`\omega _I=1/\tau _{rel}`$), can be expressed as
$$\mathrm{\Gamma }2q\mathrm{}\omega _R\frac{\omega _R\tau }{1+q(\omega _R\tau )^2}$$
(21)
with $`q=[\frac{5}{2}(1+F_0)(1+F_1/3)]^1`$. Eq. (21) describes both well known zero sound ($`\mathrm{\Gamma }\tau ^1`$, $`\omega _R\tau 1`$) and first sound ($`\mathrm{\Gamma }\tau `$, $`\omega _R\tau 1`$) regimes. In our calculations we have put the Landau parameters $`F_0=0.2`$, $`F_1=0`$. The frequency has been chosen as $`\mathrm{}\omega _R=64.7A^{1/3}`$ MeV ($`A=A_1+A_2`$) corresponding to the giant quadrupole resonance.
The solid line in Fig. 3 shows the friction coefficient
$$\beta _{ZFST}=2\mathrm{\Gamma }/\mathrm{}$$
(22)
with $`\mathrm{\Gamma }`$ given by Eq. (21). We see that Eqs. (21), (22) agree qualitatively with the BNV model, but the analytical calculation gives a more smeared-out transition between the two sounds.
In the limit of the zero sound the coefficient $`\beta `$ is
$$\beta _{ZS}=4/\tau =\beta _{ZS}^{2body}+\beta _{ZS}^{1body},$$
(23)
where $`\beta _{ZS}^{2body}=4/\tau _{2body}`$, $`\beta _{ZS}^{1body}=4/\tau _{1body}`$. This simple formula (dot-dashed line in Fig. 3) is quite close to the BNV results at $`T4.5`$ MeV. However, either the two-body (dashed line) or one-body (dotted line) contributions taken separately are strongly underpredicting the BNV friction coefficient.
In order to better understand the relative importance of the two-body and one-body mechanisms we have also performed the BNV calculations at $`E_{lab}=6.53,8.53,10.53,12.53`$ and $`14.53`$ AMeV switching off the collision term, corresponding to a pure Vlasov evolution. Fig. 4 shows the comparison of the $`Q_{zz}`$ time evolution with and without collision term for the Ni+U collision at 14.53 AMeV. In the case without collisions the damping is much reduced: we observe several large-amplitude oscillations of the quadrupole moment. We have also checked that in the Vlasov evolution the oscillations at later times, $`t>200`$ fm/c, are present for all the other studied beam energies. In Fig. 3, the coefficient $`\beta `$ in the case of collisionless dynamics is shown by full squares. The damping slowly increases with temperature and saturates at $`T=5`$ MeV. As expected, the system is always in the region of the zero sound. The absolute value of the reduced friction coefficient in the collisionless case is $`3÷5`$ times less than for a full BNV calculation and close to the wall-and-window value $`4/\tau _{1body}`$.
## V Summary and conclusions
The reduced friction coefficient was extracted for the initial stage ($`300`$ fm/c) of central <sup>64</sup>Ni + <sup>238</sup>U collisions at beam energies $`E_{lab}=6÷21`$ AMeV in the framework of the BNV transport model. The quadrupole moment $`Q_{zz}`$ has been chosen as a relevant collective variable. The quadrupole time evolution shows overdamped oscillations with a damping rate proportional to the friction coefficient. Two-body collisions play a major role in the damping of $`Q_{zz}`$.
As a function of the beam energy, the damping rate has a maximum at $`E_{lab}10`$ AMeV. The corresponding temperature of the DNS, neglecting particle emissions, is 4.6 MeV. We interpret, therefore, this temperature as a transition temperature from the zero-to-first sound propagation.
This result agrees with earlier calculations of V.M. Kolomietz et al. (Ref. ), where a value of the transition temperature $`4÷5`$ MeV was deduced on the basis of the analytical solution of the linearized BUU equation in the case of isoscalar giant resonances of multipolarities $`l=0`$ and 2 in a hot nucleus.
Our calculations are overpredicting the transition temperature $`2÷2.5`$ MeV that follows from the analysis of the prescission neutron multiplicities by J. Wilczy$`\stackrel{´}{\mathrm{n}}`$ski et al. (Ref. ). We have to stress that the considered collective modes are different, in our analysis quadrupole oscillations in the entrance channel dynamics while in the ref. the fission mode in the exit channel. In particular a relevant variance is on the time scales of the two modes, with a much larger proper time for the fission dynamics. Following the simple condition $`\omega _R\tau 1`$ for the transition, we roughly get a $`T_{tr}\sqrt{\omega _R}`$ and then we can expect quite smaller transition temperatures for the fission modes. Moreover, as already stressed before, likely our temperature assignements are a little overestimated, particularly for the higher beam energies, due to the lack of pre-equilibrium emissions in the energy balance.
Looking at the BNV simulations, the transition from zero-to-first sound appears as a presence of quadrupole vibrations at relatively late times, $`t>200`$ fm/c (see Fig. 1), for beam energies above 10 AMeV. These vibrations would accelerate the fission of the produced DNS that should be observed experimentally, as an increase of fast-fission cross sections.
## Acknowledgements
Authors are grateful to Dr. G.G. Adamian for many fruitful discussions and useful advices. Stimulating discussions with Prof. D.M. Brink are gratefully acknowledged.
## Appendix
In this Appendix we derive the relation (9) between the imaginary part $`\omega _I`$ of the complex frequency $`\omega =\omega _R+i\omega _I`$ and the reduced damping coefficient $`\beta `$ of Eq. (1).
From the Eq. (7) and the continuity equation $`\dot{\rho }=(𝐯\rho )`$ we get:
$$\dot{Q}_{zz}=\rho 𝐯q_{zz}d^3r.$$
(24)
In the case of the small-amplitude damped periodical motion of the kind
$$𝐯(𝐫,t)=\delta 𝐯_0(𝐫)\mathrm{exp}(i\omega t)$$
(25)
the Eq. (24) can be linearized with respect to small values of $`\delta v_0`$:
$$\dot{Q}_{zz}\left[\rho _0\delta 𝐯_0(𝐫)q_{zz}d^3r\right]\mathrm{exp}(i\omega t)\dot{Q}_{zz}(t=0)\mathrm{exp}(i\omega t),$$
(26)
where $`\rho _0`$ is the nonperturbed density ($`\rho =\rho _0+\delta \rho `$). The general solution of Eq. (26) is
$$Q_{zz}(t)=\text{const}+\frac{i\dot{Q}_{zz}(0)}{\omega }\mathrm{exp}(i\omega t).$$
(27)
Therefore in a linear approximation, over $`\delta v_0`$, the quadrupole moment reveals oscillations with the same frequency $`\omega =\omega _R+i\omega _I`$ as the velocity field $`𝐯(𝐫,t)`$.
The time derivative of the collective kinetic energy (2) is:
$`\dot{E}_{kin}`$ $``$ $`m\rho _0{\displaystyle \mathrm{Re}(𝐯)\mathrm{Re}(\dot{𝐯})d^3r}=`$ (29)
$`m\rho _0\left[{\displaystyle \delta v_0^2(𝐫)d^3r}\right]\mathrm{cos}(\omega _Rt)(\omega _I\mathrm{cos}(\omega _Rt)\omega _R\mathrm{sin}(\omega _Rt))\mathrm{exp}(2\omega _It).`$
Averaging Eq. (29) over the period of oscillations we come to the relation:
$$\overline{\dot{E}_{kin}}\omega _I\rho _0m\overline{(\mathrm{Re}(𝐯))^2}d^3r=2\omega _I\overline{E_{kin}},$$
(30)
where we have dropped the therm $`\mathrm{cos}(\omega _Rt)\mathrm{sin}(\omega _Rt)\mathrm{exp}(2\omega _It)`$ which changes the sign during the period of oscillation. According to the virial theorem for harmonic oscillators
$$\overline{E_{kin}}=\overline{E_{pot}}=\frac{1}{2}\overline{E},$$
(31)
where $`E`$ is the total collective energy. Therefore
$$\overline{\dot{E}}4\omega _I\overline{E_{kin}}.$$
(32)
Fit parameters from Eq. (10); potential energy "jumps" $`\mathrm{\Delta }E_{pot}`$; excitation energies and temperatures of the DNS Eqs. (11),(12); for central collisions <sup>64</sup>Ni + <sup>238</sup>U at various beam energies $`E_{lab}`$.
| $`E_{lab}`$ | $`𝒜`$ | $``$ | $`\omega _R`$ | $`\varphi _0`$ | $`\omega _I`$ | -$`\mathrm{\Delta }E_{pot}`$ | $`E^{}`$ | T |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| (AMeV) | (Afm<sup>2</sup>) | (Afm<sup>2</sup>) | c/fm | | c/fm | AMeV | AMeV | MeV |
| 6.53 | 33.81 | 24.34 | 0.024 | -6.537 | 0.0133 | 0.88 | 1.2 | 3.1 |
| 8.53 | 25.60 | 39.94 | 0.030 | -6.169 | 0.0175 | 1.28 | 1.9 | 3.9 |
| 10.53 | 21.46 | 38.45 | 0.036 | -6.197 | 0.0180 | 1.68 | 2.7 | 4.6 |
| 12.53 | 18.34 | 31.13 | 0.038 | -6.091 | 0.0163 | 2.10 | 3.4 | 5.2 |
| 14.53 | 16.42 | 27.03 | 0.039 | -5.906 | 0.0147 | 2.28 | 3.9 | 5.6 |
| 16.53 | 14.33 | 23.50 | 0.040 | -5.843 | 0.0137 | 2.46 | 4.4 | 6.0 |
| 18.53 | 12.89 | 26.73 | 0.039 | -5.589 | 0.0141 | 2.68 | 5.0 | 6.3 |
| 20.53 | 11.14 | 28.05 | 0.038 | -5.486 | 0.0147 | 2.74 | 5.4 | 6.6 |
## Figure captions
Solid lines – time dependence of the quadrupole moment $`Q_{zz}`$ for <sup>64</sup>Ni + <sup>238</sup>U central collisions at beam energies (from top to bottom) $`E_{lab}=`$ 6.53, 8.53, 10.53, 12.53, 14.53, 16.53 and 20.53 AMeV. The value of the quadrupole moment at t=0 is always 85 fm<sup>2</sup>/nucleon. Full circles show the best fits obtained using Eq. (10) for each collision energy. See text after Eq. (10) for a definition of the fit regions.
Potential interaction energy per nucleon Eq. (6) as a function of time for the <sup>64</sup>Ni(10.53 AMeV) + <sup>238</sup>U central collision.
Temperature dependence of the reduced friction coefficient $`\beta `$. BNV calculations with (without) collision term are shown by full circles (squares) connected with thin solid line to guide eye. Errorbars are due to the ambiguity caused by a finite number of test particles. The analytical result of Ref. for the Giant Quadrupole Resonance in a hot nucleus (see Eqs. (21), (22)) is shown by thick solid line. The two- and one-body contributions and their sum (see Eq. (23)) in the zero sound limit are shown by dashed, dotted and dash-dotted lines respectively.
Comparison of the time evolution of the quadrupole moment for <sup>64</sup>Ni(14.53 AMeV) + <sup>238</sup>U central collision calculated with (solid line) and without (dashed line) collision term. |
warning/0001/hep-th0001150.html | ar5iv | text | # Introduction
## Introduction
Recently a way to obtain a supersymmetric action functional for interacting branes (intersecting branes and branes ending on branes) has been proposed . The systems involving open fundamental superstrings ending on super–Dp–branes are quite generic and, on the other hand, especially interesting. The case of superstring—super-D3-brane system has been discussed briefly in (see for details).
Two types of such system are special and require separate consideration. One consists of the open superstring and a super–D9–brane (space–time filling brane). It has been elaborated in . In this contribution we present a supersymmetric action functional for the system of the open superstring ending on (the dynamical) super-D0-branes or D-particles. In distinction to the general case neither Lagrange multipliers no auxiliary space–time filling brane are necessary in this case.
Note that this dynamical system provides a supersymmetric generalization of the ’string with masses at the endpoints’ which has been considered in the early years of ’QCD string’ .
## 1 Geometric action for free super-D0-brane
The geometric action and the generalized action principle for super–Dp–branes with $`0<p<9`$ and $`p=9`$ has been constructed in respectively. However the super–D0–brane has not been considered in this framework.
The geometric action for the super–D0–brane has the form
$$S_{D0}=m_^1\stackrel{~}{}_1=_^1\left(\stackrel{~}{\mathrm{\Pi }}^{\underset{¯}{m}}u_{\underset{¯}{m}}^{(0)}(\tau )+i\left(d\stackrel{~}{\mathrm{\Theta }}^{1\underset{¯}{\mu }}\stackrel{~}{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2\stackrel{~}{\mathrm{\Theta }}^{1\underset{¯}{\mu }}d\stackrel{~}{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2\right)\right),$$
(1)
where $`m`$ is the super-D0-brane mass parameter,
$$\mathrm{\Pi }^{\underset{¯}{m}}=dX^{\underset{¯}{m}}id\mathrm{\Theta }^{1\underset{¯}{\mu }}\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}\mathrm{\Theta }^{1\underset{¯}{\nu }}id\mathrm{\Theta }_{\underset{¯}{\mu }}^2\stackrel{~}{\sigma }^{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}\mathrm{\Theta }_{\underset{¯}{\nu }}^2$$
(2)
is the basic covariant 1–form of the flat type IIA superspace,
$$\stackrel{~}{\mathrm{\Pi }}^{\underset{¯}{m}}=d\stackrel{~}{X}^{\underset{¯}{m}}id\stackrel{~}{\mathrm{\Theta }}^1\sigma ^{\underset{¯}{m}}\stackrel{~}{\mathrm{\Theta }}^1id\stackrel{~}{\mathrm{\Theta }}^2\stackrel{~}{\sigma }^{\underset{¯}{m}}\stackrel{~}{\mathrm{\Theta }}^2=d\tau \stackrel{~}{\mathrm{\Pi }}_\tau ^{\underset{¯}{m}}$$
(3)
is its pull–back on the super–D0–brane world–line $`^1`$
$$X^{\underset{¯}{m}}=\stackrel{~}{X}^{\underset{¯}{m}}(\tau ),\mathrm{\Theta }^{1\underset{¯}{\mu }}=\stackrel{~}{\mathrm{\Theta }}^{1\underset{¯}{\mu }}(\tau ),\mathrm{\Theta }_{\underset{¯}{\mu }}^2=\stackrel{~}{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2(\tau ):^1\underset{¯}{}^{(10|32)},$$
(4)
and $`u_{\underset{¯}{m}}^{(0)}`$ is a time–like unit length vector field
$$u_{\underset{¯}{m}}^{(0)}u^{(0)\underset{¯}{m}}=1.$$
(5)
It is convenient to consider $`u_{\underset{¯}{m}}^{(0)}`$ as a column of the Lorentz group valued matrix
$$u_{\underset{¯}{m}}^{\underset{¯}{a}}=(u_{\underset{¯}{m}}^0,u_{\underset{¯}{m}}^i)SO(1,9)u_{\underset{¯}{m}}^{\underset{¯}{a}}\eta ^{\underset{¯}{m}\underset{¯}{n}}u_{\underset{¯}{n}}^{\underset{¯}{b}}=\eta ^{\underset{¯}{a}\underset{¯}{b}}.$$
(6)
The conditions (6) include the normalization (5) as well as the orthogonality conditions for the vectors $`u_{\underset{¯}{m}}^{(0)},u_{\underset{¯}{m}}^i`$ (Lorentz harmonics )
$$u_{\underset{¯}{m}}^{(0)}u^{i\underset{¯}{m}}=0,u_{\underset{¯}{m}}^iu^{j\underset{¯}{m}}=\delta ^{ij}.$$
(7)
A doubly covered element for the $`SO(1,9)`$–valued matrix (6)
$$v_{\underset{¯}{\mu }}^ASpin(1,9)$$
(8)
(spinor Lorentz harmonics, see and refs. therein) is related with (6) by the conditions of $`\sigma `$–matrix conservation
$$u_{\underset{¯}{m}}^{\underset{¯}{a}}\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}=v_{\underset{¯}{\mu }}^A\sigma _{AB}^{\underset{¯}{a}}v_{\underset{¯}{\mu }}^B,u_{\underset{¯}{m}}^{\underset{¯}{a}}\stackrel{~}{\sigma }_{\underset{¯}{a}}^{AB}=v_{\underset{¯}{\mu }}^A\stackrel{~}{\sigma }_{\underset{¯}{m}}^{\underset{¯}{\mu }\underset{¯}{\nu }}v_{\underset{¯}{\mu }}^B.$$
(9)
$`A=1,\mathrm{},16`$ can be treated as $`SO(9)`$ spinor index. Then the requirement of $`SO(9)`$ gauge symmetry makes natural an identification of the harmonics with homogeneous coordinates of the coset $`SO(1,9)/SO(9)`$ (cf. with ). Note that $`SO(9)`$ group possesses a symmetric charge conjugation matrix. When it is identified with unity matrix, the difference between upper and lower $`SO(9)`$ spinor indices disappears.
Substituting the $`SO(9)`$ invariant representation for $`SO(1,9)`$ sigma-matrices
$$\sigma _{AB}^0=\delta _{AB},\sigma _{AB}^i=\mathrm{\Gamma }_{AB}^i,\stackrel{~}{\sigma }^{0AB}=\delta _{AB},\stackrel{~}{\sigma }^{iAB}=\mathrm{\Gamma }_{AB}^i,$$
(10)
one can decompose Eq. (9) into
$$u_{\underset{¯}{m}}^{(0)}\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}=v_{\underset{¯}{\mu }}^Av_{\underset{¯}{\nu }}^A,u_{\underset{¯}{m}}^i\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}=v_{\underset{¯}{\mu }}^A\mathrm{\Gamma }_{AB}^iv_{\underset{¯}{\nu }}^B.$$
(11)
$$u_{\underset{¯}{m}}^{(0)}\delta _{AB}=v_{\underset{¯}{\mu }}^A\stackrel{~}{\sigma }_{\underset{¯}{m}}^{\underset{¯}{\mu }\underset{¯}{\nu }}v_{\underset{¯}{\nu }}^B,u_{\underset{¯}{m}}^i\mathrm{\Gamma }_{AB}^i=v_{\underset{¯}{\mu }}^A\stackrel{~}{\sigma }_{\underset{¯}{m}}^{\underset{¯}{\mu }\underset{¯}{\nu }}v_{\underset{¯}{\mu }}^B.$$
(12)
Similar relations can be obtained for the inverse $`SO(1,9)/SO(9)`$ harmonics
$$v_B^{\underset{¯}{\mu }}v_{\underset{¯}{\mu }}^A=\delta _B^A,$$
(13)
$$u_{\underset{¯}{m}}^{(0)}\stackrel{~}{\sigma }^{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}=v_A^{\underset{¯}{\mu }}v_A^{\underset{¯}{\nu }},u_{\underset{¯}{m}}^i\stackrel{~}{\sigma }^{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}=v_A^{\underset{¯}{\mu }}\mathrm{\Gamma }_{AB}^iv_A^{\underset{¯}{\nu }}.$$
(14)
$$u_{\underset{¯}{m}}^{(0)}\delta _{AB}=v_A^{\underset{¯}{\mu }}\sigma _{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}v_B^{\underset{¯}{\nu }},u_{\underset{¯}{m}}^i\mathrm{\Gamma }_{AB}^i=v_A^{\underset{¯}{\mu }}\sigma _{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}v_B^{\underset{¯}{\mu }}.$$
(15)
The harmonics can be used to define a general supervielbein of the flat type $`IIA`$ superspace $`E^𝒜`$ which possesses the $`SO(9)`$ invariant decomposition $`E^𝒜=(E^{(0)},E^i;E^{A1},E^{A2})`$
$$E^{(0)}\mathrm{\Pi }^{\underset{¯}{m}}u_{\underset{¯}{m}}^{(0)},E^i\mathrm{\Pi }^{\underset{¯}{m}}u_{\underset{¯}{m}}^i,$$
(16)
$$E^{A1}d\mathrm{\Theta }^{\underset{¯}{\mu }1}v_{\underset{¯}{\mu }}^A,E^{A2}d\mathrm{\Theta }_{\underset{¯}{\mu }}^2v_A^{\underset{¯}{\mu }}.$$
(17)
A new important property of this supervielbein (in comparison with the ”coordinate” one $`(\mathrm{\Pi }^{\underset{¯}{m}},d\mathrm{\Theta }^{\underset{¯}{\mu }1},d\mathrm{\Theta }_{\underset{¯}{\mu }}^2)`$) is that it permits covariant linear combinations of the different fermionic supervielbein forms, e.g. $`E^{A1}\pm E^{A2}`$.
The structure equations of the flat type $`IIA`$ superspace can be written as
$$dE^{(0)}=iE^{A1}E^{A1}iE^{A2}E^{A2}+E^if^i,$$
(18)
$$𝒟E^idE^i+E^jA^{ji}=iE^{A1}E^{B2}\mathrm{\Gamma }_{AB}^i+E^{(0)}f^i,$$
(19)
$$𝒟E^{A1}dE^{A1}+E^{B1}\frac{1}{4}A^{ij}\mathrm{\Gamma }_{BA}^{ij}=\frac{1}{2}E^{B1}f^i\mathrm{\Gamma }_{BA}^i,$$
(20)
$$𝒟E^{A2}dE^{A2}+E^{B2}\frac{1}{4}A^{ij}\mathrm{\Gamma }_{BA}^{ij}=\frac{1}{2}E^{B2}f^i\mathrm{\Gamma }_{BA}^i.$$
(21)
Here the ’admissible derivatives’ of the harmonics (i.e. the derivatives which preserve the conditions (6), (8))
$$du_{\underset{¯}{m}}^{\underset{¯}{a}}=u_{\underset{¯}{b}\underset{¯}{m}}\mathrm{\Omega }^{\underset{¯}{b}\underset{¯}{a}}\{\begin{array}{cc}du_{\underset{¯}{m}}^{(0)}=u_{\underset{¯}{m}}^if^i,\hfill & \\ du_{\underset{¯}{m}}^i=u_{\underset{¯}{m}}^jA^{ji}+u_{\underset{¯}{m}}^{(0)}f^i,\hfill & \end{array}$$
(22)
$$dv_{\underset{¯}{\mu }}^A\frac{1}{4}\mathrm{\Omega }^{\underset{¯}{b}\underset{¯}{a}}v_{\underset{¯}{\mu }}^B(\sigma _{\underset{¯}{b}\underset{¯}{a}})_B^A=\frac{1}{2}v_{\underset{¯}{\mu }}^Bf^i\mathrm{\Gamma }_{BA}^i\frac{1}{4}A^{ij}v_{\underset{¯}{\mu }}^B\mathrm{\Gamma }_{BA}^{ij}$$
(23)
$$dv_A^{\underset{¯}{\mu }}\frac{1}{4}\mathrm{\Omega }^{\underset{¯}{b}\underset{¯}{a}}(\sigma _{\underset{¯}{b}\underset{¯}{a}})_A^Bv_B^{\underset{¯}{\mu }}=\frac{1}{2}f^i\mathrm{\Gamma }_{AB}^iv_B^{\underset{¯}{\mu }}+\frac{1}{4}A^{ij}\mathrm{\Gamma }_{AB}^{ij}v_B^{\underset{¯}{\mu }}$$
(24)
have been used. In (22), (23), (24)
$$\mathrm{\Omega }^{\underset{¯}{a}\underset{¯}{b}}(d)=\mathrm{\Omega }^{\underset{¯}{b}\underset{¯}{a}}(d)u_{\underset{¯}{m}}^{\underset{¯}{a}}du^{\underset{¯}{b}\underset{¯}{m}}=\left(\begin{array}{cc}0& f^j\\ f^i& A^{ij}\end{array}\right)$$
(25)
are $`so(1,9)`$–valued Cartan 1–forms. The forms
$$f^iu_{\underset{¯}{m}}^{(0)}du^{\underset{¯}{m}i}$$
(26)
are covariant with respect to local $`SO(9)`$ transformations and provide a basis for the coset $`SO(1,9)/SO(9)`$ while
$$A^{ij}u_{\underset{¯}{m}}^idu^{\underset{¯}{m}j}$$
(27)
transform as $`SO(9)`$ connections. From the definition (25) one can find that the Cartan forms satisfy a zero curvature conditions (Maurer–Cartan equations)
$$d\mathrm{\Omega }^{\underset{¯}{a}\underset{¯}{b}}\mathrm{\Omega }^{\underset{¯}{a}\underset{¯}{c}}\mathrm{\Omega }_{\underset{¯}{c}}^{\underset{¯}{b}}=0\{\begin{array}{cc}𝒟f^i=df^i+f^jA^{ji}=0,\hfill & \\ F^{ij}=dA^{ij}+A^{ik}A^{kj}=f^if^j,\hfill & \end{array}$$
(28)
## 2 Gauge symmetries and equations of motion for free super–D0–brane
The simplest way to vary the geometric action is to calculate an external derivative of the Lagrangian 1–form (1)
$$_1=m\left[E^{(0)}+i\left(d\mathrm{\Theta }^{1\underset{¯}{\mu }}\mathrm{\Theta }_{\underset{¯}{\mu }}^2\mathrm{\Theta }^{1\underset{¯}{\mu }}d\mathrm{\Theta }_{\underset{¯}{\mu }}^2\right)\right],$$
(29)
(cf. (1), (16), (17)) and use the seminal formula
$$\delta _1=i_\delta (d_1)+di_\delta _1.$$
(30)
Here $`i_\delta `$ can be regarded as formal contraction of differential form with variation symbol, e.g.
$$i_\delta d\mathrm{\Theta }^{1\underset{¯}{\mu }}=\delta \mathrm{\Theta }^{1\underset{¯}{\mu }},i_\delta d\mathrm{\Theta }_{\underset{¯}{\mu }}^2=\delta \mathrm{\Theta }_{\underset{¯}{\mu }}^2,i_\delta \mathrm{\Pi }^{\underset{¯}{m}}=\delta X^{\underset{¯}{m}}i\delta \mathrm{\Theta }^1\sigma ^{\underset{¯}{m}}\mathrm{\Theta }^1i\delta \mathrm{\Theta }^2\stackrel{~}{\sigma }^{\underset{¯}{m}}\mathrm{\Theta }^2.$$
(31)
The basis (31) in the space of variations is more convenient than the original one $`(\delta X^{\underset{¯}{m}},\delta \mathrm{\Theta }^{1\underset{¯}{\mu }},\delta \mathrm{\Theta }_{\underset{¯}{\mu }}^2)`$. The contractions $`i_\delta f^i,i_\delta A^{ij}`$ shall be considered as parameters of independent transformations of the harmonic variables which preserve the conditions (6), (8) (admissible variations )
$$\delta u_{\underset{¯}{m}}^{(0)}=u_{\underset{¯}{m}}^ii_\delta f^i,\delta u_{\underset{¯}{m}}^i=u_{\underset{¯}{m}}^ji_\delta A^{ji}+u_{\underset{¯}{m}}^{(0)}i_\delta f^i,$$
(32)
$$\delta v_{\underset{¯}{\mu }}^A=\frac{1}{2}v_{\underset{¯}{\mu }}^B\mathrm{\Gamma }_{BA}^ii_\delta f^i\frac{1}{4}v_{\underset{¯}{\mu }}^B\mathrm{\Gamma }_{BA}^{ij}i_\delta A^{ij},$$
(33)
$$\delta v_A^{\underset{¯}{\mu }}=\frac{1}{2}\mathrm{\Gamma }_{AB}^iv_B^{\underset{¯}{\mu }}i_\delta f^i+\frac{1}{4}\mathrm{\Gamma }_{AB}^{ij}v_B^{\underset{¯}{\mu }}i_\delta A^{ij}.$$
External derivative of the Lagrangian 1–form (29) can be written as
$$d_1=m\left[E^if^ii(E^{A1}E^{A2})(E^{A1}E^{A2})\right].$$
(34)
Thus the variation of the action (1) is
$$S_{D0}=m_^1\left(E^ii_\delta f^if^ii_\delta E^i2i(E^{A1}E^{A2})i_\delta (E^{A1}E^{A2})\right),$$
(35)
where we skipped the complete derivative term $`_^1𝑑i_\delta _1`$. The latter means that the D0-brane worldline is considered as a surface without boundary $`^1=0`$ and, hence, there are no rejections for its identification with a boundary of some surface $`^1=^{1+1}`$ (see below).
Only $`16`$ of $`32`$ fermionic variations
$$i_\delta (E^{A1}E^{A2})\delta \mathrm{\Theta }^{1\underset{¯}{\mu }}v_{\underset{¯}{\mu }}^A\delta \mathrm{\Theta }_{\underset{¯}{\mu }}^2v_A^{\underset{¯}{\mu }}$$
(36)
are involved effectively in (35). Thus the remaining $`16`$ variations
$$\kappa ^Ai_\delta (E^{A1}+E^{A2})\delta \mathrm{\Theta }^{1\underset{¯}{\mu }}v_{\underset{¯}{\mu }}^A+\delta \mathrm{\Theta }_{\underset{¯}{\mu }}^2v_A^{\underset{¯}{\mu }}$$
(37)
can be regarded as parameters of a fermionic gauge symmetry of the model. This is the famous $`\kappa `$–symmetry <sup>2</sup><sup>2</sup>2 See for the geometrical meaning of the $`\kappa `$–symmetry.. Other gauge symmetries can be found by searching for the variations whose parameters are absent in (35). They are $`SO(9)`$ symmetry ($`i_\delta A^{ij}`$) and the reparametrization ($`i_\delta E^{(0)}=\delta X^{\underset{¯}{m}}u_{\underset{¯}{m}}^{(0)}`$, $`\delta \mathrm{\Theta }^{1,2}=0`$).
Equations of motion for the super–D0-brane appear as a result of variations with respect to $`i_\delta f^i`$, $`i_\delta E^i=\delta X^{\underset{¯}{m}}u_{\underset{¯}{m}}^i`$ and $`i_\delta (E^{A1}E^{A2})`$ respectively
$$\stackrel{~}{E}^i\stackrel{~}{\mathrm{\Pi }}^{\underset{¯}{m}}\stackrel{~}{u}_{\underset{¯}{m}}^i=0,$$
(38)
$$\stackrel{~}{f}^i\stackrel{~}{u}^{(0)\underset{¯}{m}}d\stackrel{~}{u}_{\underset{¯}{m}}^i=0,$$
(39)
$$\stackrel{~}{E}^{A1}\stackrel{~}{E}^{A2}d\stackrel{~}{\mathrm{\Theta }}^{1\underset{¯}{\mu }}\stackrel{~}{v}_{\underset{¯}{\mu }}^Ad\stackrel{~}{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2\stackrel{~}{v}_A^{\underset{¯}{\mu }}=0.$$
(40)
It can be proved that these equations are equivalent to the standard equations of motion for the super–D0–brane . In the gauge $`\stackrel{~}{X}^{\underset{¯}{m}}\stackrel{~}{u}_{\underset{¯}{m}}^{(0)}=\tau `$, $`\stackrel{~}{\mathrm{\Theta }}^{1\underset{¯}{\mu }}=0`$ Eqs. (38), (39), (40) are equivalent to the set
$$d\stackrel{~}{X}^{\underset{¯}{m}}=d\tau p^{\underset{¯}{m}}/m,dp_{\underset{¯}{m}}=0,p_{\underset{¯}{m}}^2=m^2,d\stackrel{~}{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2=0,$$
(41)
which describes a massive superparticle.
## 3 Geometric action for type $`IIA`$ superstring
The geometric action, superembedding approach and generalized action principle for type $`IIB`$ superstring has been constructed in respectively. Type $`IIA`$ superstring has not been considered in this framework before.
The geometric action for type IIA superstring is
$$S_{IIA}=_{^{1+1}}\widehat{}_2=_{^{1+1}}\left(\frac{1}{2}\widehat{E}^{++}\widehat{E}^{}\widehat{B}_2\right),$$
(42)
where
$$E^{\pm \pm }\mathrm{\Pi }^{\underset{¯}{m}}U_{\underset{¯}{m}}^{\pm \pm },E^I\mathrm{\Pi }^{\underset{¯}{m}}U_{\underset{¯}{m}}^I,$$
(43)
$$B_2=i\mathrm{\Pi }^{\underset{¯}{m}}\left(d\mathrm{\Theta }^1\sigma _{\underset{¯}{m}}\mathrm{\Theta }^1id\mathrm{\Theta }^2\stackrel{~}{\sigma }_{\underset{¯}{m}}\mathrm{\Theta }^2\right)+d\mathrm{\Theta }^1\sigma ^{\underset{¯}{m}}\mathrm{\Theta }^1d\mathrm{\Theta }^2\stackrel{~}{\sigma }_{\underset{¯}{m}}\mathrm{\Theta }^2,$$
(44)
$$\widehat{\mathrm{\Pi }}^{\underset{¯}{m}}=d\widehat{X}^{\underset{¯}{m}}id\widehat{\mathrm{\Theta }}^1\sigma ^{\underset{¯}{m}}\widehat{\mathrm{\Theta }}^1id\widehat{\mathrm{\Theta }}^2\stackrel{~}{\sigma }^{\underset{¯}{m}}\widehat{\mathrm{\Theta }}^2=d\xi ^m\widehat{\mathrm{\Pi }}_m^{\underset{¯}{m}}(\xi )$$
(45)
is the pull-back of the 1–form (2) on the superstring worldsheet $`^{1+1}`$ whose embedding into the type $`IIA`$ superspace $`^{1+1}\underset{¯}{}^{(10|32)}`$ is defined by
$$X^{\underset{¯}{m}}=\widehat{X}^{\underset{¯}{m}}(\xi )\widehat{X}^{\underset{¯}{m}}(\tau ,\sigma ),\mathrm{\Theta }^{1\underset{¯}{\mu }}=\widehat{\mathrm{\Theta }}^{1\underset{¯}{\mu }}(\xi ),\mathrm{\Theta }_{\underset{¯}{\mu }}^2=\widehat{\mathrm{\Theta }}_{\underset{¯}{\mu }}^2(\xi ).$$
(46)
$`\widehat{U}_{\underset{¯}{m}}^{\pm \pm }(\xi )\widehat{U}_{\underset{¯}{m}}^0(\xi )\pm \widehat{U}_{\underset{¯}{m}}^9(\xi )`$ are light–like Lorentz harmonic vectors , i.e. the components of the $`SO(1,9)`$ valued matrix
$$U_{\underset{¯}{m}}^{\underset{¯}{a}}=(U_{\underset{¯}{m}}^0,U_{\underset{¯}{m}}^I,U_{\underset{¯}{m}}^9)=(\frac{1}{2}(U_{\underset{¯}{m}}^{++}+U_{\underset{¯}{m}}^{}),U_{\underset{¯}{m}}^I,\frac{1}{2}(U_{\underset{¯}{m}}^{++}U_{\underset{¯}{m}}^{}))SO(1,9)$$
(47)
The spinor harmonics
$$V_{\underset{¯}{\mu }}^{\underset{¯}{\alpha }}=(V_{\underset{¯}{\mu }q}^+,V_{\underset{¯}{\mu }\dot{q}}^{})^TSpin(1,9)$$
(48)
$$V_{\underset{¯}{\alpha }}^{\underset{¯}{\mu }}=(V_q^{\underset{¯}{\mu }},V_{\dot{q}}^{+\underset{¯}{\mu }})Spin(1,9)$$
(49)
$$V_{\underset{¯}{\alpha }}^{\underset{¯}{\mu }}V_{\underset{¯}{\mu }}^{\underset{¯}{\beta }}=\delta _{\underset{¯}{\alpha }}^{\underset{¯}{\beta }}:V_q^{\underset{¯}{\mu }}V_{\underset{¯}{\mu }p}^+=\delta _{qp},V_{\dot{q}}^{+\underset{¯}{\mu }}V_{\underset{¯}{\mu }\dot{p}}^{}=\delta _{\dot{q}\dot{p}},V_q^{\underset{¯}{\mu }}V_{\underset{¯}{\mu }\dot{p}}^{}=V_{\dot{p}}^{+\underset{¯}{\mu }}V_{\underset{¯}{\mu }q}^+=0$$
(50)
are related with (47) by Eqs. (9) which include, in particular,
$$U_{\underset{¯}{m}}^{++}\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}=2V_{\underset{¯}{\mu }q}^+V_{\underset{¯}{\nu }q}^+,U_{\underset{¯}{m}}^{}\stackrel{~}{\sigma }^{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}=2V_q^{\underset{¯}{\mu }}V_q^{\underset{¯}{\nu }},$$
(51)
$$U_{\underset{¯}{m}}^{}\sigma _{\underset{¯}{\mu }\underset{¯}{\nu }}^{\underset{¯}{m}}=2V_{\underset{¯}{\mu }\dot{q}}^{}V_{\underset{¯}{\nu }\dot{q}}^{},U_{\underset{¯}{m}}^{++}\stackrel{~}{\sigma }^{\underset{¯}{m}\underset{¯}{\mu }\underset{¯}{\nu }}=2V_{\dot{q}}^{+\underset{¯}{\mu }}V_{\dot{q}}^{+\underset{¯}{\nu }}.$$
(52)
The details about the stringy harmonics and Cartan forms (cf. (25))
$$\mathrm{\Omega }^{\underset{¯}{a}\underset{¯}{b}}U_{\underset{¯}{m}}^{\underset{¯}{a}}dU^{\underset{¯}{b}\underset{¯}{m}}=\left(\begin{array}{ccc}0& \frac{f^{++J}+f^J}{2}& \frac{1}{2}\omega \\ \frac{f^{++I}+f^I}{2}& A^{IJ}& \frac{f^{++I}f^I}{2}\\ \frac{1}{2}\omega & \frac{f^{++J}f^J}{2}& 0\end{array}\right)$$
(53)
can be found in Refs. .
The external derivative of the Lagrangian 2–form
$$_2=\frac{1}{2}E^{++}E^{}i\mathrm{\Pi }^{\underset{¯}{m}}\left(d\mathrm{\Theta }^1\sigma _{\underset{¯}{m}}\mathrm{\Theta }^1id\mathrm{\Theta }^2\stackrel{~}{\sigma }_{\underset{¯}{m}}\mathrm{\Theta }^2\right)+d\mathrm{\Theta }^1\sigma ^{\underset{¯}{m}}\mathrm{\Theta }^1d\mathrm{\Theta }^2\stackrel{~}{\sigma }_{\underset{¯}{m}}\mathrm{\Theta }^2$$
(54)
can be calculated with the use of (51), (52), (53) and stringy counterparts of Eqs. (22)– (28)
$$d_2=2iE^{++}E^{\dot{q}1}E^{\dot{q}1}+2iE^{}E_{\dot{q}}^{+2}E_{\dot{q}}^{+2}+$$
(55)
$$E^I\left(\frac{1}{2}E^{}f^{++I}\frac{1}{2}E^{++}f^I+2i\left(E^{+q1}E^{\dot{q}1}+E_q^{+2}E_{\dot{q}}^2\right)\gamma _{q\dot{q}}^I\right).$$
The parameters of the stringy $`\kappa `$–symmetry can be identified with the contractions of those fermionic forms which are absent in the first line of Eq. (55)
$$\kappa ^{+q}i_\delta E^{+q1}=\delta \mathrm{\Theta }^{1\underset{¯}{\mu }}v_{\underset{¯}{\mu }q}^+,\kappa ^qi_\delta E_q^2=\delta \mathrm{\Theta }_{\underset{¯}{\mu }}^2v_q^{\underset{¯}{\mu }}.$$
(56)
The second line of (55) determines, in particular, the transformations of the harmonics with respect to the $`\kappa `$–symmetry. Other gauge symmetries are $`SO(1,1)\times SO(8)`$ ($`i_\delta \omega ,i_\delta A^{IJ}`$) and the reparametrization ($`i_\delta E^{\pm \pm }=\delta X^{\underset{¯}{m}}U_{\underset{¯}{m}}^{\pm \pm }`$).
The equations of motion for the type $`IIA`$ superstring can be obtained from (55). They are
$$\widehat{E}^I\widehat{\mathrm{\Pi }}^{\underset{¯}{m}}u_{\underset{¯}{m}}^I=0,$$
(57)
$$M_2^IE^{}f^{++I}E^{++}f^I+4i\left(E^{+q1}E^{\dot{q}1}+E_q^{+2}E_{\dot{q}}^2\right)\gamma _{q\dot{q}}^I=0,$$
(58)
$$E^{++}E^{\dot{q}1}\widehat{\mathrm{\Pi }}^{\underset{¯}{m}}d\mathrm{\Theta }^{1\underset{¯}{\mu }}v_{\underset{¯}{\mu }\dot{q}}^{}u_{\underset{¯}{m}}^{++}=0,$$
(59)
$$E^{}E_{\dot{q}}^{+2}\widehat{\mathrm{\Pi }}^{\underset{¯}{m}}d\mathrm{\Theta }_{\underset{¯}{\mu }}^2v_{\dot{q}}^{+\underset{¯}{\mu }}u_{\underset{¯}{m}}^{}=0.$$
(60)
It can be proved that this set is equivalent to the standard equations of motion for the type $`IIA`$ superstring (see for the type $`IIB`$ case).
## 4 Supersymmetric action functional for type $`IIA`$ superstring with super–D0–branes at the endpoints
The main problem which should be solved to write down the action of interacting branes is: how to take into account an identification of the bosonic and fermionic superembedding functions on the intersection . However, this problem has the natural solution just for the system under consideration. Here the super–D0–brane worldline $`^1`$ should be considered as the boundary of the superstring worldsheet $`^1=^{1+1}`$. Thus one can define an embedding of $`^1`$ into $`^{1+1}`$
$$\xi ^m=\stackrel{~}{\xi }^m(\tau ):^1=^{1+1}^{1+1}$$
(61)
and identify the super–D0–brane coordinate functions $`\stackrel{~}{X}(\tau ),\stackrel{~}{\mathrm{\Theta }}^{1,2}(\tau )`$ with the images of the superstring coordinate functions $`\widehat{X}(\xi ),\widehat{\mathrm{\Theta }}^{1,2}(\xi )`$ on the boundary
$$\stackrel{~}{X}(\tau )=\widehat{X}\left(\stackrel{~}{\xi }(\tau )\right),\stackrel{~}{\mathrm{\Theta }}^{1,2}(\tau )=\widehat{\mathrm{\Theta }}^{1,2}\left(\stackrel{~}{\xi }(\tau )\right).$$
(62)
With this identification the action for the coupled system of an open fundamental superstring and super–D0–branes at the ends of the superstring is the direct sum of the actions (1) and (42)
$$S_{str+D0}=_{^{1+1}}\widehat{}_2+m_{^{1+1}}\stackrel{~}{}_1.$$
(63)
The variation of the action can be calculated as
$$\delta S_{str+D0}=_{^{1+1}}i_\delta \left(d_2\right)+m_{^{1+1}}\left(i_\delta _2+i_\delta d_1\right).$$
(64)
A possibility is to require the vanishing of the bulk and the boundary variations in (64) separately. On the other hand, following , one can introduce the following current density distribution
$$j_1=d\xi ^m\epsilon _{mn}_^2𝑑\stackrel{~}{\xi }^n(\tau )\delta ^2\left(\xi \stackrel{~}{\xi }(\tau )\right)=\epsilon _{mn}d\xi ^nj^m$$
(65)
with the property
$$_{^{1+1}}j_1\widehat{𝒜}_1=_{^{1+1}}\stackrel{~}{𝒜}_1.$$
(66)
In (66) $`\widehat{𝒜}_1`$ is an arbitrary 1-form defined on the worldsheet $`^{1+1}`$ and $`\stackrel{~}{𝒜}_1`$ is its pull–back onto $`^1=^{1+1}`$. As it is easy to define the extension of the super–D0–brane Lagrangian form to the whole worldsheet (29), we can use (66) to lift the action (63) or its variation (64) to the integral over the whole worldsheet
$$S_{str+D0}=_{^{1+1}}\widehat{}_2+mj_1\widehat{}_1,$$
(67)
$$\delta S_{str+D0}=_{^{1+1}}i_\delta \left(d_2\right)+mj_1\left(i_\delta _2+i_\delta d_1\right).$$
(68)
Here $`\widehat{}_2`$ is defined by Eq. (42), (54), $`\widehat{}_1`$ is the pull–back of the 1–form $`_1`$ (29) on the worldsheet (46). In (68) the contractions of the forms (i.e. $`i_\delta _2=i_\delta (1/2dZ^MdZ^N_{NM})=dZ^{\underset{¯}{M}}\delta Z^{\underset{¯}{N}}_{\underset{¯}{N}\underset{¯}{M}}`$ ) should be pulled back to the worldsheet and the contractions of the coordinate differentials in the second term should include the variations of the functions $`\stackrel{~}{\xi }^m(\tau )`$ (61), e.g. $`i_\delta dX^{\underset{¯}{m}}=\delta \stackrel{~}{X}^{\underset{¯}{m}}+\delta \stackrel{~}{\xi }^m(\tau )_m\widehat{X}^{\underset{¯}{m}}|_{\xi =\stackrel{~}{\xi }(\tau )}`$ (though the variations $`\delta \stackrel{~}{\xi }^m(\tau )`$ shall not produce independent equations).
## 5 D0-branes at the endpoints of bosonic string
In the pure bosonic limit our dynamical system describes a ’string with masses at the endpoints’. Such system has been studied in early years of QCD strings and partial solutions have been found . However the geometric or first order formulation as well as ’extended variational problem’ approach (67), (68) for this system are new and, in our opinion, instructive.
The geometric action for the coupled system has the form
$$S=_{^{1+1}}\frac{1}{2}\widehat{E}^{++}\widehat{E}^{}+mj_1\widehat{E}^{(0)}.$$
(69)
Its variation with respect to harmonic variables
$$\delta _hS=_{^{1+1}}\left(\frac{1}{2}\widehat{E}^I\widehat{E}^{}i_\delta f^{++I}\frac{1}{2}\widehat{E}^I\widehat{E}^{}i_\delta f^{++I}+mj_1\widehat{E}^ii_\delta f^i\right)$$
(70)
produces the same algebraic embedding equations as in the case of a free string and free D0-brane(s)
$$\widehat{E}^Id\widehat{X}^{\underset{¯}{m}}(\xi )U_{\underset{¯}{m}}^I(\xi )=0,$$
(71)
$$\stackrel{~}{E}^id\widehat{X}^{\underset{¯}{m}}(\xi (\tau ))u_{\underset{¯}{m}}^i(\tau )=0.$$
(72)
This provides the possibility to simplify the variation with respect to the coordinate functions considering it modulo Eqs. (71), (72)
$$\delta S|_{\widehat{E}^I=0=\stackrel{~}{E}^i}=\frac{1}{2}_{^{1+1}}\left(\widehat{M}_2^IU_{\underset{¯}{m}}^I+j_1\left(\widehat{E}^{++}U_{\underset{¯}{m}}^{}\widehat{E}^{}U_{\underset{¯}{m}}^{++}2mf^iu_{\underset{¯}{m}}^i\right)\right)\delta \widehat{X}^{\underset{¯}{m}}$$
(73)
Here
$$M_2^IE^{}f^{++I}E^{++}f^I$$
(74)
is the pure bosonic limit of the l.h.p. of the free superstring equation (58).
When considered together, Eqs. (71) and (72) relate the images of the stringy harmonics on the worldsheet boundary with the harmonics of the D0-branes. Indeed, Eq. (72) implies $`d\stackrel{~}{X}^{\underset{¯}{m}}=E^{(0)}u^{(0)\underset{¯}{m}}`$. Substituting this relation into the pull-back of Eq. (71) on the boundary one arrives at
$`u^{(0)\underset{¯}{m}}(\tau )U_{\underset{¯}{m}}^I(\xi (\tau ))=0`$. This implies
$$u^{(0)\underset{¯}{m}}(\tau )=w^{}(\tau )U^{++\underset{¯}{m}}(\xi (\tau ))+U^{++\underset{¯}{m}}(\xi (\tau ))/4w^{}(\tau )$$
(75)
with some indefinite function $`w^{}(\tau )`$ (compensator for $`SO(1,1)`$ symmetry). The relative coefficient in (75) is fixed by normalization conditions. Now, using the $`SO(9)`$ gauge symmetry, we can chose the super–D0–brane harmonics to be expressed through the images of stringy ones by
$$u^{i\underset{¯}{m}}(\tau )=(U^{++\underset{¯}{m}}(\xi (\tau )),u^{(9)\underset{¯}{m}}(\tau )),$$
(76)
$$u^{(9)\underset{¯}{m}}(\tau )=w^{}(\tau )U^{++\underset{¯}{m}}(\xi (\tau ))U^{++\underset{¯}{m}}(\xi (\tau ))/4w^{}(\tau ).$$
Then the set of Cartan forms (26) splits as $`f^i=(f^I,f^{(9)})`$ and, after some algebraic manipulations, the equations of motion can be written in the form
$$M_2^IE^{}f^{++I}E^{++}f^I=mj_1f^I,$$
(77)
$$j_1\left(E^{(0)}mf^{(9)}\right)=0.$$
(78)
When $`m0`$ the latter equation evidently implies
$$f^{(9)}u^{(0)\underset{¯}{m}}du_{\underset{¯}{m}}^{(9)}=\frac{1}{m}E^{(0)}.$$
(79)
For $`m\mathrm{}`$ we can neglect the left hand side of Eq. (77) and the right hand side of Eq.(79). Thus we arrive at the free equations of motion (39) for the D0-branes. This means that D0–branes (or ’quarks’ ) with infinite mass(es) do not feel the influence of the open string. When $`m0`$ Eq. (77) becomes the free string equation, while (78) implies that $`j_1=0`$, i.e. that the worldsheet has no boundary and, thus, the string is closed.
## Concluding Remarks
The analyzes of the gauge symmetries of the action (63), (67) and the supersymmetric equations which follow from it will be the subject of a forthcoming article. We expect that they shall provide an important insights for future study of the generic system of interacting superbranes.
Another direction of the development of the present results is to elaborate the generalized action principle and the superembedding approach the super–D0–brane (see for the super-Dp-branes with $`0<p<9`$ and for $`p=9`$). The basis for such study is provided by the geometric action (1).
## Acknowledgments
The author is grateful to D. Sorokin, M. Tonin, B. Julia for useful conversations and for the hospitality at the Padova Section of INFN (Padova) and Laboratoire de Physique Theorique de l’Ecole Normale Superieure (Paris), where a part of this work has been done. A partial support from the INTAS Grant 96-308 and the Ukrainian GKNT Grant 2.5.1/52 is acknowledged. |
warning/0001/hep-th0001176.html | ar5iv | text | # Pre-Big Bang Cosmology and Quantum Fluctuations
## 1 INTRODUCTION
In the framework of string theory, the Pre-Big Bang scenario provides an alternative to the standard inflationary paradigm. In the Pre-Big Bang model <sup>1</sup><sup>1</sup>1 An updated collection of papers on the PBB scenario is available at http://www.to.infn.it/$``$ gasperin/. the inflationary solutions, driven by the kinetic energy of the dilaton, emerge naturally via the duality symmetries of string theory.
It has, however, been argued that, even though the two classical moduli of the open ($`𝒦=1`$), homogeneous and isotropic solution lie deeply inside the perturbative regime, the vacuum quantum fluctuations drastically modify the classical behaviour preventing the occurrence of an appreciable amount of inflation.
Quantum fluctuations in a non-spatially flat Universe are considerably harder to study than in the flat case . In Ref. we thoroughly studied the quantum fluctuations around the $`𝒦=1`$ solution . In that work we showed that the perturbation equations can be exactly integrated in terms of standard hypergeometric functions. We found that particle production (i.e. the amplification of vacuum fluctuations) is strongly suppressed at very early times and remains small through the whole perturbative PBB phase, and hence, does not impede the occurrence of PBB inflation.
## 2 THE SECOND-ORDER ACTION
The (string-frame) open homogeneous, isotropic PBB-type solution was first found in and then rederived and discussed in . The solution contains two arbitrary moduli, $`L`$ and $`\varphi _{\mathrm{in}}`$, reflecting the symmetries of the classical equations under a constant shift of the dilaton and a constant rescaling of the metric. These two parameters are to be chosen appropriately (see Refs. ) in order to ensure the occurrence of a sufficient amount of PBB inflation. Such a solution describes a universe which is almost trivial (Milne-like) at $`\eta \mathrm{}`$ and inflating in $`\mathrm{}<\eta <𝒪(1)`$, having an initial curvature $`𝒪(L^2)`$ and coupling $`𝒪(\mathrm{exp}(\varphi _{\mathrm{in}}/2))`$, until it enters the strong curvature and/or strong coupling regime at $`\eta \eta _1`$. The critical value $`\eta _1`$ is easily determined in terms of the integration constants $`L`$ and $`\varphi _{\mathrm{in}}`$ as $`(\eta _1)=\mathrm{max}(e^{\varphi _{\mathrm{in}}/\sqrt{3}},(\mathrm{}_s/L)^{1+1/\sqrt{3}})`$.
It is well known that studying perturbations is technically simpler in the so-called Einstein-frame which is related to the string-frame by a conformal transformation. The action is
$`S^{(E)}={\displaystyle \frac{1}{2\mathrm{}_P^2}}{\displaystyle d^4x\sqrt{g}\left(R(g)\frac{1}{2}(\varphi )^2\right)},`$ (1)
where $`\varphi _{\mathrm{today}}`$ is the present value of the dilaton, $`\mathrm{}_P\sqrt{8\pi G}=\mathrm{exp}(\varphi _{\mathrm{today}}/2)\mathrm{}_s0.1\mathrm{}_s`$ refers to the present value of the Planck-length with $`\mathrm{}=1`$. Usually one computes perturbations in the Einstein frame and then transforms the results back to the string frame for a physical interpretation.
The $`𝒦=1`$ solution is:
$`a(\eta )`$ $`=`$ $`\mathrm{}(\mathrm{sinh}\eta \mathrm{cosh}\eta )^{\frac{1}{2}}`$
$`\varphi (\eta )`$ $`=`$ $`\sqrt{3}\mathrm{ln}(\mathrm{tanh}\eta )+\varphi _{\mathrm{in}},\eta <0,`$ (2)
where the modulus $`\mathrm{}`$ is given by $`\mathrm{}^2=L^2\mathrm{exp}(\varphi _{\mathrm{today}}\varphi _{\mathrm{in}})`$. Generic perturbations are defined by
$`g_{\mu \nu }=g_{\mu \nu }^{(0)}+\delta g_{\mu \nu },\varphi =\varphi ^{(0)}+\delta \varphi `$ (3)
where superscript $`(0)`$ refers to the background solution and we shall use isotropic-spatial coordinates.
## 3 QUANTUM FLUCTUATIONS
### 3.1 Tensor perturbations
Since the tensor metric perturbations are automatically gauge-invariant and decoupled from the scalar perturbations, they are easier to study. They are defined as
$`\delta g_{\mu \nu }^{(\mathrm{T})}=\mathrm{diag}(0,a^2h_{ij}),`$ (4)
where the symmetric three-tensor $`h_{ij}`$ satisfies the transverse-traceless (TT) conditions.
We then find:
$`\delta ^{(2)}S^{(T)}`$ $`=`$ $`{\displaystyle \frac{1}{4\mathrm{}_P^2}}{\displaystyle }d^4x\sqrt{\gamma }a^2(h^{ij}h_{ij}^{}`$ (5)
$``$ $`^lh^{ij}_lh_{ij}2𝒦h^{ij}h_{ij}).`$
By expanding the tensor perturbations in TT tensor-pseudospherical harmonics (as $`𝒦=1`$) , we eventually get the simple equation
$`u_{nlm}^{\prime \prime }+\left(n^2+{\displaystyle \frac{1}{12}}\varphi ^2\right)u_{nlm}=0,`$ (6)
where $`u_{nlm}ah_{nlm}`$ is the canonical variable of perturbation. For the background (2) Eq. (6) can be exactly solved in terms of the standard hypergeometric function as
$`u_N(\eta )=C_1[\mathrm{csch}^2(2\eta )]^{\frac{in}{4}}\times `$
$`F[{\displaystyle \frac{1in}{4}},{\displaystyle \frac{1in}{4}},{\displaystyle \frac{2in}{2}},\mathrm{csch}^2(2\eta )]`$
$`+C_2\mathrm{c}.\mathrm{c}.,`$ (7)
where $`N`$ stands for the collection of indices $`(nlm)`$ and $`C_{1,2}`$ are (classically arbitrary) integration constants. At early times, $`n^2\varphi ^2`$, and thus $`u`$ is a free canonical field. Hence, imposing the standard commutation relations, as $`\eta \mathrm{}`$, we get
$`u_N(\eta )u_N^{\mathrm{}}(\eta ){\displaystyle \frac{2\mathrm{}_P}{\sqrt{n}}}e^{in\eta }.`$ (8)
Since $`F[a,b,c,0]=1`$, Eq. (8) fixes the integration constants as $`|C_1|=2\mathrm{}_P/\sqrt{n},C_2=0`$. The deviation from a trivial plane-wave behaviour can easily be computed from the small argument limit of $`F`$. We find
$`u_N(\eta )=u_N^{\mathrm{}}(\eta )\left(1+\alpha _ne^{4\eta i\beta _n}\right),`$ (9)
where $`\alpha _n,\beta _n`$ are $`n`$-dependent constants fixed from the Taylor expansion of the hypergeometric function. It is worth noting that the correction to the vacuum amplitude dies off as $`e^{4\eta }`$, i.e. as $`t^4`$ in terms of cosmic time $`te^\eta `$.
We can also estimate the behaviour of the solution near the singularity, i.e. for $`\eta 0`$. By virtue of the small $`\eta `$ behaviour $`a\mathrm{}|\eta |^{1/2}`$, we find
$`|h_N|2\sqrt{{\displaystyle \frac{2}{\pi }}}{\displaystyle \frac{\mathrm{}_P}{\mathrm{}}}\sqrt{\mathrm{coth}\left({\displaystyle \frac{n\pi }{2}}\right)}\mathrm{ln}|\eta |.`$ (10)
We shall come back to this result after deriving a similar expression for scalar perturbations.
### 3.2 Scalar perturbations
Consider now scalar metric-dilaton perturbations (3). The scalar part of metric perturbations is defined by
$`\delta g_{\mu \nu }^{(\mathrm{S})}a^2(\eta )\left(\begin{array}{cc}2\phi & _iB\\ _iB& 2(\psi \gamma _{ij}+_i_jE)\end{array}\right).`$ (11)
In the second-order action the variables $`B,\phi `$ are Lagrange multipliers, providing two constraints. We can introduce the gauge-invariant variable $`\mathrm{\Psi }`$ by (see Ref. )
$`\mathrm{\Psi }={\displaystyle \frac{4}{\varphi ^{}}}\left[\psi +(BE^{})\right],`$ (12)
and, after using the constraints, the action reads
$`\delta ^{(2)}S^{(S)}={\displaystyle \frac{1}{2\mathrm{}_P^2}}{\displaystyle }d^4xa^2\sqrt{\gamma }\times `$
$`(^2+3𝒦)\mathrm{\Psi }\left[_\eta ^2^2+2(^{}+𝒦)\right]\mathrm{\Psi }.`$ (13)
One can now make use of the constraints to eliminate the variable $`(BE^{})`$ from the action (13) in terms of $`\phi ,\psi `$ and $`\delta \varphi `$. The latter variables are not independent either, being related by a linear combination of the two constraints. After its implementation the action (13) contains only true degrees of freedom.
As was in the case of tensor perturbations, we introduce a canonical field $`\mathrm{\Psi }_c`$ and expand it as
$`\mathrm{\Psi }_ca\mathrm{\Psi }={\displaystyle 𝑑n\underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\mathrm{\Psi }_{nlm}(\eta )Q_{nlm}(𝐱)},`$ (14)
where $`Q_{nlm}(𝐱)`$ are the scalar pseudospherical harmonics . We then get a simple equation for $`\overline{\mathrm{\Psi }}_N\sqrt{n^2+4}\mathrm{\Psi }_N`$, namely
$`\overline{\mathrm{\Psi }}_N^{\prime \prime }+(n^2{\displaystyle \frac{1}{4}}\varphi ^2)\overline{\mathrm{\Psi }}_N=0,`$ (15)
where we must impose, as $`\eta \mathrm{}`$,
$`\overline{\mathrm{\Psi }}_N(\eta )\overline{\mathrm{\Psi }}_N^{\mathrm{}}(\eta ){\displaystyle \frac{\mathrm{}_P}{\sqrt{n}}}e^{in\eta }.`$ (16)
Eq. (15) can again be transformed (for the background (2)) into a hypergeometric equation. We hence find
$`\overline{\mathrm{\Psi }}_N(\eta )=\stackrel{~}{C_1}[\mathrm{csch}^2(2\eta )]^{\frac{in}{4}}\times `$
$`F[{\displaystyle \frac{1in}{4}},{\displaystyle \frac{3in}{4}},{\displaystyle \frac{2in}{2}},\mathrm{csch}^2(2\eta )]`$
$`+\stackrel{~}{C_2}\mathrm{c}.\mathrm{c}.,`$ (17)
where, as before, we have to take $`|\stackrel{~}{C_1}|=\mathrm{}_P/\sqrt{n},\stackrel{~}{C_2}=0`$. Corrections to the free plane wave can be easily computed and, again, are suppressed by four powers of $`1/t`$:
$`\overline{\mathrm{\Psi }}_N(\eta )=\overline{\mathrm{\Psi }}_N^{\mathrm{}}(\eta )\left(1+\stackrel{~}{\alpha }_ne^{4\eta i\stackrel{~}{\beta }_n}\right),`$ (18)
where $`\overline{\mathrm{\Psi }}_N^{\mathrm{}}`$ is given by (16) and $`\stackrel{~}{\alpha }_n,\stackrel{~}{\beta }_n`$ are $`n`$-dependent constants fixed from the expansion of the hypergeometric function.
Estimating the behaviour of (17) near $`\eta 0`$ , we obtain:
$`|\overline{\mathrm{\Psi }}_N|\mathrm{}_P\sqrt{{\displaystyle \frac{n^2+1}{2\pi }}}\sqrt{\mathrm{coth}\left({\displaystyle \frac{n\pi }{2}}\right)}\times `$
$`\left(|\eta |^{3/2}\mathrm{ln}|\eta |+{\displaystyle \frac{2}{n^2+1}}|\eta |^{1/2}\right).`$ (19)
## 4 CONCLUSIONS
Let us choose the off-diagonal gauge , defined by setting $`\psi =E=0`$ in (11). By using Eq. (12) one can reconstruct the scalar field fluctuation $`\delta \varphi `$ from $`\mathrm{\Psi }`$ as
$`\delta \varphi =\mathrm{\Psi }^{}+{\displaystyle \frac{𝒦^{}}{}}\mathrm{\Psi },`$ (20)
implying that $`\delta \varphi `$ represents, in this gauge, a gauge-invariant object.
In the presence of spatial curvature, the field $`v=a\delta \varphi `$ plays the role of the canonical field in the far past, when $`\eta `$ is large and negative. Eq. (20) tells us that the behaviour of $`v`$ in the far past follows directly from that of $`\overline{\mathrm{\Psi }}_N`$, given in Eqs. (16) and (18):
$`v^{\mathrm{}}(\eta ){\displaystyle \frac{\mathrm{}_P}{\sqrt{n}}}\sqrt{{\displaystyle \frac{2in}{2+in}}}e^{in\eta }.`$ (21)
Corrections to (21) are again suppressed as $`t^4`$
$`v(\eta )=v^{\mathrm{}}(\eta )\left(1+\widehat{\alpha }_ne^{4\eta i\widehat{\beta }_n}\right),`$ (22)
where $`\widehat{\alpha }_n,\widehat{\beta }_n`$ are $`n`$-dependent constants.
The behaviour of $`\delta \varphi `$ near $`\eta 0`$ is :
$`|\delta \varphi _N|{\displaystyle \frac{\mathrm{}_P}{\mathrm{}}}\sqrt{{\displaystyle \frac{n^2+1}{2\pi }}}\sqrt{{\displaystyle \frac{\mathrm{coth}(\frac{n\pi }{2})}{n^2+4}}}\mathrm{ln}|\eta |.`$ (23)
Lastly, let us compare the energy contained in the quantum fluctuations of the dilaton and that in the classical solution near the singularity. Note that the expansion (19) can be trusted only up to some maximum $`n`$ for which $`1n_{\mathrm{max}}1/|\eta |`$. Consequently, the ratio of the kinetic energy densities near $`|\eta |0`$ (up to constant prefactors of $`𝒪(1)`$) becomes
$`{\displaystyle \frac{_\mathrm{Q}}{_\mathrm{C}}}={\displaystyle \frac{d^3x\sqrt{\gamma }a^2(\delta \varphi ^{})^2}{d^3x\sqrt{\gamma }a^2\varphi _{}^{}{}_{}{}^{2}}}{\displaystyle \frac{\mathrm{}_P^2}{\mathrm{}^2}}{\displaystyle ^{n_{\mathrm{max}}}}{\displaystyle \frac{dn}{n}}n^3.`$ (24)
We can express the above result in terms of the value of the physical Hubble parameter $`H(\eta )/a`$ at horizon crossing of the scale $`n`$, $`H_{\mathrm{HC}}(n)`$, i.e.
$`H_{\mathrm{HC}}(n){\displaystyle \frac{1}{\eta a}}(\eta 1/n)n^{3/2}/\mathrm{}.`$ (25)
Thus Eq. (24) takes the suggestive form
$`{\displaystyle \frac{_\mathrm{Q}}{_\mathrm{C}}}=\mathrm{}_P^2{\displaystyle ^{n_{\mathrm{max}}}}{\displaystyle \frac{dn}{n}}H_{\mathrm{HC}}^2(n).`$ (26)
In order to draw physical conclusion we should transform the results back to the string frame. However, in our case, this is hardly necessary. As far as the importance of vacuum fluctuations is concerned, as $`\eta 0`$, the final result (26) expresses the relative importance of quantum and classical fluctuations near the singularity in terms of a frame-independent quantity: the ratio of the effective Planck length to the size of the horizon. Since, by definition of the perturbative dilaton phase, the Hubble radius is always larger than the string scale, the relative importance of quantum fluctuations is always bounded by the ratio $`\mathrm{}_P/\mathrm{}_s`$ which is always less than one in the perturbative phase.
Let us now come to the more subtle issue of the far-past behaviour of tensor and scalar quantum fluctuations. Computations may be done in either frame, since the dilaton is approximately constant in the far past. Our results, expressed in Eqs. (9) and (22), show that corrections to the trivial quantum fluctuations are of relative order $`e^{4\eta }t^4`$, i.e. of order $`t^3`$ relative to the (homogeneous) classical perturbation. This suggests that quantum effects do not modify appreciably classical behaviour in the far past, contrary to the claim of . This result is also supported by the structure of the superstring one-loop effective-action (which is well-defined thanks to the string cutoff). Because of supersymmetry, neither a cosmological term nor a renormalization of Newton’s constant are generated at one-loop, but only terms containing at least four derivatives. Thus, quantum corrections to early-time classical behaviour are of relative order $`t^6`$, i.e just like our corrections $`(\delta \varphi ^{}/\varphi ^{})^2`$. Note also that that generating a cosmological constant by quantum corrections would upset completely the whole PBB scenario. |
warning/0001/math0001018.html | ar5iv | text | # Path-wise solutions of SDEs driven by Lévy processes
## Introduction
In this paper I give a path-wise method for solving the following integral equation:
$$Y_t=Y_0+_0^tf(Y_t)𝑑X_tY_0=a^d.$$
(1)
when the driving process is a Lévy process.
Typically, a Lévy process a.s. has unbounded variation. The integral does not exist in a Lebesgue-Stieltjes sense. However, the integral still makes sense as a random variable due to the stochastic calculus of semi-martingales developed by the Strasbourg school .
The semi-martingale integration theory is not complete though. There are processes of interest which do not fit into the semi-martingale framework, for example the fractional Brownian motion. An alternative integral is provided by the path-wise approach studied by Lyons , and Dudley . The basis of their papers is that of Young , who showed that the integral
$$_0^tf𝑑g$$
(2)
is defined whenever $`f`$ and $`g`$ have finite $`p`$ and $`q`$-variation for $`1/p+1/q>1`$ (and they have no common discontinuities). For a comprehensive overview of the theory we recommend the lecture notes of Dudley and Norvaiša .
Recently in , a system of linear Riemann-Stieltjes integral equations is solved when the integrator has finite $`p`$-variation for some $`0<p<2`$. These results are contained in Theorem 1.1 where we allow non-linearity of the vector field $`f`$. This is because our approach is an extension of the method of , .
The approach that I follow distinguishes two cases. The first is when the process has finite $`p`$-variation a.s., for some $`p<2`$. We use the Young integral . In (1) is solved when $`X_t`$ is a continuous path of finite $`p`$-variation for some $`p<2`$.
The second case is when the process has finite $`p`$-variation a.s., for some $`p>2`$. The Young integral is only defined when $`f`$ and $`g`$ have finite $`p`$ and $`q`$-variation for $`1/p+1/q>1`$. So an iteration scheme on the space of paths with finite $`p`$-variation does not work. However, Lyons defined an integral against a continuous function of $`p`$-variation for some $`p>2`$ . The integral is developed in the space of geometric multiplicative functionals (described in Appendix A). The key idea is that we enhance the path by adding an area function to it. If there is sufficient control of the pair, path and area, then the integral is defined. The canonical example in is Brownian motion. The area process enhancing the Brownian motion is the Lévy area \[8, Ch.7, Sect.55\]. I show that there is an area process of a Lévy process which has finite $`(p/2)`$-variation a.s..
In order to solve (1) for a discontinuous function I add fictitious time during which linear segments remove the discontinuities, creating a continuous path. By solving for the continuous path and then removing the fictitious time we recover a solution for the discontinuous path. This is called a geometric solution. A second type of solution is derined from the geometric solution which we call the forward solution.
The first section treats the case where the discontinuous driving path has finite $`p`$-variation for some $`p<2`$. The second section treats the case where the path has finite $`p`$-variation for some $`p>2`$ only. The main proofs of the second section are deferred to the third section. In the appendix I prove the homeomorphic flow property for the solutions when the driving path is continuous. This is used in proving that forward solutions can be recovered from geometric solutions.
## 1 Discontinuous processes - $`p<2`$
In this section we extend the results of to allow the driving path of (1) to have discontinuities. The results are applied to sample paths of some Lévy processes, those that have finite $`p`$-variation a.s. for some $`p<2`$. Throughout this section $`p[1,2)`$ unless otherwise stated.
First, we determine the solution’s behaviour when the integrator jumps. There are two possibilities to consider: the first is an extension of the Lebesgue-Stieltjes integral; the second is based on a geometric approach.
Suppose that the discontinuous integrator has bounded variation. The solution $`y`$ would jump
$$y_ty_t=f(y_t)(x_tx_t)$$
at a jump time $`t`$ of $`x`$. If $`x`$ has finite $`p`$-variation for some $`1<p<2`$ we insert these jumps at the discontinuities of $`x`$. We call a path $`y`$ with the above jump behaviour a forward solution.
The other jump behaviour we consider is the following: When a jump of the integrator occurs we insert some fictitious time during which the jump is traversed by a linear segment, creating a continuous path on an extended time frame. Then we solve the differential equation driven by the continuous path. Finally we remove the fictitious time component of the solution path. We call this a geometric solution because the solution has an ’instantaneous flow’ along an integral curve at the jump times. This jump behaviour has been considered before by and .
The disadvantage of the first approach is that the solution does not, generally, generate a flow of diffeomorphisms .
In this section we prove the following theorem:
###### Theorem 1.1
Let $`x_t`$ be a discontinuous function of finite $`p`$-variation for some $`p<2`$. Let $`f`$ be an $`\alpha `$-Lipschitz vector field for some $`\alpha >p`$. Then there exists a unique geometric solution to the integral equation
$$y_t=y_0+_0^tf(y_t)𝑑x_ty_0=a^d.$$
(3)
With the above assumptions, there exists a unique forward solution as well.
Before proving the theorem we recall the definitions of $`p`$-variation and $`\alpha `$-Lipschitz:
###### Definition 1.1
The $`p`$-variation of a function $`x(s)`$ over the interval $`[0,t]`$ is defined as follows:
$$x_{_{p,[0,t]}}=\left\{\underset{\pi \pi [0,t]}{sup}\underset{\pi }{}|x(t_k)x(t_{k1})|^p\right\}^{\frac{1}{p}}$$
where $`\pi [0,t]`$ is the collection of all finite partitions of the interval $`[0,t]`$.
###### Remark 1.1
This is the strong $`p`$-variation. Usually probabilists use the weaker form where the supremum is over partitions restricted by a mesh size which tends to zero.
###### Definition 1.2
A function $`f`$ is in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >1`$ if
$$f_{\mathrm{}}<\mathrm{}\text{and}\frac{f}{x_j}\text{Lip}(\alpha 1)j=1,\mathrm{},d.$$
Its norm is given by
$$f_{\mathrm{Lip}(\alpha )}f_{\mathrm{}}+\underset{j=1}{\overset{d}{}}\frac{f}{x_j}_{\mathrm{Lip}(\alpha 1)}\text{for}\alpha >1.$$
This is Stein’s definition of $`\alpha `$-Lipschitz continuity for $`\alpha >1`$. It extends the classical definition: $`f`$ is in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha (0,1]`$ if
$$|f(x)f(y)|K|xy|^\alpha $$
with norm
$$f_{\mathrm{}}+\underset{xy}{sup}\frac{|f(x)f(y)|}{|xy|^\alpha }.$$
### 1.1 Geometric Solutions.
In this subsection we define a parametrisation for a càdlàg path $`x`$ of finite $`p`$-variation. The parametrisation adds fictitious time allowing the traversal of the discontinuities of the path $`x`$. We prove that the resulting continuous path $`x^\delta `$ has the same $`p`$-variation that $`x`$ has. We solve (3) driven by $`x^\delta `$ using the method of Lyons . Then we get a geometric solution of (3) by removing the fictitious time (i.e. by undoing the parametrisation).
###### Definition 1.3
Let $`x`$ be a càdlàg path of finite $`p`$-variation. Let $`\delta >0,`$ for each $`n1,`$ let $`t_n`$ be the time of the $`n`$’th largest jump of $`x`$. We define a map $`\tau ^\delta :[0,T][0,T+\delta _{i=1}^{\mathrm{}}|j(t_i)|^p]`$ (where $`j(u)`$ denotes the jump of the path $`x`$ at time u) in the following way:
$$\tau ^\delta (t)=t+\delta \underset{n=1}{\overset{\mathrm{}}{}}|j(t_n)|^p\chi _{\{t_nt\}}(t).$$
(4)
The map $`\tau ^\delta :[0,T][0,\tau ^\delta (T)]`$ extends the time interval into one on where we define the continuous process $`x^\delta (s)`$
$`x^\delta (s)`$
$`=`$ $`\{\begin{array}{cc}x(t)\hfill & \text{if}s=\tau ^\delta (t),\hfill \\ x(t_n^{})+(s\tau ^\delta (t_n^{}))j(t_n)\delta ^1|j(t_n)|^p\hfill & \text{if}s[\tau ^\delta (t_n^{})\tau ^\delta (t_n)).\hfill \end{array}`$ (7)
###### Remarks 1.1
1. $`(s,x_s^\delta ),s[0,\tau ^\delta (T)]`$ is a parametrisation of the driving path $`x`$.
2. The terms $`|j(t_n)|^p`$ in (4) ensure that the addition of the fictitious time does not make $`\tau ^\delta (t)`$ explode.
3. In Figure 1 we see an example of a parametrisation of a discontinuous path $`x_s`$ in terms of the pair $`(t(s),y(s))`$.
The next proposition shows that the above parametrisation has the same $`p`$-variation as the original path, on the extended time frame $`[0,\tau ^\delta (T)]`$.
###### Proposition 1.1
Let $`x`$ be a càdlàg path of finite $`p`$-variation. Let $`x^\delta `$ be a parametrisation of $`x`$ as above. Then
$$x^\delta _{p,[0,\tau ^\delta (T)]}=x_{p,[0,T]}\delta >0.$$
Proof. Let $`\pi _0`$ be a partition of $`[0,\tau ^\delta (T)]`$. Let
$$V_{x^\delta }(\pi _0)=\underset{\pi _0}{}|x^\delta (t_i)x^\delta (t_{i1})|^p$$
We show that we increase the value of $`V_p(\pi _0)`$ by moving points lying on the jump segments to the endpoints of those segments.
Let $`t_{i1},t_i,t_{i+1}`$ be three neighbouring points in the partition $`\pi _0`$ such that $`t_i`$ lies in a jump segment. Consider the following term:
$$|x_{t_i}^\delta x_{t_{i1}}^\delta |^p+|x_{t_{i+1}}^\delta x_{t_i}^\delta |^p.$$
(8)
We show that (8) is dominated by replacing $`x_{t_i}^\delta `$ by one of $`x_l^\delta `$ and $`x_r^\delta `$, where $`l`$ and $`r`$ denote the left and right endpoint of the jump segment containing $`t_i`$.
For simplicity we set $`a=x_{t_{i1}}^\delta ,b=x_{t_{i+1}}^\delta `$ and $`c=x_l^\delta `$. Let
$$L\{c+kx:k(0,1),c,x^d,x0\},a,b^d\backslash L.$$
Let the function $`f:[0,1](0,\mathrm{})`$ be defined by
$$f(k)=|ad|^p+|db|^p,d=c+kx.$$
Then $`fC^2[0,1]`$ and one can show that $`f^{\prime \prime }0`$ on $`(0,1)`$ when $`p1`$. To conclude the proof we move along the partition replacing $`t_i`$ which lie in the jump segments by new points $`t_i^{}`$ that increase $`V_{x^\delta }(\pi _0)`$. The partition $`\pi _0`$ is replaced by a partition $`\pi _0^{}`$ whose points lie on the pre-image of $`[0,\tau ^\delta (T)]`$. Therefore we have
$$V_{x^\delta }(\pi _0)V_{x^\delta }(\pi _0^{})=V_x(\pi _0^{}).$$
Hence $`x^\delta _{p,[0,\tau ^\delta (T)]}=x_{p,[0,T]}`$. $`\mathrm{}`$
###### Theorem 1.2
Let $`x`$ be a càdlàg path with finite $`p`$-variation for some $`p<2`$. Let $`f`$ be a $`\mathrm{Lip}(\gamma )`$ vector field on $`^n`$ for some $`\gamma >p`$. Then there exists a unique geometric solution $`y`$, having finite $`p`$-variation which solves the differential equation
$$dy_t=f(y_t)dx_ty_0=a^n.$$
(9)
Proof. Let $`x^\delta `$ be the parametrisation given in (1.3). The theorem of section three of proves that there is a continuous solution $`y^\delta `$ which solves (3) on $`[0,\tau ^\delta (T)]`$. Then $`(s,y_s^\delta )`$ is a parametrisation of a càdlàg path $`y`$ on $`[0,T]`$.
The solution is well-defined. To see this, consider two parametrisations of $`x`$ and note that there exists a monotonically increasing function $`\lambda _s`$ such that
$$(s,x_s^\delta )=(\lambda _s,x_{\lambda _s}^\nu ).\mathrm{}$$
### 1.2 Forward Solutions.
In this subsection we show how to recover forward solutions from geometric solutions. The idea behind our approach is to correct the jump behaviour of the geometric solution using a Taylor series expansion Lemma 1.1. The correction terms are controlled by
$$\underset{i=1}{\overset{\mathrm{}}{}}|x_{t_i}x_{t_i^{}}|^2$$
which is finite due to the finite $`p`$-variation of the path $`x`$.
In the case where the driving path has only a finite number of jumps we note that the forward solution can be recovered trivially. It is enough to mark the jump times of $`x`$ and solve the differential equation on the components where $`x`$ is continuous, inserting the forward jump behaviour when the jumps occur. It remains to show that the forward solution exists when the driving path has a countably infinite number of jumps. The method we use requires the following property of the geometric solution:
###### Theorem 1.3
Let $`x`$ be a continuous path of finite $`p`$-variation for some $`p>1`$. Let $`f`$ be in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p`$. The maps $`(\pi _t)_{t0}:^n^n`$ obtained by varying the initial condition of the following differential equation generate a flow of homeomorphisms:
$$d\pi _t=f(\pi _t)dx_t\pi _0=Id,(\text{the identity map}).$$
(10)
We leave the proof of Theorem 1.3 until Appendix A. We note the uniform estimate
$$\underset{0tT}{sup}\left|\pi _t^a\pi _t^b\right|C(T)\left|ab\right|.$$
(11)
The following lemma will enable estimates to be made when the geometric jumps are replaced by the forward jumps:
###### Lemma 1.1
Let $`x`$ be a càdlàg path with finite $`p`$-variation. Let $`f`$ be in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p`$. Let $`\mathrm{\Delta }y_i`$ (resp. $`\mathrm{\Delta }z_i`$) denote the geometric (resp. forward) solution’s jump which correspond to $`\mathrm{\Delta }x_i`$, the $`i`$’th largest jump of $`x`$. Then we have the following estimate on the difference of the two jumps:
$$\mathrm{\Delta }y_i\mathrm{\Delta }z_i_{\mathrm{}}K|\mathrm{\Delta }x_i|^2$$
where the constant $`K`$ depends on $`f_{\mathrm{Lip}(\alpha )}`$.
Proof. Parametrise the path $`x`$ so that it traverses its discontinuity in unit time. Solve geometrically over this interval with the solution having initial point $`a`$. Note that the forward jump is the first order Taylor approximation to the geometric jump. Then
$`y_1(a)=y_0(a)`$ $`+`$ $`{\displaystyle \frac{dy_s(a)}{ds}}|_{s=0}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{d{}_{}{}^{2}y_{s}^{}(a)}{ds^2}}|_{s=\theta }\text{for some}\mathrm{\hspace{0.33em}0}<\theta <1`$ (12)
$`=`$ $`z_1(a)+{\displaystyle \frac{1}{2}}{\displaystyle \frac{d{}_{}{}^{2}y_{s}^{}(a)}{ds^2}}|_{s=\theta }.`$
We estimate the second order term by
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d{}_{}{}^{2}y_{s}^{}(a)}{ds^2}}_{\mathrm{}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d}{ds}}f(y_s(a))(\mathrm{\Delta }x_i)_{\mathrm{}}`$ (13)
$``$ $`{\displaystyle \frac{1}{2}}f_{\mathrm{}}f_{\mathrm{}}|\mathrm{\Delta }x_i|^2`$
$``$ $`{\displaystyle \frac{1}{2}}f_{\mathrm{Lip}(\alpha )}^2|\mathrm{\Delta }x_i|^2`$
Both $`f_{\mathrm{}}`$ and $`f_{\mathrm{}}`$ are finite because $`f`$ is $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p1`$. $`\mathrm{}`$
###### Theorem 1.4
Let $`x`$ be a càdlàg path with finite $`p`$-variation. Let $`f`$ be in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p`$. Then there exists a unique forward solution to the following differential equation:
$$dz_t=f(z_t)dx_tz_0=a.$$
(14)
Proof. By Theorem 1.3 there exists a unique homeomorphism $`y`$ which solves
$$dy_t=f(y_t)dx_ty_0=a$$
in a geometric sense.
Label the jumps of $`x`$ by $`j_x=\{j_i\}_{i=1}^{\mathrm{}}`$ according to their decreasing size. Let $`z^n`$ denote the path made by replacing the geometric jumps of $`y`$ corresponding to $`\{j_i\}_{i=1}^n`$ by the forward jumps $`\{f()(\mathrm{\Delta }x_i)\}_{i=1}^n`$. We show that the $`(z^n)_{n1}`$ have a uniform limit.
We order the corrected jumps chronologically, say $`\{t_i\}_{i=1}^n`$. Then we estimate the following term using Lemma 1.1 and the uniform bound on the growth of $`y`$ given in (11):
$`|z_s^n(a)y_s(a)|`$ $``$ $`{\displaystyle \underset{i=1}{\overset{n}{}}}|y_{t_i,s}(z_{t_i}^n(a))y_{t_i,s}(y_{t_{i1},t_i}(z_{t_{i1}}^n(a)))|`$ (15)
$``$ $`C(T){\displaystyle \underset{i=1}{\overset{n}{}}}|z_{t_i}^n(a)y_{t_{i1},t_i}(z_{t_{i1}}^n(a))|`$
$``$ $`C^2(T)K{\displaystyle \underset{i=1}{\overset{\mathrm{}}{}}}|\mathrm{\Delta }x_i|^2.`$
So we have the uniform estimate
$$z^ny_{\mathrm{}}K(C_3(T),f_{\mathrm{Lip}(\alpha )})\underset{i=1}{\overset{\mathrm{}}{}}|\mathrm{\Delta }x_i|^2<\mathrm{}n1.$$
(16)
We use an analogous bound to get Cauchy convergence of $`\{z^n\}_{n1}`$. Let $`m,r1`$.
$$z^mz^{m+r}_{\mathrm{}}K(C(T,z^m),f_{\mathrm{Lip}(\alpha )})\underset{i=m+1}{\overset{\mathrm{}}{}}|\mathrm{\Delta }x_i|^2.$$
One notes that $`\{C(T,z^m)\}`$ are uniformly bounded, because of the boundedness of $`C(T)=C(T,y)`$ and the Lipschitz condition on $`f`$. Therefore we have the following estimate:
$$z^mz^{m+r}_{\mathrm{}}L\underset{i=m+1}{\overset{\mathrm{}}{}}|\mathrm{\Delta }x_i|^2.$$
This implies that $`\{z^n\}`$ are Cauchy in the supremum norm because $`x`$ has finite $`p`$-variation $`(p<2)`$ which implies that $`_{m+1}^{\mathrm{}}|\mathrm{\Delta }x_i|^2`$ tends to zero as $`m`$ increases. $`\mathrm{}`$
###### Remark 1.2
Theorems 1.4 and 1.2 combine to prove Theorem 1.1.
###### Corollary 1.1
With the above notation, $`z`$ has finite $`p`$-variation.
Proof. Let $`s<t[0,T]`$.
$`|z_tz_s|`$ $`|(z_tz_s)(y_ty_s)|+|y_ty_s|`$
where $`(y_ty_s)`$ is the increment of the geometric solution starting from $`z_s`$ driven by the path $`x_t`$ on the interval $`[s,T]`$. Then
$`|(z_tz_s)(y_ty_s)|`$ $`C{\displaystyle \underset{\begin{array}{c}j_x|_{[s,t]}\end{array}}{}}|\mathrm{\Delta }x_i|^2\text{ and }|y_ty_s|x_{p,[s,t]},`$
which implies that
$`|z_tz_s|^p`$ $`2^{p1}\left\{C^p\left({\displaystyle \underset{\begin{array}{c}j_x|_{[s,t]}\end{array}}{}}|\mathrm{\Delta }x_i|^2\right)^p+x_{p,[s,t]}^p\right\},`$
hence
$`z_{p,[0,T]}`$ $`2^{(p1)/p}\left\{C^p\left({\displaystyle \underset{\begin{array}{c}j_x|_{[0,T]}\end{array}}{}}|\mathrm{\Delta }x_i|^2\right)^p+x_{p,[0,T]}^p\right\}^{1/p}<\mathrm{}.`$
$`\mathrm{}`$
### 1.3 $`p`$-variation of Lévy processes
In this subsection we apply Theorem 1.1 to Lévy processes which have finite $`p`$-variation a.s..
Lévy processes are the class of processes with stationary, independent increments which are continuous in probability. The class includes Brownian motion, although this process is atypical due to its continuous sample paths. Typically a Lévy process will be a combination of a deterministic drift, a Gaussian process and a jump process. For further information on Lévy processes we direct the reader to .
The regularity of the sample paths of a Lévy process has been studied intensively. In the 1960’s several people worked on the question of characterising the sample path $`p`$-variation. The following theorem, due to Monroe, gives the characterisation:
###### Theorem 1.5
\[15, Theorem 2\] Let $`(X_t)_{t0}`$ be a Lévy process in $`^n`$ without a Gaussian part. Let $`\nu `$ be the Lévy measure. Let $`\beta `$ denote the index of $`X_t`$, that is
$$\beta inf\{\alpha >0:_{|y|1}|y|^\alpha \nu (dy)<\mathrm{}\}$$
(17)
and suppose that $`1\beta 2`$. If $`\gamma >\beta `$ then
$$\left(X_\gamma <\mathrm{}\right)=1$$
(18)
where the $`\gamma `$-variation is considered over any compact interval.
###### Remark 1.3
Note that all Lévy processes with a Gaussian part only have finite $`p`$-variation for $`p>2`$.
###### Corollary 1.2
Let $`(X_t)_{t0}`$ be a Lévy process with index $`\beta <2`$ and no Gaussian part. Let $`f`$ be a vector field in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p`$. Then, a.s., the following stochastic differential equation has a unique forward and a unique geometric solution:
$$dY_t=f(Y_t)dX_tY_0=a.$$
Proof. The corollary follows immediately from Theorems 1.5 and 1.1. $`\mathrm{}`$
## 2 Discontinuous processes - $`p>2`$
The goal of this section is to extend (Corollary 1.2) to let any Lévy process be the integrator of (1).
One problem we have is that the Young integral is no longer useful because we use a Picard iteration scheme which fails condition (2) when $`p>2`$. However, we can use the method from . To define the integral we need to provide more information about the sample path. We do this by defining an area process of the Lévy process. Then we prove that the enhanced process (path and area) has finite $`p`$-variation Definition A.3.
We parametrise the enhanced process in an analogous manner to (1.3) (adding fictitious time). Then we solve (1) in a geometric sense using the method for continuous paths $`(p>2)`$ given in . Finally, forward solutions are obtained by jump correction as before.
Before enhancing $`(X_t)_{t0}`$ we give an example which shows that there exist Lévy measures with index two. So a Lévy process does not need a Gaussian part to have, a.s., finite $`p`$-variation only for $`p>2`$.
###### Example 2.1
One can define the following measures on $``$:
$`\nu _k(dx)`$ $`|x|^{3+1/k}dx|x|((k+1)^{3(k+1)},k^{3k}]J_k`$
$`\eta _m(dx)`$ $`{\displaystyle \underset{k=1}{\overset{m}{}}}\nu _k(dxJ_k(J_k)).`$
We show that $`\eta lim_m\mathrm{}\eta _m`$ is a Lévy measure. The integrability condition
$$_{|x|1}|x|^2\eta (dx)<\mathrm{}.$$
(19)
must be satisfied.
$`{\displaystyle _{|x|1}}|x|^2\eta _m(dx)`$ $`=2{\displaystyle _0^1}{\displaystyle \underset{k=1}{\overset{m}{}}}x^{1+1/k}\chi _{J_k}(x)dx=2{\displaystyle \underset{k=1}{\overset{m}{}}}\left[kx^{1/k}\right]_{(k+1)^{3(k+1)}}^{k^{3k}}`$
$`=2{\displaystyle \underset{k=1}{\overset{m}{}}}k\left\{k^3(k+1)^{3(1+1/k)}\right\}`$
$`2{\displaystyle \underset{k=1}{\overset{m}{}}}k\left\{k^32^{3(1+1/k)}k^{3(1+1/k)}\right\}`$
$`=2{\displaystyle \underset{k=1}{\overset{m}{}}}k^2\left\{12^{3(1+1/k)}k^{3/k}\right\}`$
$`<C{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}k^2<\mathrm{},`$
where $`C`$ is some suitable constant. We take the limit as $`m`$ tends to infinity on the left hand side to prove (19).
Now we show that
$$_{|x|1}|x|^\alpha \eta (dx)=\mathrm{}$$
(20)
for all $`\alpha <2`$. Fix $`\alpha <2`$. Define the following number:
$`m(\alpha )`$ $`inf\{k:\alpha +1/k<2\}<\mathrm{}\text{as}\alpha <2.`$
Let $`m>m(\alpha )`$. Then
$`{\displaystyle _{|x|1}}|x|^\alpha \eta _m(dx)`$
$`2{\displaystyle \underset{k=m(\alpha )}{\overset{m}{}}}{\displaystyle \frac{1}{(\alpha +1/k2)}}\left\{k^{3k(\alpha +1/k2)}(k+1)^{3(k+1)(\alpha +1/k2)}\right\}`$
$`=2{\displaystyle \underset{k=m(\alpha )}{\overset{m}{}}}{\displaystyle \frac{1}{(2(\alpha +1/k))}}\left\{(k+1)^{3(k+1)(\alpha +1/k2)}k^{3k(\alpha +1/k2)}\right\}`$
$`{\displaystyle \frac{2}{2\alpha }}{\displaystyle \underset{m(\alpha )}{\overset{m}{}}}\left\{(k+1)^{3(k+1)(2(\alpha +1/k))}k^{3k(2(\alpha +1/k))}\right\}`$
$`\mathrm{}\text{as}m\mathrm{}.`$
This proves that the index $`\beta `$ of $`\eta `$ equals two. Theorem 1.5 implies that the pure jump process associated to the Lévy measure $`\eta `$ a.s. has finite $`p`$-variation for $`p>2`$ only.
The following theorem gives a construction of the Lévy area of the Lévy process $`(X_t)_{t0}`$. The Lévy area process and the Lévy process form the enhanced process which we need in order to use the method of Lyons .
###### Theorem 2.1
The $`d`$-dimensional Lévy process $`(X_t)_{t0}`$ has an anti-symmetric area process
$$(A_{s,t})^{ij}\frac{1}{2}_s^tX_u^idX_u^jX_u^jdX_u^ii,j=1,2.\mathrm{a}.\mathrm{s}.$$
The proof is deferred to Section 3.
###### Theorem 2.2
The Lévy area of the Lévy process $`(X_t)_{t0}`$ a.s. has finite $`(p/2)`$-variation for $`p>2`$. That is
$$\underset{\pi }{sup}\left(\underset{\pi }{}|A_{t_{k1},t_k}|^{p/2}\right)^{2/p}<\mathrm{}\mathrm{a}.\mathrm{s}.$$
where the supremum is taken over all finite partitions $`\pi `$ of $`[0,T]`$.
The proof is deferred to Section 3.
Now we parametrise the sample paths of $`(X_t)_{t0}`$ as before (1.3).
###### Proposition 2.1
Parametrising the process $`(X_t)_{t0}`$ does not affect the area process’ $`(p/2)`$-variation.
Proof. The proof is similar to the proof of Proposition 1.1. One can show that if $`\lambda `$ lies in a jump segment then
$$|A_{s,\lambda }|^{(p/2)}+|A_{\lambda ,t}|^{(p/2)}s<\lambda <t$$
is maximised when $`\lambda `$ is moved to one of the endpoints of the jump segment. $`\mathrm{}`$
With the parametrisation of the path and the area we can define the integral in the sense of Lyons . Consequently we have the following theorem:
###### Theorem 2.3
Let $`(X_t)_{t0}`$ be a Lévy process with finite $`p`$-variation for some $`p>2`$. Let $`f`$ be in $`\mathrm{Lip}(\alpha )`$ for some $`\alpha >p`$. Then there exists, with probability one, a unique geometric and a unique forward solution to the following integral equation:
$$Y_t=Y_0+_0^tf(Y_t)𝑑X_tY_0=a^d.$$
(21)
###### Remark 2.1
When constructing the forward solution it is necessary that the sum
$$\underset{n=1}{\overset{\mathrm{}}{}}|\mathrm{\Delta }X_n|^2$$
remains finite. This is guaranteed by the requirement on Lévy measures to satisfy
$$_{|x|1}|x|^21\nu (dx)<\mathrm{}.$$
## 3 Proofs of Theorem 2.1 and Theorem 2.2
For clarity throughout this section we assume that the Lévy process $`(X_t)_{t0}`$ is two dimensional and takes the following form:
$$X_t=B_t+_{|x|1}x(N_t(dx)t\nu (dx)).$$
(22)
That is, $`(X_t)_{t0}`$ is a Gaussian process with a compensated pure jump process, whose Lévy measure is supported on $`(x^2:|x|1)`$.
###### Proposition 3.1
The $`d`$-dimensional Lévy process $`(X_t)_{t0}`$ has an anti-symmetric area process
$$(A_{s,t})^{ij}\frac{1}{2}_s^tX_u^idX_u^jX_u^jdX_u^ii,j=1,2.\mathrm{a}.\mathrm{s}.$$
For fixed $`s<t`$ we obtain the area process by the following limiting procedure:
$$(A_{s,t})^{ij}=\underset{n\mathrm{}}{lim}\underset{m=0}{\overset{n}{}}\underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}A_{k,m}^{i,j}\mathrm{a}.\mathrm{s}.$$
where $`A_{k,m}^{ij}`$ is the area of the $`(ij)`$-projected triangle with vertices
$$X(u_{(k+1)/2,m1}),X(u_{(k1)/2,m1}),X(u_{k,m})$$
where $`u_{k,m}s+k2^m(ts)`$. Also we have the second order moment estimate
$$𝔼\left((A_{s,t}^{ij})^2\right)C(\nu )(ts)^2.$$
(23)
Proof. We define $`A_{s,t}(n)`$
$`A_{s,t}(n)`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{k=0}{\overset{2^n1}{}}}(X^{(1)}(u_{k,n})X^{(1)}(s))(X^{(2)}(u_{k+1,n})X^{(2)}(u_{k,n}))`$
$`(X^{(2)}(u_{k,n})X^{(2)}(s))(X^{(1)}(u_{k+1,n})X^{(1)}(u_{k,n}))`$
$`={\displaystyle \underset{k=0}{\overset{2^n1}{}}}B_{k,n}`$
where $`B_{k,n}`$ is the (signed) area of the triangle with vertices
$$X(s),X(u_{k,n}),X(u_{k+1,n}).$$
By considering the difference between $`A_{s,t}(n)`$ and $`A_{s,t}(n+1)`$ we see that
$$B_{2k,n+1}+B_{2k+1,n+1}B_{k,n}$$
is the area of the triangle with vertices
$$X(u_{k,n}),X(u_{k+1,n}),X(u_{2k+1,n+1})$$
which we denote by $`A_{k,n}`$. We re-order $`A_{s,t}(n)`$
$`A_{s,t}(n)`$ $`={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}\left(X(u_{k,m})d_{k,m}\right)`$
$`\left(X(u_{(k+1)/2,m1})X(u_{(k1)/2,m1})\right)`$
$`={\displaystyle \frac{1}{2}}{\displaystyle \underset{m=0}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}A_{k,m}`$
where $`d_{k,m}1/2(X(u_{(k+1)/2,m1})+X(u_{(k1)/2,m1}))`$. The convergence to the area process is completed using martingale methods.
Let $`𝔉_n\sigma (X(u_{k,n}):k=0,\mathrm{},2^n).`$ Then
###### Lemma 3.1
$$𝔼\left(X(u_{k,m})|𝔉_{m1}\right)=d_{k,m}\mathrm{a}.\mathrm{s}.$$
(24)
Proof. For ease of presentation we let
$`U_1`$ $`X(u_{k,m})X(u_{(k1)/2,m1})`$
$`U_2`$ $`X((u_{(k+1)/2,m1})X(u_{k,m}).`$
Then
$`𝔼\left(X(u_{k,m})d_{k,m}|𝔉_{m1}\right)`$
$`=𝔼\left(X(u_{k,m})d_{k,m}|X(u_{(k1)/2,m1}),X(u_{(k+1)/2,m1})\right)`$
$`={\displaystyle \frac{1}{2}}𝔼\left(U_1U_2|X(u_{(k1)/2,m1}),X(u_{(k+1)/2,m1})\right)`$
Using the stationarity and the independence of the increments of $`X`$ we see that $`U_1`$ and $`U_2`$ are exchangeable, that is
$$\left(U_1A,U_2B\right)=\left(U_2A,U_1B\right)A,B𝔅(^2).$$
The exchangeability extends to the random variables
$$\left(U_i|X(u_{(k1)/2,m1}),X(u_{(k+1)/2,m1})\right)i=1,2.$$
We deduce that
$$𝔼(U_1U_2|X(u_{(k1)/2,m1}),X(u_{(k+1)/2,m1}))=0.\mathrm{}$$
Returning to the proof of Proposition 3.1, we compute the variance of $`A_{k,m}`$. This will be used to show that
$$\underset{n1}{sup}𝔼\left(A_{s,t}(n)^2\right)<\mathrm{}.$$
$`𝔼(`$ $`A_{k,m}^2)`$
$`=𝔼(`$ $`([X^{(1)}(u_{k,m})d_{k,m}^{(1)}][U_1^{(2)}+U_2^{(2)}][X^{(2)}(u_{k,m})d_{k,m}^{(2)}][U_1^{(1)}+U_2^{(1)}])^2)`$
$`={\displaystyle \frac{1}{4}}𝔼\left(\{(U_1^{(1)}U_2^{(1)})(U_1^{(2)}+U_2^{(2)})(U_1^{(2)}U_2^{(2)})(U_1^{(1)}+U_2^{(1)})\}^2\right)`$
$`={\displaystyle \frac{1}{4}}𝔼\left((U_1^{(1)}U_2^{(2)})^22U_1^{(1)}U_2^{(2)}U_2^{(1)}U_1^{(2)}+(U_2^{(1)}U_1^{(2)})^2\right)`$
$`(1)+(2)+(3).`$
We use the independence of the increments and Itô’s formula for discontinuous semi-martingales to compute $`(1),(2)`$ and $`(3)`$.
$`(1)`$ $`=𝔼\left((U_1^{(1)}U_2^{(2)})^2\right)=𝔼\left((U_1^{(1)})^2\right)𝔼\left((U_2^{(2)})^2\right).`$
By applying Itô’s formula and using the stationarity of the Lévy process we find that
$`(3)=(1)=2^{2m}(ts)^2{\displaystyle _{|x|1}}|x_1|^2\nu (dx){\displaystyle _{|x|1}}|x_2|^2\nu (dx).`$
Another application of Itô’s formula gives
$`(2)`$ $`=2𝔼\left(U_1^{(1)}U_2^{(2)}U_2^{(1)}U_1^{(2)}\right)=2𝔼\left(U_1^{(1)}U_1^{(2)}\right)𝔼\left(U_2^{(2)}U_2^{(1)}\right)`$
$`=2^{2m+1}(ts)^2\left({\displaystyle _{|x|1}}x_1x_2\nu (dx)\right)^2.`$
Collecting the terms together we have the following expression:
$`𝔼\left(A_{k,m}^2\right)=C_0(\nu )\mathrm{\hspace{0.33em}2}^{2m+1}(ts)^2`$
where
$`C_0(\nu )\left\{{\displaystyle _{|x|1}}|x_1|^2\nu (dx){\displaystyle _{|x|1}}|x_2|^2\nu (dx)\left({\displaystyle _{|x|1}}x_1x_2\nu (dx)\right)^2\right\}.`$
Now we estimate the following term:
$`𝔼\left(A_{s,t}^2(n)\right)`$ $`=𝔼\left(\left({\displaystyle \underset{m=1}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}A_{k,m}\right)^2\right)`$
which through conditioning and independence arguments equals
$`=𝔼\left({\displaystyle \underset{m=1}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}A_{k,m}^2\right)`$
$`=C_0(\nu ){\displaystyle \underset{m=1}{\overset{n}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}2^{2m+1}(ts)^2`$
$`C_0(\nu ){\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}{\displaystyle \underset{\begin{array}{c}k=1,\\ odd\end{array}}{\overset{2^m1}{}}}2^{2m+1}(ts)^2C(\nu )(ts)^2.`$
We use the martingale convergence theorem to deduce that a.s. there is a unique limit of $`A_{s,t}(n)`$. Furthermore the last calculation implies that there is a moment estimate of the area process given by
$$𝔼\left(A_{s,t}^2\right)C(\nu )(ts)^2.\mathrm{}$$
We note that there is another way that one could define an area process of a Lévy process. One could define the area process for the truncated Lévy processes and look for a limit as the small (compensated) jumps are put in. Using the above construction one can define $`A_{s,t}^ϵ`$ for a fixed pair of times, corresponding to the Lévy process $`X^ϵ`$. With the $`\sigma `$-fields $`(𝔊^ϵ)_{ϵ>0}`$ defined by
$$𝔊^ϵ\sigma (X^\delta :\delta >ϵ)\text{for}ϵ>0$$
we have the following proposition:
###### Proposition 3.2
$`(A_{s,t}^ϵ)_{ϵ>0}`$ form a $`(𝔊^ϵ)`$-martingale.
Proof. Let $`\eta >ϵ>0`$. By considering the construction of the area given above for the truncated processes $`X^\eta `$ and $`X^ϵ`$ we look at the difference at the level of the triangles $`A_{k,n}^\eta `$ and $`A_{k,n}^ϵ`$.
$`𝔼\left(A_{k,n}^ϵA_{k,n}^\eta |𝔊^\eta \right)`$ $`=𝔼(A_{k,n}^{\eta ,ϵ}`$
$`+(X_{k,n}^{\eta ,ϵ}d_{k,n}^{\eta ,ϵ})(X_{(k+1)/2,n1}^\eta X_{(k1)/2,n1}^\eta )`$
$`+(X_{k,n}^\eta d_{k,n}^\eta )(X_{(k+1)/2,n1}^{\eta ,ϵ}X_{(k1)/2,n1}^{\eta ,ϵ})\left|𝔊^\eta \right)`$
where the superscript $`\eta ,ϵ`$ signifies that the process is generated by the part of the Lévy measure whose support is $`(ϵ,\eta ]`$. Using the spatial independence of the underlying Lévy process we have
$`=𝔼\left(A_{k,n}^{\eta ,ϵ}\right)+𝔼\left((X_{k,n}^{\eta ,ϵ}d_{k,n}^{\eta ,ϵ})\right)(X_{(k+1)/2,n1}^\eta X_{(k1)/2,n1}^\eta )`$
$`+(X_{k,n}^\eta d_{k,n}^\eta )𝔼\left((X_{(k+1)/2,n1}^{\eta ,ϵ}X_{(k1)/2,n1}^{\eta ,ϵ})\right)`$
$`=0.\mathrm{}`$
With the uniform control on the second moment of the martingale
$$𝔼\left((A_{s,t}^ϵ)^2\right)C(\nu )(ts)^2ϵ>0$$
we conclude that $`A_{s,t}^ϵ`$ converges a.s. as $`ϵ0`$.
The algebraic identity
$$A_{s,u}=A_{s,t}+A_{t,u}+\frac{1}{2}[X_{s,t},X_{t,u}]s<t<u$$
(25)
for the anti-symmetric area process $`A`$ generated by a piecewise smooth path $`X`$ extends to the area process of the Lévy process. This is due to (25) holding for the area processes $`A^ϵ`$ of the truncated Lévy processes $`X^ϵ`$.
###### Proposition 3.3
The Lévy area of the Lévy process $`(X_t)_{t0}`$ has finite $`(p/2)`$-variation for $`p>2`$ a.s.. That is
$$\underset{\pi }{sup}\left(\underset{\pi }{}|A_{t_{k1},t_k}|^{p/2}\right)^{2/p}<\mathrm{}\mathrm{a}.\mathrm{s}.$$
where the supremum is taken over all finite partitions $`\pi `$ of $`[0,T]`$.
Proof. In Proposition 3.1 we constructed the area process for a pair of times, a.s.. This can be extended to a countable collection of pairs of times, a.s.. In the proof below we assume that the area process has been defined for the times
$$k2^nT,(k+1)2^nTk=0,1,\mathrm{},2^n1,n1.$$
The proof follows the method of estimation used in . To estimate the area process for two arbitrary times $`u<v`$ we split up the interval $`[u,v]`$ in the following manner:
We select the largest dyadic interval $`[(k1)2^nT,k2^n]`$ which is contained within $`[u,v]`$. Then we add dyadic intervals to either side of the initial interval, which are chosen maximally with respect to inclusion in the interval $`[u,v]`$. Continuing in this fashion we label the partition according to the lengths of the dyadics. We note that there are at most two dyadics of the same length in the partition which we label $`[l_{1,k},r_{1,k}]`$ and $`[l_{2,k},r_{2,k}]`$ where $`r_{1,k}l_{2,k}`$. Then
$$[u,v]=\underset{k=1}{\overset{\mathrm{}}{}}\underset{i=1,2}{}[l_{i,k},r_{i,k}].$$
We estimate $`A_{u,v}`$ using the algebraic formula (25).
$`A_{l_{1,m},r_{2,m}}`$ $`={\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{i=1,2}{}}A_{l_{i,k},r_{i,k}}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \underset{1ab2}{}}{\displaystyle \underset{1j<km}{}}[X_{r_{a,k}}X_{l_{a,k}},X_{r_{b,j}}X_{l_{b,j}}].`$
Noting that
$`{\displaystyle \underset{1ab2}{}}{\displaystyle \underset{1j<km}{}}\left|[X_{r_{a,k}}X_{l_{a,k}},X_{r_{b,j}}X_{l_{b,j}}]\right|`$
$`={\displaystyle \underset{1ab2}{}}{\displaystyle \underset{1j<km}{}}|(X_{r_{a,k}}X_{l_{a,k}})(X_{r_{b,j}}X_{l_{b,j}})`$
$`(X_{r_{b,j}}X_{l_{b,j}})(X_{r_{a,k}}X_{l_{a,k}})|`$
$`{\displaystyle \underset{1ab2}{}}{\displaystyle \underset{1j<km}{}}\left|X_{r_{a,k}}X_{l_{a,k}}\right|\left|X_{r_{b,j}}X_{l_{b,j}}\right|`$
$`\left({\displaystyle \underset{k=1}{\overset{m}{}}}{\displaystyle \underset{i=1,2}{}}\left|X_{r_{i,k}}X_{l_{i,k}}\right|\right)^2`$
we have the estimate:
$$|A_{u,v}|^{p/2}2^{(p/2)1}\left[\left(\underset{k=1}{\overset{\mathrm{}}{}}\underset{i=1,2}{}|A_{l_{i,k},r_{i,k}}|\right)^{p/2}+\frac{1}{2}\left(\underset{k=1}{\overset{\mathrm{}}{}}\underset{i=1,2}{}|X_{r_{i,k}}X_{l_{i,k}}|\right)^p\right].$$
(26)
Using Hölder’s inequality, with $`p>2`$ and $`\gamma >p1`$, we have
$`|A_{u,v}|^{p/2}`$ $`2^{(p/2)1}[\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{\gamma /((p/2)1)}\right)^{(p/2)1}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma ({\displaystyle \underset{i=1,2}{}}|A_{l_{i,k},r_{i,k}}|)^{p/2}`$
$`+{\displaystyle \frac{1}{2}}\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^{\gamma /(p1)}\right)^{p1}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma ({\displaystyle \underset{i=1,2}{}}|X_{r_{i,k}}X_{l_{i,k}}|)^p]`$
$`C_1(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1,2}{}}|A_{l_{i,k},r_{i,k}}|^{p/2}`$
$`+C_2(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1,2}{}}|X_{r_{i,k}}X_{l_{i,k}}|^p.`$ (27)
One can uniformly bound $`|A_{u,v}|^{p/2}`$ for any pair of times $`u<v[0,T]`$ by extending the estimate in (27) over all the dyadic intervals at each level $`n`$, that is,
$`|A_{u,v}|^{p/2}`$ $`C_1(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}|A_{l_{i,k},r_{i,k}}|^{p/2}`$
$`+C_2(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}|X_{r_{i,k}}X_{l_{i,k}}|^p.`$
If the right hand side is finite a.s. then the area can be defined for any pair of times.
The $`(p/2)`$-variation of the Lévy area can be estimated by the same bound.
$`\underset{\pi }{sup}{\displaystyle \underset{\pi }{}}|A_{u,v}|^{p/2}`$ $`C_1(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}|A_{l_{i,k},r_{i,k}}|^{p/2}`$
$`+C_2(p,\gamma ){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}|X_{r_{i,k}}X_{l_{i,k}}|^p.`$ (28)
We use (23) to control the first sum
$$𝔼\left(|A_{s,t}|^{p/2}\right)C(ts)^{p/2}\text{for}p4.$$
So we have
$`𝔼\left({\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}|A_{l_{i,k},r_{i,k}}|^{p/2}\right)`$ $`C{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{i=1}{\overset{2^n}{}}}(2^nT)^{p/2}`$
$`=C{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma 2^{n((p/2)1)}`$
$`<\mathrm{}\text{for}p>2.`$
This implies that the first term in the right hand side of (3) is a.s. finite. Now we consider the second term of (3).
###### Lemma 3.2
$$\underset{n=1}{\overset{\mathrm{}}{}}n^\gamma \underset{k=0}{\overset{2^n1}{}}|X_{(k+1)2^nT}X_{k2^nT}|^p<\mathrm{}\mathrm{a}.\mathrm{s}.$$
Before proving the lemma we recall a result of Monroe .
###### Definition 3.1
Let $`B_t`$ be a Brownian motion defined on a probability space $`(\mathrm{\Omega },𝔉,)`$. A stopping time $`T`$ is said to be minimal if for any stopping time $`ST`$, $`B(T)\stackrel{(d)}{=}B(S)`$ implies that a.s. $`S=T`$.
###### Theorem 3.1
\[14, Theorem 11\] Let $`(M_t)_{t0}`$ be a right continuous martingale. Then there is a Brownian motion $`(\mathrm{\Omega },𝔊_t,B_t)`$ and a family $`(T_t)`$ of $`𝔊_t`$-stopping times such that the process $`B_{T_t}`$ has the same finite distributions as $`M_t`$. The family $`T_t`$ is right continuous, increasing, and for each $`t`$, $`T_t`$ is minimal. Moreover, if $`M_t`$ has stationary independent increments then so does $`T_t`$.
###### Remark 3.1
It should be noted that the stopping times $`T_t`$ are not generally independent of $`B_t`$. However, in the case of $`\alpha `$-stable processes $`0<\alpha <2`$ one can use subordination to gain independence of the stopping times .
Proof of Lemma 3.2 Let $`(\tau _t)_{t0}`$ denote the collection of minimal stopping times for which
$$X_t\stackrel{(d)}{=}B_{\tau _t}.$$
The proof will be completed once it has been shown that
$$\underset{n=1}{\overset{\mathrm{}}{}}n^\gamma \underset{k=0}{\overset{2^n1}{}}|B_{\tau ((k+1)2^nT)}B_{\tau (k2^nT)}|^p<\mathrm{}\mathrm{a}.\mathrm{s}.$$
(29)
The following inequality holds because Brownian motion is $`(1/p^{})`$-Hölder continuous a.s. for $`p^{}>2`$:
$`|B_{\tau ((t_{k+1,n})}B_{\tau (t_{k,n})}|^p`$ $`C|\tau (t_{k+1,n})\tau (t_{k,n})|^{\frac{p}{p^{}}}`$ (30)
$`k=0,\mathrm{},2^n1,n1,\mathrm{a}.\mathrm{s}.`$
where $`t_{k,n}k2^nT`$ and $`2<p^{}<p`$.
\[15, Theorem 1\] shows that the index of the process $`\tau (s)`$ is half that of the Lévy process. Therefore, with probability one, $`\tau (s)`$ has finite $`(1+\delta )`$-variation for all $`\delta >0`$.
###### Theorem 3.2
\[14, Theorem 5\] If $`\tau `$ is a minimal stopping time and $`𝔼(B_\tau )=0`$, then $`𝔼(\tau )=𝔼(B_\tau ^2)`$.
Consequently the process $`(\tau _t)_{t0}`$ can be controlled in the following way:
$$𝔼\left(\tau _t\right)=𝔼\left(B_{\tau _t}^2\right)=𝔼\left(X_t^2\right)=t_{|x|<1}|x|^2\nu (dx)$$
(31)
where $`\nu `$ is the Lévy measure corresponding to the process $`X_t`$. From (31) and Theorem 3.1 we note that the process $`\tau _t`$ is a Lévy process whose Lévy measure, say $`\mu `$, satisfies the following:
$$_0^1x\mu (dx)<\mathrm{}.$$
From this result we deduce that the process $`\tau _t`$ a.s. has bounded variation. From \[16, Theorem 5\] we note that there is a positive constant $`A`$ such that
$$\left(\tau _tAt,t0\right)=1.$$
From the above bound and using the fact that $`\tau `$ has stationary independent increments one can show
$`\left(\tau (t_{k+1,n})\tau (t_{k,n})A(t_{k+1,n}t_{k,n})=A2^n|\tau (t_{k,n})\right)`$ $`=1,`$
$`\left({\displaystyle \underset{n1}{}}{\displaystyle \underset{k0}{\overset{2^n1}{}}}\left(|\tau (t_{k+1,n})\tau (t_{k,n})|A\mathrm{\hspace{0.17em}2}^n\right)\right)`$ $`=1.`$
Returning to (30) we see that
$`|B_{\tau ((t_{k+1,n})}B_{\tau (t_{k,n})}|^p`$ $`C|\tau (t_{k+1,n})\tau (t_{k,n})|^{\frac{p}{p^{}}}`$
$`CA\mathrm{\hspace{0.17em}2}^{n(\frac{p}{p^{}})}`$
which implies that
$`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma {\displaystyle \underset{k=0}{\overset{2^n1}{}}}|B_{\tau ((k+1)2^nT)}B_{\tau (k2^nT)}|^p`$ $`CA{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}n^\gamma 2^{n(\frac{p}{p^{}}1)}<\mathrm{}`$
due to $`p^{}`$ being chosen in the interval $`(2,p)`$. $`\mathrm{}`$
This lemma concludes the proof that the bound in (3) is finite, which shows that the area process a.s. has finite $`(p/2)`$-variation. $`\mathrm{}`$
In this section we have proved that the area process exists and has finite $`(p/2)`$-variation when $`(X_t)_{t0}`$ has the form (22). To prove Theorems 2.1, 2.2 we note that a general Lévy process has the form
$$X_t=at+B_t+L_t+\underset{\begin{array}{c}0s<t\\ |\mathrm{\Delta }X_s|1\end{array}}{}\mathrm{\Delta }X_s\mathrm{a}.\mathrm{s}.$$
So, we need to add area corresponding to the drift vector and the jumps of size greater than one. However, this part of the Lévy process has bounded variation and is piecewise smooth so there is no problem defining its area. Similarly, it has a.s. finite $`(p/2)`$-variation.
## Appendix A Homeomorphic flows
In this section we give a proof that the solutions, generated by (1) as the initial condition is varied, form a flow of homeomorphisms when the integrator is a continuous function. The proof modifies the one given in for the existence and uniqueness of solution to (1). The main idea is that one uniformly bounds a sequence of iterated maps which have projections giving the convergence of the solutions with two different initial points and bounding the difference of the solutions.
First, we need some notation.
###### Definition A.1
Let $`T^{(n)}(^d)`$ denote the truncated tensor algebra of length $`n`$ over $`^d`$. That is
$$T^{(n)}(^d)\underset{i=0}{\overset{n}{}}(^d)^i$$
where $`(^d)^0=`$ and $`T^{(\mathrm{})}(^d)`$ denotes the tensor algebra over $`^d`$.
Let $`\mathrm{\Delta }=[0,T]\times [0,T]`$. A map $`X:\mathrm{\Delta }T^{(n)}(^d)`$ will be called a multiplicative functional of size $`n`$ if for all times $`s<t<u`$ in $`[0,T]`$ the following relation holds in $`T^{(n)}(^d):`$
$$X_{st}X_{tu}=X_{su}$$
and $`X_{st}^{(0)}1`$.
A map $`X:\mathrm{\Delta }T^{(n)}(^d)`$ is called a classical multiplicative functional if $`tX_tX_{0t}^{(1)}`$ is continuous and piecewise smooth and
$$X_{st}^{(i)}=_{s<u_1<\mathrm{}<u_i<t}𝑑X_{u_1}\mathrm{}𝑑X_{u_i}$$
(32)
where the right hand side is a Lebesgue-Stieltjes integral. We denote the set of all classical multiplicative functionals in $`T^{(n)}(^d)`$ by $`S^{(n)}(^d)`$.
###### Definition A.2
We call a continuous function $`\omega :\mathrm{\Delta }^+`$ a control function if it is super-additive and regular, that is,
$`\omega (s,t)+\omega (t,u)`$ $`\omega (s,u)`$ $`s<t<u[0,T]`$
$`\omega (s,s)`$ $`=0`$ $`s[0,T].`$
###### Example A.1
Let $`X`$ be a path of strong finite $`p`$-variation. Then we can define the following control function:
$$\omega (s,t)X_{p,[s,t]}^p.$$
(33)
###### Definition A.3
A functional $`X=(1,X^{(1)},\mathrm{},X^{(n)})`$ defined on $`T^{(n)}(^d)`$ where $`n=[p]`$ is said to have finite $`p`$-variation if there is a control function $`\omega `$ such that
$$|X_{st}^{(i)}|\frac{\omega (s,t)^{i/p}}{\beta (i/p)!}(s,t)\mathrm{\Delta },i=1,\mathrm{},n$$
(34)
for some sufficiently large $`\beta `$ and $`x!\mathrm{\Gamma }(x+1)`$.
###### Theorem A.1
\[10, Theorem 2.2.1\] Let $`X^{(n)}`$ be a multiplicative functional of degree $`n`$ which has finite $`p`$-variation, with $`n[p]`$ ($`[p]`$ denotes the integer part of $`p`$). Then for $`m>n`$ there is a unique multiplicative extension $`X^{(m)}`$ in $`T^{(m)}(^d)`$ which has finite $`p`$-variation.
###### Remark A.1
The above theorem shows that once a sufficient number of low order integrals associated to a path $`X_t`$ have been defined, then the remaining iterated integrals of $`X_t`$ are defined.
###### Definition A.4
We call a multiplicative functional $`X:\mathrm{\Delta }T^{(n)}(^d)`$ geometric if there is a control function $`\omega `$ such that for any positive $`ϵ`$ there exists a classical multiplicative functional $`Y(ϵ)`$ which approximates $`X`$ in the following way:
$$\left|\left(X_{st}Y_{st}(ϵ)\right)^{(i)}\right|ϵ\omega (s,t)^{i/p}i=1,\mathrm{},n=[p].$$
We denote the class of geometric multiplicative functionals with finite $`p`$-variation by $`\mathrm{\Omega }G(^d)^p`$.
###### Example A.2
Let $`W_t`$ be an $`^d`$-valued Brownian motion. Then the following functional $`W`$ defined on $`T^{(2)}(^d)`$ belongs to $`\mathrm{\Omega }G(^d)^p`$ for any $`p>2`$.
$$W_{st}(1,W_tW_s,_{s<u_1<u_2<t}dW_{u_1}dW_{u_2})$$
(35)
where $`dW_u`$ denotes the Stratonovich integral. It should be noted that if one replaced the Stratonovich differential in (35) by the Itô differential then one would not get an element of $`\mathrm{\Omega }G(^d)^p`$. This is due to the quadratic variation term which occurs in the symmetric part of the area process
$$W_{st}^{(2)}=_{s<u_1<u_2<t}𝑑W_{u_1}𝑑W_{u_2}.$$
It was shown in that one had sufficient control of the above functional to generate path-wise solutions to SDEs driven by a Brownian motion. This control was derived from a moment condition in the same spirit as Kolmogorov’s criterion for Hölder continuous paths. The moment condition was verified for the above area by the use of known stochastic integral results, though one could also derive it from a construction depending on the linearly interpolated Brownian motion.
There are two stages to defining the integral against a geometric multiplicative functional. The first gives a functional which is almost multiplicative (see for definition). The second associates, uniquely, a multiplicative functional to the almost multiplicative functional.
###### Theorem A.2
There is a unique geometric multiplicative functional $`Y`$ which we call the integral of the 1-form $`\theta `$ against the geometric multiplicative functional $`X`$. We denote this by
$$Y_{st}_s^t\theta (X_u)\delta X.$$
###### Corollary A.1
One has the following control on the $`p`$-variation of $`Y`$:
$$\left|\left(_s^t\theta (X_u)\delta X\right)^{(i)}\right|\left(C\omega (s,t)\right)^{i/p}/\left(\beta \left(i/p\right)!\right)i=1,\mathrm{},[p]$$
(36)
where $`C`$ depends on $`p,f_{\mathrm{Lip}(\gamma )},\gamma ,\lambda ,\beta ,L`$ and $`[p]`$.
The estimate is derived from estimating both the almost multiplicative functional and the difference of it from the integral.
We now state two lemmas which help prove that the solutions of (1) are homeomorphic flows when the initial condition is varied.
###### Lemma A.1
Let $`X`$ be in $`\mathrm{\Omega }G(^d)^p`$ controlled by a regular $`\omega _0`$. Let $`f:^n\mathrm{hom}(^d,^n)`$ be a $`\mathrm{Lip}(\gamma )`$ map for some $`\gamma >p`$. Let $`Y_{st}^{(i)},i=1,2`$ denote the element in $`\mathrm{\Omega }G(^n)^p`$ which solves the rough integral equation
$$Y_{st}^{(i)}=_s^tf(Y^{(i)})\delta X$$
with initial condition $`Y_0^{(i)}=a_i,i=1,2`$. Let $`W_{st}`$ be the multiplicative functional which records the difference in the multiplicative functionals $`Y_{st}^{(1)}`$ and $`Y_{st}^{(2)}`$. Then
$$\left|W_{st}^{(i)}\right|\theta ^i\frac{\omega (s,t)^{(i/p)}}{\beta (i/p)!}i1,$$
(37)
where $`\theta =|a_1a_2|`$, $`\omega C\omega _0`$, the constant $`C`$ depends on $`p,`$ $`f_{\mathrm{Lip}(\gamma )},`$ $`\beta ,`$ $`\gamma `$. The bound holds for all times $`st`$ on the interval $`J\{u:\omega (0,u)1\}`$.
###### Lemma A.2
With the assumptions of Lemma A.1 one can estimate the difference of the increments of $`Y_{st}^{(1)}`$ and $`Y_{st}^{(2)}`$ for any pair of times $`0s<t`$ which satisfy $`\omega (s,t)1`$ as follows:
$`|Y_{st}^{(1)}Y_{st}^{(2)}|`$ $`\theta \mathrm{exp}\left[{\displaystyle \frac{1}{\beta (1/p)!}}\left(\omega (0,s)+\omega (0,s)^{(1/p)}\right)\right]{\displaystyle \frac{\omega (s,t)^{(1/p)}}{\beta (1/p)!}}`$
In particular for any $`t>0`$ one has:
$$|Y_t^{(1)}Y_t^{(2)}||a_1a_2|C(t).$$
(38)
Now we can prove that the solutions form a flow of homeomorphisms as the initial condition is varied.
Proof of Theorem 1.3 The continuity of solutions follows from Lemma A.2. It remains to show that the inverse map exists and is continuous. This can be checked by repeating all the previous arguments using the reversed path $`(X_{ts})_{0st}`$ as the integrator. $`\mathrm{}`$
The induction part of the proof of Lemma A.1 will require the following lemma about rescaling:
###### Lemma A.3
Let $`X`$ be a multiplicative functional in $`T^{([p])}(^d)`$ which is of finite $`p`$-variation controlled by $`\omega `$. Let $`(X,Y)`$ be an extension of $`X`$ to $`T^{([p])}(^d^n)`$ of finite $`p`$-variation controlled by $`K\omega `$. Then $`(X,\varphi Y)`$ is controlled by
$$\mathrm{max}\{1,\varphi ^{kp/i}K:1ki[p]\}\omega $$
where $`\varphi `$. In particular, if $`\varphi K^{[p]/p}1`$ then $`(X,\varphi Y)`$ is controlled by $`\omega `$.
Proof of Lemma A.1 We set up an iteration scheme of multiplicative functionals which we will bound uniformly, by induction. A projection of the sequence proves that a Picard iteration scheme converges to the solutions of (1) starting from $`a_1`$ and $`a_2`$. Another projection shows that the difference of these solutions is bounded.
Let $`ϵ>0`$ and $`\eta >1`$. Let $`V_{st}^{(1)}`$ be the geometric multiplicative functional given by
$`V_{st}^{(1)}`$ $`(Z_{st}^{(1)(1)},Y_{st}^{(1)(1)},Y_{st}^{(1)(0)},Z_{st}^{(2)(1)},Y_{st}^{(2)(1)},Y_{st}^{(2)(0)},W_{st}^{(1)},ϵ^1X_{st})`$
$`=({\displaystyle _s^t}f(a_1)\delta Xa_1,{\displaystyle _s^t}f(a_1)\delta X,a_1,{\displaystyle _s^t}f(a_2)\delta Xa_2,`$
$`{\displaystyle _s^t}f(a_2)\delta X,a_2,{\displaystyle _s^t}f(a_1)f(a_2)\delta X,ϵ^1X_{st}).`$
The iteration step is a two stage process. Given $`V^{(m)}`$ we set
$$\stackrel{~}{V}^{(m+1)}=k_\theta ^m(V^{(m)})\delta V^{(m)}$$
where $`k_\theta ^m`$ is the $`1`$-form on $`((^n)^7^d)`$ given by
$`k_\theta ^m(a_1,\mathrm{},a_8)(dA_1,\mathrm{},dA_8)`$
$`=(a_1g(a_2,a_3)dA_8,dA_3+\eta ^mdA_1,dA_2,a_4g(a_5,a_6)dA_8,`$
$`dA_6+\eta ^mdA_4,dA_5,\theta ^1g(a_2,a_4)dA_8,dA_8).`$
$`g(x,y)`$ is the 1-form appearing in \[9, Lemma 3.2\] which satisfies the following relation with respect to $`f`$:
$$f^i(x)f^i(y)=\underset{j}{}(xy)^jg^{ij}(x,y).$$
$`\stackrel{~}{V}^{(m+1)}`$ is well defined because $`g`$ and $`k_\theta ^m`$ are both $`\mathrm{Lip}(\gamma )`$ for some $`\gamma >p1`$.
We define $`V^{(m+1)}`$ to be the geometric multiplicative functional obtained by rescaling the first and fourth components of $`\stackrel{~}{V}^{(m+1)}`$ by $`ϵ\eta `$ and the seventh component by $`ϵ`$.
The uniform bound on the iterates $`(V^{(m)})_{m1}`$ will be obtained by induction. $`X`$ is controlled by a regular $`\omega _0`$ so there exists a constant $`C`$ such that $`V^{(1)}`$ is controlled by $`\omega C\omega _0`$. Suppose that $`V^{(k)}(km)`$ are controlled by $`\omega `$. From (Corollary A.1) there is a constant $`C_1`$ such that $`\stackrel{~}{V}^{(m+1)}`$ is controlled by $`C_1\omega `$. If we choose $`ϵ>0,\eta >1`$ such that $`ϵC_1^{\frac{[p]}{p}}`$ and $`ϵ\eta C_1^{\frac{[p]}{p}}`$, then Lemma A.3 implies that $`V^{(m+1)}`$ is controlled by $`\omega `$, completing the induction step.
The uniform control on the iterates $`V^{(m)}`$ ensures the convergence of $`\{Y^{(i)(m)}\}_{m1}`$ to the solutions of
$`dY_t^{(i)}`$ $`=f(Y_t^{(i)})dX_tY_0^{(i)}=a_i,i=1,2.`$
Through the definition of $`\{{}_{}{}^{\theta }W_{}^{(m)}\}_{m1}`$, the sequence at the level of the paths will converge to the scaled difference of the two solutions $`\theta ^1(Y^{(2)}Y^{(1)})`$. For $`s,t`$ in $`J`$ one has
$$|{}_{}{}^{\theta }W_{st}^{(i)}|\frac{\omega (s,t)^{i/p}}{\beta (i/p)!}i=1,\mathrm{},[p],$$
which implies that
$$|W_{st}^{(i)}|\theta ^i\frac{\omega (s,t)^{i/p}}{\beta (i/p)!}i=1,\mathrm{},[p].\mathrm{}$$
Proof of Lemma A.2 We define the following set of times:
$`t_00`$ $`\text{and }t_jinf\{u>t_{j1}:\omega (t_{j1},u)=1\}j\{1,\mathrm{},n(s)\}`$ (39)
where $`n(s)\mathrm{max}\{j:t_js\}\text{and}t_{n(s)+1}=s.`$
We solve the differential equation starting from $`s`$ and use (37) to show that
$$|W_{st}^k|K(s)^k\frac{\omega (s,t)^{(k/p)}}{\beta (k/p)!}$$
where $`K(s)`$ is an upper bound on the supremum over all the possible differences of the paths $`|Y_s^{(1)}Y_s^{(2)}|,`$ at time $`s`$. The bound $`K(s)`$ is derived recursively by considering the analogous upper bound for the difference of the solutions to the differential equation over the time interval $`[t_{i1},t_i]`$ given below:
$`|Y_{t_i}^{(1)}Y_{t_i}^{(2)}|`$ $`|Y_{t_{i1}}^{(1)}Y_{t_{i1}}^{(2)}|+|W_{t_{i1}t_i}|`$
$`|Y_{t_{i1}}^{(1)}Y_{t_{i1}}^{(2)}|\left(1+{\displaystyle \frac{\omega (t_{i1},t_i)^{(1/p)}}{\beta (1/p)!}}\right)`$
which implies that
$`K(t_j)`$ $`K(t_{j1})\left\{1+{\displaystyle \frac{\omega (t_{j1},t_j)^{(1/p)}}{\beta (1/p)!}}\right\}j=1,\mathrm{},n(s)+1.`$
Therefore
$`|W_{st}^k|`$ $`K(t_0)^k{\displaystyle \underset{j=1}{\overset{n(s)+1}{}}}\left\{1+{\displaystyle \frac{\omega (t_{j1},t_j)^{(1/p)}}{\beta (1/p)!}}\right\}^k{\displaystyle \frac{\omega (s,t)^{(k/p)}}{\beta (k/p)!}}`$
$`\theta ^k\mathrm{exp}\left[k\left({\displaystyle \underset{j=1}{\overset{n(s)}{}}}{\displaystyle \frac{\omega (t_{j1},t_j)^{(1/p)}}{\beta (1/p)!}}+{\displaystyle \frac{\omega (t_{n(s)},s)^{(1/p)}}{\beta (1/p)!}}\right)\right]{\displaystyle \frac{\omega (s,t)^{(k/p)}}{\beta (k/p)!}},`$
noting that $`\omega (t_{j1},t_j)=1`$ and using the sub-additivity of $`\omega `$ we obtain
$`\theta ^k\mathrm{exp}\left[{\displaystyle \frac{k}{\beta (1/p)!}}\left(\omega (0,s)+\omega (0,s)^{(1/p)}\right)\right]{\displaystyle \frac{\omega (s,t)^{(k/p)}}{\beta (k/p)!}}.`$
By considering the above bound at the level of the paths $`(k=1)`$ and repeatedly using the triangle inequality one deduces (38). $`\mathrm{}`$
Department of Statistics
University of Oxford
1 South Parks Road
Oxford OX1 3TG
England
E-mail: williams@stats.ox.ac.uk |
warning/0001/cond-mat0001411.html | ar5iv | text | # Spinless fermions and charged stripes at the strong-coupling limit
\[
## Abstract
Spinless fermions on a lattice with nearest-neighbor repulsion serve as a toy version Hubbard model, and have a symmetry-broken even/odd superlattice at half-filling. At infinite repulsion, doped holes form charged stripes which are antiphase walls (as noted by Mila in 1994). Exact-diagonalization data for systems up to 36 sites around 1/4 filling, and also for one or two holes added to a stripe of length up to 12, indicate stability of the stripe-array state against phase separation. In the boson version of the model, the same behavior can be stabilized by addition of a four-fermi term.
\]
Forty years of study have not yet produced a complete understanding of the phase diagram of the Hubbard model, the simplest nontrivial paradigm of interacting spinfull fermions. The spinless lattice fermion model is a simpler and more tractable analog which retains many Hubbard-model properties, much as the Ising model stands in for the $`n`$-component magnet in critical phenomena: understanding of the spinless model may provide fresh viewpoints of the Hubbard model, or new tests of known methods. Spinless models also arise naturally for ferromagnetic materials in which one of the spin-split bands is completely full or completely empty, such as magnetite or “half-metallic” manganites .
Our aim is to promote the systematic study of this model’s phase diagram, which is almost untouched in the literature . As a beginning, this paper argues that, in the strong-coupling limit, the spinless model possesses a phase with the quantum-fluctuating, hole-rich antiphase domain walls known as “stripes.” Such stripes are an active topic in the Hubbard or $`tJ`$ models, particularly since stripes were observed in cuprates and seem related to incommensurate correlations found in high-temperature superconductors.
Let us take a square lattice model with Hamiltonian
$$=t\underset{ij}{}(c_i{}_{}{}^{}c_{j}^{}+c_j{}_{}{}^{}c_{i}^{})+V\underset{ij}{}\widehat{n}_i\widehat{n}_j$$
(1)
Here $`c_i^{}`$ and $`c_i`$ are creation/annihilation operators on site $`i`$, $`\widehat{n}_ic{}_{}{}^{}{}_{i}{}^{}c_{i}^{}`$, and “$`\mathrm{}`$” counts each nearest-neighbor pair once. In most places we will consider hard-core bosons in parallel with fermions . From here on, we take $`V/|t|=\mathrm{}`$ so neighboring particles are simply forbidden, and $`t`$ is the only energy scale. (This constraint amounts to adopting a hardcore radius just over 1 lattice constant.)
Phase diagram as function of $`n`$ Consider first the dilute limit, $`n0.15`$. When $`V\mathrm{}`$, the Hartree-Fock approximation gives absurd results; in reality the renormalized interaction of two particles is of order $`t`$ and a Bose or Fermi liquid is expected, as when the on-site repulsion $`U\mathrm{}`$ in the Hubbard model .
At the other extreme, the dense limit ($`n=1/2`$, half-filling) admits only the two microstates with the $`\sqrt{2}\times \sqrt{2}`$ checkerboard pattern, called the “CDW” (“charge-density wave”) order . An Ising symmetry breaking between even/odd lattices is exhibited, the spinless model’s cartoon of the Heisenberg antiferromagnetic order at half-filling in the large-$`U`$ Hubbard model.
Stripes in a hard-core model – Now consider light hole doping, $`1/2n1`$. An isolated hole is immobile and gains no hopping energy in the CDW state (see Fig. 1). As Mila observed, a droplet including $``$3 holes can fluctuate but is still confined to a circumscribed rectangle with edges along the $`45^{}`$ directions, since a particle is prevented from hopping away from a CDW domain surface oriented along $`\{11\}`$ (Fig. 1).
The natural way to dope holes is a “stripe”, an antiphase domain wall with charge 1/2 hole per unit length. This permits hops (arrows in Fig. 1) which implement stripe fluctuations. A single stripe’s path can be parametrized as a unique function $`y(x)`$ (hopping never generates overhangs). Then $`y(x+1)y(x)=\pm 1`$, and the steps up or down can be represented by a string of corresponding $``$ and $``$ arrows. A particle hop has the effect $``$ and so the Hamiltonian of a single stripe maps exactly to the spin-1/2 XX chain with exchange $`J_{}=t`$ for the X and Y spin components. That model is exactly soluble by a well-known mapping, whereby the up (down) spins map to noninteracting spinless fermions (empty sites), respectively, in one dimension. It follows that the energy of the (coarse-grained) stripe is
$$𝑑x[\sigma _0+\frac{1}{2}K(dy/dx)^2+\mathrm{}],$$
(2)
where $`\sigma _0=(2/\pi )t`$ and the stripe stiffness $`K=(\pi /2)t`$. As usual, the sound velocity of the stripe’s capillary waves is the Fermi velocity $`v=2t/\mathrm{}`$ of the 1D noninteracting fermions. Knowing $`K`$ and $`v`$ allows us to compute the fluctuations of the Fourier mode at each wavevector $`q`$, as for any harmonic string: $`|y_q|^2=\mathrm{}v/2K|q|`$. The general result is that one such “Gaussian” quantum-fluctuating stripe has divergent fluctuations,
$$|y(x)y(0)|^2=(v/2\pi K)\mathrm{ln}|x|+\mathrm{const},\text{as }x\mathrm{}$$
(3)
where $`v/2\pi K=2/\pi ^2`$ in this case.
The thermodynamic phase at $`1/2n1`$ could then be an array of stripes all parallel (on average) to either the $`\widehat{x}`$ or $`\widehat{y}`$ axis. They have only a contact interaction, so the array’s long-range order depends on stripe collisions, which surely exist since isolated stripes have divergent fluctuations (eq. (3)).
Prior analytic work suggested that spinless fermions in $`d=2`$, when doped away from half-filling, develop an incommensurate ordering wavevector slightly off from $`(\pi ,\pi )`$). Expanding around the $`d=\mathrm{}`$ limit, small doping led to coexistence between the half-filled CDW and a slightly incommensurate state (but not at $`V/t\mathrm{}`$). We conjecture these incommensurate phases, in $`d=2`$, consist of stripe arrays. In the hard-core boson model near half-filling, in a regime $`0<V/t<\mathrm{}`$, the uniform CDW phase is asserted to phase-separate upon doping . The dense coexisting state is just as plausible, a priori, to be a stripe array as the phase-separated state that was assumed. .
Stability estimates – The key question is whether (or when) the stripe-array is stable, compared to a phase-separated state in which the CDW state and the dilute (=hole-rich) liquid coexist. In the case of the Hubbard model, it was argued that doping invariably leads to phase separation except when it is “frustrated” by the long-range Coulomb force. Contrarily, it was argued that holes in fluctuating stripes may gain more kinetic energy than they would in a phase-separated state .
To decide the issue of coexistence, one first plots the energy per site $`E^{\mathrm{liq}}(n)`$ and $`E^{\mathrm{sa}}(n)`$ for the low-density liquid and the stripe-array, respectively, which should look like Fig. 2 for either fermions or hardcore bosons. We have
$$E^{\mathrm{liq}}(n)=(4t)n+A_2n^2+A_3n^3+\mathrm{}$$
(4)
where the leading coefficient is the bottom of the single-particle band. With increasing density, the energy $`E^{\mathrm{liq}}(n)`$ turns upwards and becomes small around $`n=0.3`$ as the hopping becomes “jammed” (neighbor sites become forbidden due to other adjacent particles). The matrix elements contribute with the same sign in the boson ground state but can’t in the fermion case, so $`E_{}^{\mathrm{liq}}{}_{\mathrm{fermions}}{}^{}>E_{}^{\mathrm{liq}}{}_{\mathrm{bosons}}{}^{}`$.
On the other hand, in the conjectured stripe array near half filling,
$$E^{\mathrm{sa}}(n)[\sigma _0+\varphi (d)](12n)$$
(5)
where $`\sigma _0`$ is the energy per unit $`x`$-length from (2); the second factor is the length of stripe per unit area. The mean stripe separation is $`d(12n)^1`$, and $`\varphi (d)`$ parametrizes the energy cost per unit length from collisions of adjacent stripes, so $`\varphi (d)0`$ as $`d\mathrm{}`$. Thus, the chemical potential in the limit of separated stripes is $`\mu ^{}dE^{\mathrm{sa}}(n)/dn|_{n=1/2}=(4/\pi )t=1.273t`$. We emphasize that the form and the leading coefficients in (4) and (5) are the same for fermions and hardcore bosons. \[Indeed, the fermion and boson models are identical in the single-stripe sector, since the $`V=\mathrm{}`$ constraint prevents any permutations; this identity extends to the single stripe with one extra hole, in which case only even permutations are accessible .\]
There are three necessary conditions for the stability of the stripe array: we have strong evidence for each of them, from exact diagonalizations of systems with 20 to 72 sites. These are far too small for direct observation of a fluctuating stripe array, or of the coexistence of the liquid and dense phases; yet they are large enough to yield some of the parameters which the phase diagram can be calculated from. From here on we use units $`t1`$.
The first stability condition is that stripes repel, i.e. $`\varphi (d)>0`$; stripe attraction would suggest instability to a domain of liquid phase, which is scarcely distinguishable from a bundle of self-bound stripes. In both the fermion and boson cases we diagonalized $`L\times L^{}`$ systems, for $`L=4`$, and $`L^{}=6,8,10`$, as well as $`L=L^{}=6`$, doped with $`L`$ holes so that two stripes run in the short direction, and thus $`d=L^{}/2`$. Define a stripe interaction per unit length $`\varphi _{\mathrm{eff}}(L^{}/2)(E2E^{\mathrm{str}}(L))/2L`$, where $`E^{\mathrm{str}}(L)`$ is the energy of one isolated stripe of finite (even) length $`L`$. ($`E^{\mathrm{str}}(L)`$ is calculated exactly using the 1D spinless-fermion representation.) Indeed, $`\varphi _{\mathrm{eff}}(d)`$ was positive and (for $`L=4`$) decreasing with $`d`$.
Next, a stripe can contain extra holes, which move as quasiparticles with an excitation gap $`\mathrm{\Delta }`$. The second stability condition is
$$\mathrm{\Delta }>\mu ^{}$$
(6)
If not, further doping would add holes to existing stripes rather than form new ones, again suggesting a tendency to form phase-separated droplets. From diagonalizations we measured the excitation energy of one added hole, $`\mathrm{\Delta }(L,L^{})`$, with even $`L`$ and odd $`L^{}`$ in the range $`[4,8]`$ for fermions, or $`[4,12]`$ for bosons; as mentioned above, this is strictly independent of statistics . From this we extrapolated, first $`L^{}\mathrm{}`$ using $`\mathrm{\Delta }(L,L^{})=\mathrm{\Delta }(L)+A_\mathrm{\Delta }(L)e^{L^{}/l(L)}`$ and then $`L\mathrm{}`$ using $`\mathrm{\Delta }(L)=\mathrm{\Delta }+B_\mathrm{\Delta }/L`$. We found $`\mathrm{\Delta }=0.65(5)`$, which comfortably satisfies (6).
We also analyzed the two-hole energy $`\mathrm{\Delta }_2(L)`$, for $`L=4,6,8,10`$ only; with two holes, the $`L`$ dependence is more like $`1/2^L`$ than $`1/L`$. Extrapolating to $`L=\mathrm{}`$, yielded $`\mathrm{\Delta }_2=1.42(7)`$ for bosons and $`1.44(9)`$ for fermions. Note $`\mathrm{\Delta }_22\mathrm{\Delta }0.1`$, i.e. hole binding is insignificant when $`L10`$; we think it is a real effect in a large system, since holes on a stripe can be collected into a (hole-free) vertical segment of the stripe. (Since the stripe’s $`90^{}`$ kinks cost energy, an array of parallel stripes will still be the thermodynamic phase in a large system.)
The third condition is the crucial one: as shown in Fig. 2, the dense phase coexisting with the liquid must not be the CDW, which would preempt a stripe-array phase. That is,
$$\mu ^{}>\mu ^{\mathrm{LC}}.$$
(7)
where $`\mu ^{\mathrm{LC}}`$ is the slope of the trial liquid-CDW tie-line tangent to $`E^{\mathrm{liq}}(n)`$ and passing through $`(n=1/2,E=0)`$.
To test (7), the equation of state $`E^{\mathrm{liq}}(n)`$ is required. We exactly diagonalized all rectangular lattices with $`L,L^{}4`$ and $`LL^{}=20`$ to 36 sites, and with occupation in the range $`0.20n<0.3`$, and fitted the results to (4). We obtained $`(A_2,A_3)(9.45\pm 0.6,5\pm 2)`$ for bosons and $`(11.25\pm 0.6,1\pm 2)`$ for fermions. This implies $`\mu _{}^{\mathrm{LC}}{}_{\mathrm{bosons}}{}^{}=1.33(2)`$ and $`\mu _{}^{\mathrm{LC}}{}_{\mathrm{fermions}}{}^{}=1.25(2)`$. (Here the errors are estimated by varying the subset of data used for the fit.) For either bosons or fermions, coexistence with the CDW would occur at $`n_c0.24`$.
Hence, stripes are unstable in the boson case and (very likely) stable in the fermion case, but close enough to the boundary in either case that the balance can be tipped either way by the small perturbation $`\stackrel{~}{t}_c`$, discussed later.
Exotic states? – Like the Hubbard model, the spinless fermion model may be extended by adding other hopping terms to the Hamiltonian, which might stabilize additional phases. Many of these terms have the form
$$\stackrel{~}{t}_x(c{}_{}{}^{}{}_{j}{}^{}c_{i}^{}+c{}_{}{}^{}{}_{i}{}^{}c_{j}^{})\widehat{n}_k,$$
(8)
where $`(i,j,k)`$ are three different sites arranged as in Fig. 3, and $`x`$ is “$`a`$”, “$`b`$”, or “$`c`$” for the hops shown in the corresponding parts of Fig. 3. For example, when $`V`$ is large but finite, hops are possible to a neighbor’s neighbor with $`\stackrel{~}{t}_a=\stackrel{~}{t}_b=t^2/V`$ as in Fig. 3(a) and (b), analogous to similar terms of order $`t^2/U`$ when the $`t`$-$`J`$ model is derived from the Hubbard model. This spinless analog of the $`t`$-$`J`$ model, in which virtual states with neighbor pairs are projected out, will be the natural starting point to study phenomena at large (but not infinite) $`V`$, e.g. the mobility of lone holes.
In the fermion model, one could artificially take $`\stackrel{~}{t}_bt`$ (analogous to $`Jt`$ in the $`t`$-$`J`$ model); then the term Fig. 3 (b) naturally favors superconductivity. Namely, fermions form tightly bound $`p`$-wave pairs, separated by $`\sqrt{2}`$; these composite bosons hop with bandwidth $`8\stackrel{~}{t}_b`$, and Bose-condense in the usual fashion. Thus $`\stackrel{~}{t}_b`$ is analogous to negative $`U`$ in the Hubbard model, in that superconductivity is put in “by hand”. But it is a plausible speculation that, in the highly correlated liquid at $`n0.25`$, BCS superconductivity appears even with $`\stackrel{~}{t}_bt`$.
Finally, consider the hopping $`\stackrel{~}{t}_c`$ of Fig. 3 (c) which just modifies the amplitude of already possible hops. This tends to stabilize (destabilize) stripes according to whether $`\stackrel{~}{t}_c`$ has the same (opposite) sign as $`t`$, since every allowed hop in a stripe is surrounded by particles on all four possible “$`k`$” sites. (Compare Fig. 3(c) with Fig. 1). Hence the stripe energy $`\sigma _0`$ and $`\mu ^{}`$ get multiplied by a factor $`(1+4\stackrel{~}{t}_c/t)`$. On the other hand, assuming that each “$`k`$” site is about $`1/3`$ occupied in the liquid at $`n=0.25`$, it follows that $`E^{\mathrm{liq}}(n)`$ is multiplied by about $`(1+4\stackrel{~}{t}_c/3t)`$. If so, the critical perturbation where $`\mu ^{\mathrm{LC}}=\mu ^{}`$ (so the stripe phase appears or vanishes) is only $`\stackrel{~}{t}_c/t0.02`$ for bosons or $`0.007`$ for fermions, using our values of $`\mu ^{\mathrm{LC}}`$ quoted above.
Discussion – To establish the occurrence of stripes in the $`V=\mathrm{}`$ system, we addressed, by exact diagonalizations, (i) stripe-stripe interactions; (ii) the energy of a single hole, as well as hole-hole interactions, on a stripe; and (iii) the medium density liquid regime. (iv) hole-hole interactions on a stripe The enormous reduction of Hilbert space due to the nearest-neighbor exclusion (at $`V=\mathrm{}`$), as well as the lack of spin, permits numerical explorations at system sizes much larger than would be possible in the Hubbard model – vital not just for studying stripes, but any microscopically inhomogenous states. Monte Carlo simulation of the stripe phase is straightforward for the hardcore boson case. Boson results are valid for fermions too, when the stripe separation $`d`$ is large and the density of “extra holes” on each stripe is low, since particles do not exchange in this limit . As for the fermion case, the new “meron-cluster” Monte Carlo algorithm cancels the sign problem for a limited class of models including spinless fermions, but not the basic Hubbard model .
The ultimate aim of microscopic simulations should be to extract macroscopic parameters, e.g. the stripe stiffness $`K`$ or the stripe contact repulsion. This is more straightforward than in the spin-full (Hubbard or $`t`$-$`J`$) case, where the inter-stripe domains contain gapless spin-wave excitations . These parameters may be input to analytic explorations of the interesting anisotropic conductivity of the quantum-fluctuating stripe array , also simpler in the spinless case.
More broadly, it is a challenge to test for the exotic phases we mentioned in connection with Fig. 3. In the medium-density regime $`n0.2`$, strong correlations of some sort are essential to minimize the hopping energy. These might be prosaic, e.g. a $`\sqrt{5}\times \sqrt{5}`$ superlattice, but the following possibilities are realizable, in principle, even in a spinless model: (i) orbital magnetism (spontaneous circulating currents around plaquettes); (ii) $`p`$-wave superconductivity (see the speculations on Fig. 3(b)); or (iii) the analog of spin-charge separation, the spinon being replaced by a spinless particle that carries Fermi statistics . If it transpires that such states are hard to stabilize without spin, that would shed additional light on the Hubbard model; contrariwise, if they are stabilized, they may be easier to study in the spinless case, free from any background of low-energy spin excitations.
A crude comparison may be made of the spinless-fermion model with the infinite-$`U`$ Hubbard model in the dilute regime. In that case, each fermion excludes one site (its own) from half of the other fermions, not counting the exclusion built in by Fermi statistics. In the present spinless model each fermion excludes four sites from all other fermions, so in a sense the hole-rich metal phase is “jammed” eight times more effectively than in the Hubbard case. We expect, then, that kinetic-energy-stabilized stripes are far more robust in the present model than in the large-$`U`$ Hubbard case. In fact they are practically marginal in the present model, so this encourages the opinion that stripes are not stable in the short-range Hubbard (or $`t`$$`J`$) model.
We acknowledge support by the National Science Foundation under grants DMR-9981744 and PHY94-07194. C. L. H. thanks R. McKenzie, G. Uhrig, D. Scalapino, D. Khomski, M. Troyer, and G. G. Batrouni for helpful discussions. |
warning/0001/math0001065.html | ar5iv | text | # Incidence algebras of simplicial complexes
## Introduction
This paper, being motivated by physical problems, brings together the issues which traditionally belong to ‘disjoint’ areas of mathematics: combinatorics and differential moduli and, besides that, uses the notation from quantum mechanics.
It is shown that simplicial complexes resemble differential manifolds from the algebraic point if view, namely, their incidence algebras are similar to algebras of exterior differential forms on manifolds: they are graded, and possess an analog of Cartan differential.
To make the paper self-consistent, I begin it with an outline of basic definitions and results.
### Dirac notation.
Let $``$ be a finite-dimensional linear space with a basis labelled by an index set $`𝒦`$, and $`^{}`$ be its dual. Write down the elements of $``$ as
$$hh=\underset{P𝒦}{}c_P\mathbf{|}P\mathbf{}$$
Since the dimension of $``$ is finite, the same index set $`𝒦`$ is used to label the dual basis in the space $`^{}`$
$$h^{}^{}h^{}=\underset{P𝒦}{}c_P\mathbf{}P\mathbf{|}$$
such that
$$\mathbf{}P\mathbf{}Q\mathbf{}=\delta _{PQ}=\{\begin{array}{cc}1,\hfill & \text{if }P=Q\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(1)
The elements of $`^{}`$ are called bra-vectors and the elements of $``$ are ket-vectors (the terms derived from splitting the word ‘bracket’).
Let $`A:`$ be a linear operator and $`A^{}:^{}^{}`$ be its adjoint. In the dual bases we have for any $`P,Q𝒦`$:
$$\left(A^{}(\mathbf{}P\mathbf{|})\right)\mathbf{|}Q\mathbf{}=\mathbf{}P\mathbf{|}\left(A(\mathbf{|}Q\mathbf{})\right)$$
(2)
therefore with no confusion we can use the same symbol, say, $`a`$ for both $`A`$ and its adjoint $`A^{}`$:
$$\begin{array}{cccc}A:& \hfill \mathbf{|}Q\mathbf{}& & a\mathbf{|}Q\mathbf{}\hfill \\ A^{}:& \hfill \mathbf{}P\mathbf{|}& & \mathbf{}P\mathbf{|}a\hfill \end{array}$$
and the identity (2) reads
$$(\mathbf{}P\mathbf{|}a)\mathbf{|}Q\mathbf{}=\mathbf{}P\mathbf{|}(a\mathbf{|}Q\mathbf{})=\mathbf{}P\mathbf{|}a\mathbf{|}Q\mathbf{}=a_{pq}$$
Then $`A`$ can be written down as
$$A=\underset{P,Q𝒦}{}a_{PQ}\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}$$
and the product is calculated in accordance with (1):
$$AB=\left(\underset{P,Q𝒦}{}a_{PQ}\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}\right)\left(\underset{R,S𝒦}{}b_{RS}\mathbf{|}R\mathbf{}\mathbf{}S\mathbf{|}\right)=\underset{P,Q,S𝒦}{}a_{PQ}b_{QS}\mathbf{|}P\mathbf{}\mathbf{}S\mathbf{|}$$
(3)
This notation was introduced by P.A.M.Dirac for state vectors in quantum mechanics.
### Incidence algebras.
The Dirac notation turns out to be natural for incidence algebras. Let $`𝒦`$ be an arbitrary finite poset. Denote by $``$ its linear span
$$=span\{\mathbf{|}P\mathbf{}:P𝒦\}=\left\{\underset{P𝒦}{}c_P\mathbf{|}P\mathbf{}\right\}$$
with the coefficients taken from a field with characteristic zero.
###### Definition 1.
The incidence algebra of a poset $`𝒦`$ is the following linear span
$$\mathrm{\Omega }=\mathrm{\Omega }(𝒦)=span\{\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}:P,Q𝒦\text{and}PQ\}$$
(4)
and the product defined on the basic elements according to (4):
$$\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}\mathbf{|}R\mathbf{}\mathbf{}S\mathbf{|}=\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{}R\mathbf{}\mathbf{}S\mathbf{|}=\delta _{QR}\mathbf{|}P\mathbf{}\mathbf{}S\mathbf{|}$$
This definition of product is correct due to the transitivity of partial orders:
$$\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|},\mathbf{|}Q\mathbf{}\mathbf{}S\mathbf{|}\mathrm{\Omega }PQ\text{and}QSPS\mathbf{|}P\mathbf{}\mathbf{}S\mathbf{|}\mathrm{\Omega }$$
Incidence algebras were introduced by Rota . It was proved by Stanley that a poset $`𝒦`$ can be reconstructed from its incidence algebra $`\mathrm{\Omega }(𝒦)`$ up to a poset isomorphism.
Meanwhile, poset homomorphisms (namely, monotone mappings) induce no homomorphism of incidence algebras. However, and this is the contents of this paper, the situation drastically changes when the class of posets is restricted to simplicial complexes.
### The category $`𝔖`$ of simplicial complexes.
For the sake of self-consistency I give a brief account of the standard theory, mainly to introduce the notation. Let $`V`$ be a non-empty finite set, call the elements of $`V`$ vertices.
###### Definition 2.
A collection $`𝒦`$ of non-empty subsets of $`V`$ is called (abstract) simplicial complex with the set of vertices $`V`$ whenever
* $`vV\{v\}𝒦`$
* $`P𝒦,QVQPQ𝒦`$
Evidently, $`𝒦`$ is a poset with respect to set inclusion. The elements of $`𝒦`$ are called simplices.
###### Definition 3.
Let $`𝒦`$, $`𝒦^{}`$ be two simplicial complexes with the sets of vertices $`V`$, $`V^{}`$, respectively. A mapping $`\pi :𝒦^{}𝒦`$ is called simplicial if
* vertices are mapped on vertices: $`V^{}\pi V`$
* $`.\pi |_V^{}`$ completely determines $`\pi `$ on the whole $`𝒦^{}`$: $`P^{}𝒦^{}P^{}\pi =_{v^{}P^{}}\{v^{}\}\pi `$
Remark. Sometimes (in particular, in the rest of this paper) the notion of simplicial mapping is referred to a mapping $`\pi `$ between vertices, then the second condition reads:
$$\{v_0^{},\mathrm{},v_n^{}\}𝒦^{}\{v_0^{}\pi ,\mathrm{},v_n^{}\pi \}𝒦^{}$$
(5)
For instance, let $`𝒦^{}=\text{}`$ and $`𝒦=\text{}`$ and let $`\pi _1`$, $`\pi _2`$ be
$$\begin{array}{ccc}\text{}& & \text{}\\ \pi _1:𝒦^{}𝒦& & \pi _2:𝒦^{}𝒦\end{array}$$
then $`\pi _1`$ is simplicial and $`\pi _2`$ is not (since $`\{1^{},2^{}\}\pi _2=\{1,3\}𝒦`$).
Denote by $`𝔖`$ the category whose objects are simplicial complexes and whose arrows are simplicial mappings
$$𝔖=(\text{simplicial complexes},\text{simplicial mappings})$$
It follows immediately from the definition that simplicial mappings are monotone with respect to set inclusion and that $`𝔖`$ is (not a full) sub-category of the category $`𝔓𝔒𝔖𝔈𝔗=(\text{posets},\text{monotone mappings})`$
### Dirac notation for homological operations.
Fix an enumeration of the vertices of $`𝒦`$, then for any simplex $`P=\{v_0,\mathrm{},v_n\}𝒦`$ and any its vertex $`v_i`$ the incidence coefficient is defined
$$ϵ_{v_iP}=(1)^i$$
(6)
A face $`P_v`$ of a simplex $`P`$ is its subset $`P\{v\}`$, we write
$$P_v=Pv;P=p_v+v$$
(7)
The dimension of a simplex $`P`$ is the number of its vertices minus one:
$$dimP=\mathrm{card}P1$$
(8)
Denote by $`𝒦^n`$ the $`n`$-skeleton of $`𝒦`$ — the set of its simplices of dimension $`n`$
$$𝒦^n=\{P𝒦:dimP=n\}$$
and consider the linear spans
$$^n=span𝒦^n=\left\{\underset{P𝒦^n}{}c_P\mathbf{|}P\mathbf{}\right\}$$
###### Definition 4.
The border operator $`\delta :^n^{n1}`$ acts as
$$\delta \mathbf{|}P\mathbf{}=\underset{vP}{}ϵ_{vP}\mathbf{|}P_v\mathbf{}$$
(9)
and then extends to the space $``$ (assuming $`\delta \mathbf{|}v\mathbf{}=0,vV`$):
$$=^n=span𝒦$$
It is proved that $`\delta ^2`$ is always zero. Due to Dirac notation the same symbol $`\delta `$ is used to denote its adjoint, called coborder operator acting from $`^n`$ to $`^{n+1}`$, so
$$P𝒦^n\{\begin{array}{ccc}\hfill \mathbf{}P\mathbf{|}\delta & & ^{n+1}\hfill \\ \hfill \delta \mathbf{|}P\mathbf{}& & ^{n1}\hfill \end{array}$$
(10)
### Finite-dimensional differential moduli.
Recall the basic definitions concerning finite-dimensional analogs of moduli of exterior differential forms. Let $`𝒜`$ be a semisimple finite-dimensional commutative algebra.
###### Definition 5.
A differential module $`𝒟`$ over a basic algebra $`𝒜`$ is a triple
$$𝒟=(\mathrm{\Omega },𝒜,\text{d})$$
where $`\mathrm{\Omega }`$ is a graded algebra
$$\mathrm{\Omega }=\mathrm{\Omega }^0\mathrm{\Omega }^1\mathrm{},\mathrm{\Omega }^0=𝒜$$
(11)
equipped with the Kähler differential $`\text{d}:\mathrm{\Omega }^n\mathrm{\Omega }_{n+1}`$ such that for any $`\omega ^r\mathrm{\Omega }^r`$, $`\omega ^s\mathrm{\Omega }^s`$
$$\begin{array}{c}\text{d}^2=0\hfill \\ \text{d}(\omega ^r\omega ^s)=\text{d}\omega ^r\omega ^s+(1)^r\omega ^r\text{d}\omega ^s\hfill \end{array}$$
(12)
The second equality is called graded Leibniz rule.
### Universal differential envelope.
Given an algebra $`𝒜`$, any differential module over it can be obtained as a quotient of a universal object $`\mathrm{\Omega }_u=\mathrm{\Omega }_u(𝒜)`$, called universal differential envelope of $`𝒜`$ over appropriate differential ideal $``$. Recall the necessary definitions (for details the Reader is referred to ).
Consider the tensor product $`𝒜𝒜`$ and define the operator $`m:𝒜𝒜𝒜`$ as follows:
$$m(ab)=ab$$
Then consider its kernel
$$\mathrm{\Omega }_u^1=\mathrm{ker}m$$
define by induction
$$\mathrm{\Omega }_u^{n+1}=\mathrm{\Omega }_u^n_𝒜\mathrm{\Omega }_u^1$$
and form the sum
$$\mathrm{\Omega }_u=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_u^n$$
(13)
Define the Kähler differential $`\text{d}_u`$ first on $`\mathrm{\Omega }_u^0=𝒜`$:
$$\text{d}_ua=\mathrm{𝟏}aa\mathrm{𝟏}$$
and then extend it by induction to higher degrees using the identities (12), for instance
$$\text{d}_u(a\text{d}_ub)=\text{d}_ua\text{d}_ub$$
(14)
###### Definition 6.
The universal differential envelope $`\mathrm{\Omega }_u=\mathrm{\Omega }_u(𝒜)`$ of the algebra $`𝒜`$ is the differential module $`(\mathrm{\Omega }_u,𝒜,\text{d}_u)`$.
## 1 Differential structure of incidence algebras
In this section we introduce a particular representation for universal differential envelopes of finite-dimensional algeras called stories semantics. Let $`𝒦`$ be an arbitrary (yet structureless) finite set, denote by $`𝒜`$ the algebra of all complex-valued functions on $`𝒦`$.
### Stories semantics for universal differential envelope.
Call the elements of $`𝒦`$ statements, and consider first all possible sequences of statements of finite length. A homogeneous $`\mathrm{n}`$-story is a sequence $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}`$ whose no neighbor statements are the same. Denote by $`\mathrm{\Omega }_S^n`$ the linear span of all homogeneous $`n`$-stories:
$$\mathrm{\Omega }_S^n=span\left\{\mathbf{}P_0,\mathrm{},P_n\mathbf{}:i=1,\mathrm{},nP_{i1}P_i\right\}$$
(15)
and define the product of two stories as follows:
$$\mathbf{}P_0,\mathrm{},P_n\mathbf{}\mathbf{}Q_0,\mathrm{},Q_m\mathbf{}=\{\begin{array}{cc}\mathbf{}P_0,\mathrm{},P_nQ_1,\mathrm{},Q_m\mathbf{},\hfill & \text{if }P_n=Q_0\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(16)
which, being extended by linearity to the direct sum $`_{n=0}^{\mathrm{}}\mathrm{\Omega }_S^n`$ makes it graded algebra, call it stories algebra.
Now let us write down the explicit form of the Kähler differential. For $`\mathbf{}P\mathbf{}𝒜`$ we have:
$$\text{d}_u\mathbf{}P\mathbf{}=\underset{QP}{}(\mathbf{}QP\mathbf{}\mathbf{}PQ\mathbf{})$$
and for arbitrary $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}\mathrm{\Omega }_u^n`$
$$\begin{array}{c}\text{d}_u\mathbf{}P_0,\mathrm{},P_n\mathbf{}=\underset{Q:QP_0}{}\mathbf{}QP_0,\mathrm{},P_n\mathbf{}+\hfill \\ \hfill +\underset{k=1}{\overset{n}{}}(1)^k\underset{Q:P_{k1}QP_k}{}\mathbf{}P_0,\mathrm{},P_{k1}QP_k,\mathrm{},P_n\mathbf{}+\\ \hfill +(1)^{n+1}\underset{Q:P_nQ}{}\mathbf{}P_0,\mathrm{},P_nQ\mathbf{}\end{array}$$
(17)
###### Lemma 1.
The universal differential envelope of a finite-dimensional semisimple algebra $`𝒜`$ is isomorhic (as graded algebra) to the stories algebra:
$$\mathrm{\Omega }_u=\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Omega }_S^n$$
### Simplicial differential ideals.
Let us specify the structure of the set $`𝒦`$ assuming it to be a simplicial complex. Our goal is to find the explicit form of the differential ideal $``$ giving rise to its incidence algebra $`\mathrm{\Omega }(𝒦)`$. To gather it, first let us classify the stories in a way taking into account that $`𝒦`$ is a simplicial complex.
###### Definition 7.
A story $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}`$ is called fair whenever $`P_{i1}=P_iv`$ — see (7). Otherwise the story $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}`$ is unfair.
Form the linear span $`_N=_{n=0}^{\mathrm{}}_N^n`$, where
$$_N^n=span\{\text{unfair }n\text{-stories}\}=span\{\mathbf{}P_0,\mathrm{},P_n\mathbf{}:\neg (iP_{i1}=P_iv)\}$$
(18)
Fix an enumeration of vertices of $`𝒦`$ and with each fair story $`w`$ associate a number $`ϵ_w=\pm 1`$ as follows:
$$ϵ_w=\underset{i=1}{\overset{n}{}}ϵ_{v_iP_i}$$
where $`ϵ_{v_iP_i}`$ is the incidence coefficient (6). Let $`P_0`$, $`P_n`$ be two simplices such that $`P_0P_n`$, and $`dimP_0dimP_n`$ is exactly $`n`$. For any two fair $`n`$-stories $`w=\mathbf{}P_0P_1,\mathrm{},P_{n1}P_n\mathbf{}`$ and $`w^{}=\mathbf{}P_0P_1^{},\mathrm{},P_{n1}^{}P_n\mathbf{}`$ having common initial and final statements form the difference $`ϵ_wwϵ_w^{}w^{}`$, and consider the linear hull of all such differences
$$_S^n=\underset{w}{span}\{ϵ_wwϵ_w^{}w^{}:w,w^{}\text{as described above}\}$$
Take for each $`n`$ the sum
$$^n=_N^n_S^n$$
and form the following graded linear space
$$=\underset{n=1}{\overset{\mathrm{}}{}}^n$$
(19)
###### Lemma 2.
$``$ is a differential ideal in $`\mathrm{\Omega }_u`$.
###### Proof.
Evidently $`_N=_n_N^n`$ is an ideal in $`\mathrm{\Omega }_u`$. Besides that, a product of an element of $`_S^n`$ and a fair $`m`$-story is always in $`_S^{m+n}`$, therefore $``$ is an ideal. It remains to prove that the ideal $``$ is differential: $`\text{d}_u(^n)^{n+1}`$.
First let $`w_N^n`$, consider $`\text{d}_uw`$ as the sum (17). The 1st and the 3rd summands of (17) are always in $`_N^{n+1}`$, the same for any term from the middle sum of (17) with the only possible exception when $`w=\mathbf{}P_0,\mathrm{},P_iP_{i+1},\mathrm{},P_n\mathbf{}`$ such that both $`\mathbf{}P_0,\mathrm{},P_i\mathbf{}`$ and $`\mathbf{}P_{i+1},\mathrm{},P_n\mathbf{}`$ are fair stories, while $`P_i=P_{i1}uv`$. Then
$$\text{d}_uw=\nu +(1)^{i+1}(\rho +\rho ^{})$$
where $`\nu `$ is the sum of elements of (17) from $`_N^{n+1}`$, and $`\rho `$, $`\rho ^{}`$ are the following fair stories:
$$\begin{array}{ccc}\hfill \rho & =& \mathbf{}P_0,\mathrm{},P_i,P_{i+1}u,P_{i+1},\mathrm{},P_n\mathbf{}\hfill \\ \hfill \rho ^{}& =& \mathbf{}P_0,\mathrm{},P_i,P_{i+1}v,P_{i+1},\mathrm{},P_n\mathbf{}\hfill \end{array}$$
for which $`ϵ_\rho =ϵ_\rho ^{}`$, therefore $`\rho +\rho ^{}_S^{n+1}`$ and
$$\text{d}_u_N^n_N^{n+1}_S^{n+1}$$
Now let $`w=ϵ_\rho \rho ϵ_\rho ^{}\rho ^{}_S^n`$. That is,
$$\rho =\mathbf{}P_0,P_1,\mathrm{},P_{n1},P_n\mathbf{},\rho ^{}=\mathbf{}P_0,P_1^{},\mathrm{},P_{n1}^{},P_n\mathbf{}$$
Consider the three sums (17) for $`\text{d}_uw`$. All the terms in the 2nd sum (17) will be in $`_N^{n+1}`$. The terms from the first sum will also belong to $`_N^{n+1}`$ with the only exception — the terms of the form $`k=ϵ_\tau \tau ϵ_\tau ^{}\tau ^{}`$, where
$$\tau =\mathbf{}P_0v,P_0\mathbf{}\rho ,\tau ^{}=\mathbf{}P_0v,P_0\mathbf{}\rho ^{}$$
but for them $`ϵ_\tau =ϵ_{vP_0}ϵ_\rho `$ and $`ϵ_\tau ^{}=ϵ_{vP_0}ϵ_\rho ^{}`$, therefore $`ϵ_\tau ϵ_\tau ^{}=ϵ_\rho ϵ_\rho ^{}`$ and $`k_S^{n+1}`$, so $`\text{d}_u_S^n_N^n_S^n`$. This completes the proof:
$$\text{d}_u(_N^n_S^n)_N^{n+1}_S^{n+1}$$
###### Theorem 3.
The quotient $`\mathrm{\Omega }_u/`$ and the incidence algebra $`\mathrm{\Omega }(𝒦)`$ are isomorphic algebras.
###### Proof.
Consider the mapping $`\sigma :\mathrm{\Omega }_u\mathrm{\Omega }`$ defined on any story $`w=\mathbf{}P_0,\mathrm{},P_n\mathbf{}\mathrm{\Omega }_u^n`$ as
$$\sigma (w)=\{\begin{array}{cc}ϵ_w\mathbf{|}P_0\mathbf{}\mathbf{}P_n\mathbf{|},\hfill & \text{if }w\text{ is a fair story}\hfill \\ 0\hfill & \text{otherwise}\hfill \end{array}$$
(20)
then $`\mathrm{ker}\sigma =_N_S=`$, so $`\mathrm{\Omega }\mathrm{\Omega }_u/`$ as graded linear spaces. To verify that $`\sigma `$ preserves products, let $`w=\mathbf{}P_0,\mathrm{},P_m\mathbf{}`$ and $`w^{}=\mathbf{}Q_0,\mathrm{},Q_n\mathbf{}`$ be two fair stories. If $`P_mQ_0`$ everything is trivial. Suppose $`P_m=Q_0`$, then $`ww^{}=\mathbf{}P_0,\mathrm{},P_mQ_1,\mathrm{},Q_n\mathbf{}`$ and
$$\sigma (ww^{})=ϵ_{(ww^{})}\mathbf{|}P_0\mathbf{}\mathbf{}Q_n\mathbf{|}=ϵ_wϵ_w^{}\mathbf{|}P_0\mathbf{}\mathbf{}P_m\mathbf{|}\mathbf{|}Q_0\mathbf{}\mathbf{}Q_n\mathbf{|}=\sigma (w)\sigma (w^{})$$
### The induced differential structure on $`\mathrm{\Omega }(𝒦)`$.
Having a projection of the universal differential envelope $`\mathrm{\Omega }_u`$ onto the incidence algebra $`\mathrm{\Omega }(𝒦)`$ we first conclude that $`\mathrm{\Omega }(𝒦)`$ necessarily has the structure of a differential module over $`𝒜`$, namely, that induced by the projection $`\sigma `$ onto the quotient. Now, having the expression (20) for $`\sigma `$, let us explicitly calculate the form of the differential d on $`\mathrm{\Omega }(𝒦)`$ according to the formula
$$\text{d}=\sigma ^1\text{d}_u\sigma $$
(21)
###### Theorem 4.
The differential in the incidence algebra $`\mathrm{\Omega }(𝒦)`$ of a simplicial complex $`𝒦`$ has the following form. Let $`\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}\mathrm{\Omega }^n`$, then
$$\text{d}\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}=\mathbf{|}\delta P\mathbf{}\mathbf{}Q\mathbf{|}(1)^n\mathbf{|}P\mathbf{}\mathbf{}Q\delta \mathbf{|}$$
(22)
where $`\delta `$ is the symbol (10) for both border and coborder operations.
###### Proof.
Begin with $`\mathrm{\Omega }^0=𝒜`$. In this case $`\sigma ^1=.\text{id}|_𝒜`$, therefore
$$\text{d}\mathbf{|}P\mathbf{}\mathbf{}P\mathbf{|}=\sigma (\text{d}_u\mathbf{}P\mathbf{})=\underset{vP}{}ϵ_{vP}\mathbf{|}Pv\mathbf{}\mathbf{}P\mathbf{|}\underset{u:P+u𝒦}{}ϵ_{u(P+u)}\mathbf{|}P\mathbf{}\mathbf{}P+u\mathbf{|}$$
According to (9), the first term is $`\mathbf{|}\delta P\mathbf{}\mathbf{}P\mathbf{|}`$ while the second is $`\mathbf{|}P\mathbf{}\mathbf{}P\delta \mathbf{|}`$, so
$$\text{d}\mathbf{|}P\mathbf{}\mathbf{}P\mathbf{|}=\mathbf{|}\delta P\mathbf{}\mathbf{}P\mathbf{|}\mathbf{|}P\mathbf{}\mathbf{}P\delta \mathbf{|}$$
Now let $`\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}\mathrm{\Omega }^1`$, then $`P=Qv`$ for some $`vV`$ and
$$\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}=\mathbf{|}P\mathbf{}\mathbf{}P\mathbf{|}ϵ_{vQ}\text{d}(\mathbf{|}Q\mathbf{}\mathbf{}Q\mathbf{|})$$
therefore it follows from (14) $`ϵ_{vQ}=\mathbf{}P\mathbf{|}\delta \mathbf{|}Q\mathbf{}`$ that
$$\begin{array}{c}\text{d}\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}=\text{d}\mathbf{|}P\mathbf{}\mathbf{}P\mathbf{|}ϵ_{vQ}\text{d}(\mathbf{|}Q\mathbf{}\mathbf{}Q\mathbf{|})=(\mathbf{|}\delta P\mathbf{}\mathbf{}P\mathbf{|}\mathbf{|}P\mathbf{}\mathbf{}P\delta \mathbf{|})ϵ_{vQ}(\mathbf{|}\delta Q\mathbf{}\mathbf{}Q\mathbf{|}\mathbf{|}Q\mathbf{}\mathbf{}Q\delta \mathbf{|})=\hfill \\ \hfill ϵ_{vQ}(\mathbf{|}\delta P\mathbf{}\mathbf{}P\mathbf{|}\delta \mathbf{|}Q\mathbf{}\mathbf{}Q\mathbf{|}\mathbf{|}Q\mathbf{}\mathbf{}P\mathbf{|}\delta \mathbf{|}Q\mathbf{}\mathbf{}Q\delta \mathbf{|})=\mathbf{|}\delta P\mathbf{}\mathbf{}Q\mathbf{|}+\mathbf{|}P\mathbf{}\mathbf{}Q\delta \mathbf{|}\end{array}$$
For higher degrees the formula (22) is proved by induction. First note that d enjoys the Leibniz rule as both $`\sigma `$, $`\sigma ^{}`$ in (21) preserve products. Then represent any $`\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}\mathrm{\Omega }^n`$ as a product
$$\mathbf{|}P\mathbf{}\mathbf{}Q\mathbf{|}=\mathbf{|}P\mathbf{}\mathbf{}Qv\mathbf{|}\mathbf{|}Qv\mathbf{}\mathbf{}Q\mathbf{|}\mathrm{\Omega }^{n1}\mathrm{\Omega }^n$$
and carry out a routine calculation.
### Summary of this section.
It was established that the incidence algebra $`\mathrm{\Omega }(𝒦)`$ of any simplicial complex $`𝒦`$ is a differential module over the algebra $`𝒜`$ of all functions on $`𝒦`$ (see Definition 5). The algebra $`\mathrm{\Omega }(𝒦)`$ was represented as a quotient of the universal differential envelope $`\mathrm{\Omega }_u(𝒜)`$ over the simplicial differential ideal $``$ (19). The form of the Kähler differential for $`\mathrm{\Omega }(𝒦)`$ is given in (22).
## 2 Functorial properties
As it was already mentioned, an arbitrary monotone mapping $`\pi :𝒦^{}𝒦`$ between two posets $`𝒦`$ and $`𝒦^{}`$ produces no homomorphism of their incidence algebras. However, if we consider the category $`𝔖`$ of simplicial complexes, the situation becomes completely different. In this section I show that that the correspondence $`𝒦\mathrm{\Omega }(𝒦)`$ is a contravariant functor from the category $`𝔖`$ to the category $`𝔇𝔐`$ of differential moduli over commutative algebras.
### The category $`𝔇𝔐`$.
The objects of the category $`𝔇𝔐`$ are differential moduli over semisimple commutative algebras — see Definition 5. The morphisms of $`𝔇𝔐`$ are differentiable mappings.
###### Definition 8.
Let $`𝒟=(\mathrm{\Omega },𝒜,\text{d})`$, $`𝒟^{}=(\mathrm{\Omega }^{},𝒜^{},\text{d}^{})`$ be two differential moduli. A mapping $`\varphi :\mathrm{\Omega }\mathrm{\Omega }^{}`$ is called differentiable iff
$$\begin{array}{c}\varphi \text{is a homomorphism of graded algebras}\hfill \\ \varphi \text{d}^{}=\text{d}\varphi \hfill \end{array}$$
(23)
###### Lemma 5.
Any differentiable mapping is completely defined by its values on the basic algebra $`𝒜`$.
###### Proof.
By induction, let $`w\mathrm{\Omega }^1`$, then $`w=_ia_i\text{d}b_i`$ with $`a_i,b_i𝒜`$. Then $`\varphi (w)=\varphi (a_i)\varphi (\text{d}b_i)=\varphi (a_i)\text{d}^{}\varphi (b_i)`$. When $`w^{n+1}\mathrm{\Omega }^{n+1}`$, it reads $`w^{n+1}=_ia_i\text{d}w_i^n`$ with $`a_i𝒜`$ and $`w_i^n\mathrm{\Omega }^n`$. Applying of $`\varphi `$ and using the second property (23) completes the proof. ∎
So, the notion of differentiable mappings can be referred to homomorphisms $`\varphi :𝒜𝒜^{}`$ between basic algebras. Let $`𝒟=(\mathrm{\Omega },𝒜,\text{d})`$, $`𝒟^{}=(\mathrm{\Omega }^{},𝒜^{},\text{d}^{})`$ be two differential moduli. A mapping $`\varphi :𝒜𝒜^{}`$ completely determines both the homomorhisms $`\varphi :\mathrm{\Omega }\mathrm{\Omega }`$ and $`\overline{\varphi }:\mathrm{\Omega }_u\mathrm{\Omega }_u`$. Then $`\varphi `$ is differentiable if and only if the following holds:
$$\overline{\varphi }()^{}$$
(24)
Let the basic algebras $`𝒜,𝒜^{}`$ are represented by functions on the sets $`𝒦`$, $`𝒦^{}`$, respectively. Then the dual mapping $`\pi :𝒦^{}𝒦`$ determines $`\varphi `$ completely, and the explicit form of the mapping $`\overline{\varphi }`$ is the following: for any $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}\mathrm{\Omega }_u^n`$
$$\varphi \mathbf{}P_0,\mathrm{},P_n\mathbf{}=\underset{P_i^{}:iP_i^{}\pi =p_i}{}\mathbf{}P_0^{},\mathrm{},P_n^{}\mathbf{}$$
(25)
### Functorial properties.
To prove that the correspondence $`𝒦\mathrm{\Omega }(𝒦)`$ is a functor we have to provide a correspondence between the arrows of the categories. As stated above, any (set-theoretical) mapping $`\pi :𝒦^{}𝒦`$ gives rise to a homomorphism (25) $`\varphi =\pi ^{}`$ of the basic algebras.
###### Theorem 6.
Let $`𝒦`$, $`𝒦^{}`$ be two simplicial complexes, and let a mapping $`\pi :𝒦^{}𝒦`$ be simplicial. Then its dual $`\varphi =\pi ^{}`$ is differentiable.
###### Proof.
By virtue of (24) it suffices to prove that simplicial ideals are mapped into simplicial ideals in the target algebra. Let $`W^{}=\mathbf{}P_0^{},\mathrm{},P_n^{}\mathbf{}`$ be a fair story. Let us prove that its image
$$W^{}\pi =\mathbf{}P_0^{}\pi ,\mathrm{},P_n^{}\pi \mathbf{}$$
being a sequence of elements of $`𝒦`$ is either a fair story or not a story. $`P_i^{}=P_{i1}^{}+v_i`$ for any $`i=1,\mathrm{},n`$, and $`ijv_iv_j`$. Denote $`P_i=P_i^{}\pi `$ and $`v_i=v_i^{}\pi `$ (recall that $`V^{}\pi V`$). For each $`i`$ we have exactly two possibilities: either $`v_iP_{i1}`$ or $`v_iP_{i1}`$ (and then according to (5) $`P_i=P_{i1}+v_i𝒦`$). If $`v_iP_i`$ for some for some $`i`$ then $`W^{}\pi `$ contains $`P_i=P_{i1}`$ (and therefore is not a story) otherwise $`iP_i=P_{i1}+v_i`$, so $`W`$ is a fair story. So,
$$\{\text{fair stories}\}\pi \{\text{fair stories}\}$$
(26)
Return to simplicial ideals. Let $`\mathbf{}P_0,\mathrm{},P_n\mathbf{}`$ be an unfair story, then $`\varphi (\mathbf{}P_0,\mathrm{},P_n\mathbf{})`$ is according to (25) a sum of unfair stories (otherwise (26) is violated).
Now let $`i=ϵ_wϵ_{\stackrel{~}{w}}_S`$ such that $`P_n=P_0+v_1+\mathrm{}+v_n`$ in $`w`$ and $`P_n=P_0+v_{j_1}+\mathrm{}+v_{j_n}`$ in $`\stackrel{~}{w}`$. Let $`w^{}=P_0^{},\mathrm{},P_n^{}`$ be a fair story in $`𝒦^{}`$ such that $`w^{}\pi =w`$, then necessarily
* $`P_n^{}=P_0^{}+v_1^{}+\mathrm{}+v_n^{}`$
* all $`v_j^{}`$ are disjoint
* $`jv_j^{}\pi =v_j`$ — there is 1–1 correspondence
Make a new fair story $`\stackrel{~}{w}^{}`$ from $`w^{}`$ performing the same permutation of vertices as that making $`\stackrel{~}{w}`$ from $`w`$, and we obtain for it $`\stackrel{~}{w}^{}\pi =\stackrel{~}{w}`$ and $`ϵ_{\stackrel{~}{w}^{}}=ϵ_w^{}`$. So, any fair preimage of $`i_S`$ will be in $`_S^{}`$
### Summary of this section.
The correspondence $`𝒦\mathrm{\Omega }(𝒦)`$ is proved to be a contravariant functor from the category $`𝔖=\{\text{simplicial complexes},\text{simplicial mappings}\}`$ into the category $`𝔇𝔐`$ of differential moduli over finite-dimensional semisimple commutative algebras.
## 3 Concluding remarks
It was shown that the incidence algebras of simplicial complexes possess the structure of differential moduli, and it was established that, contrary to posets of general form, the correspondence between simplicial complexes and their incidence algebras is a contravariant functor.
It occurs that simplicial complexes possess the natural structure of discrete differential manifolds — finite sets equiped with differential calculi — see, e.g. for a review.
Aside of purely mathematical context, the presented results have physical application being a basis for discrete approximations of spacetime structure . When simplicial complexes are treated as coarse-grained spacetime patterns, our results enable the possibility of linking more and more refined approximations between each other.
### Ackowledgments.
The author appreciates numerous remarks made by the participants of Friedmann Seminar on theoretical physics (St. Petersburg, Russia) headed by A.A. Grib. Deep and profound discussiond with my Italian colleagues G. Landi, F. Lizzi and the members of the Theory Group of the INFN Section in Naples were helpful. Dr. I. Raptis (Pretoria) made his valuable impact explaining the physical meaning of the construction.
A financial support from the state research program ‘Universities of Russia’ is appreciated. |
warning/0001/cond-mat0001401.html | ar5iv | text | # s+d pairing in orthorhombic phase of copper-oxides
## Abstract
A microscopical theory of electronic spectrum and superconductivity is formulated within the two-dimensional anisotropic $`tJ`$ model with $`t_xt_y`$ and $`J_xJ_y`$. Renormalization of electronic spectrum and superconductivity mediated by spin-fluctuations are investigated within the Eliashberg equation in the weak coupling approximation. The gap function has $`d+s`$ symmetry with the extended $`s`$-wave component being proportional to the asymmetry $`t_yt_x`$. Some experimental consequences of the obtained results are discussed.
Recently the $`d`$-wave symmetry of superconducting pairing in cuprates was unambiguously confirmed by observation of half-integer magnetic flux quantum . In a tetragonal phase of copper oxides the $`s`$-wave component must be strongly suppressed due to on-site Coulomb correlations. For the 2D $`t`$-$`J`$ model it follows from the constraint of no double occupancy on a single site given by the identity :
$$\widehat{a}_{i,\sigma }\widehat{a}_{i,\sigma }=\frac{1}{N}\underset{k_x,k_y}{}\widehat{a}_{𝐤,\sigma }\widehat{a}_{𝐤,\sigma }=0.$$
(1)
for the projected electron operators $`\widehat{a}_{i,\sigma }=a_{i,\sigma }(1n_{i,\sigma })`$. The anomalous correlation function $`\widehat{a}_{k,\sigma }\widehat{a}_{k,\sigma }`$ is proportional to the gap function, $`\mathrm{\Delta }(k_x,k_y)`$, multiplied by a positive function symmetric in respect to the $`D_{4h}`$ point group which defines the symmetry of the Fermi surface (FS). Therefore to satisfy the condition (1) the gap function should have a lower symmetry, e.g. $`B_{1g}`$, ”$`d`$-wave” symmetry (see, e.g. ): $`\mathrm{\Delta }_d(k_x,k_y)=\mathrm{\Delta }_d(k_y,k_x).`$
In the orthorhombic phase the FS has a lower symmetry, e.g., $`D_{2h}`$, and the condition (1) can be fulfilled for a gap function of the same symmetry (within the $`E_{1g}`$ irreducible representation) which can be written in a general form (”$`d+s`$”):
$$\mathrm{\Delta }(k_x,k_y)=\mathrm{\Delta }_d(k_x,k_y)+ϵ\mathrm{\Delta }_s(k_x,k_y).$$
(2)
where $`\mathrm{\Delta }_s(k_x,k_y)=\mathrm{\Delta }_s(k_y,k_x)`$ is ”the extended $`s`$-wave” component.
In the present paper we calculate superconducting $`T_c`$ for the 2D $`t`$-$`t^{}`$-$`J`$ model within the theory developed by us in , both in tetragonal and orthorhombic phases. The orthorhombic ($`D_{2h}`$) distortion is taken into account by introducing the asymmetric hopping parameters $`t_{ij}`$ and the exchange interaction $`J_{ij}`$ for the nearest neighbors (n.n.) in the form: $`t_{x/y}=t(1\pm \alpha )`$, $`J_{x/y}=J(1\pm \beta )`$ where the asymmetry parameters are supposed to be small quantities: $`\alpha \beta 0.1`$. For the next n.n., $`t_{ij}=t^{}.`$
The Dyson equations for the matrix Green function (GF) in the Nambu notation was obtained by the equation of motion method for the Hubbard operators as described in . For estimation of the role of orthorhombic deformation we consider here only the weak-coupling approximation. However, to take into account strong electronic correlations in the $`t`$-$`J`$ model due to restricted hopping in the singly occupied subband we write the single-electron spectral density in the form: $`A(k,\omega )Z_k\delta (\omega +\mu E_k)+A_{inc}(\omega )`$. The quasiparticle weight $`Z_k`$ and the incoherent part $`A_{inc}(\omega )`$ are coupled by the sum rule for the spectral density: $`_{\mathrm{}}^+\mathrm{}𝑑\omega A(k,\omega )=1n/2.`$ To fix the value of $`Z_k`$ we assume that the FS for quasiparticles with spectrum $`E_k`$ obeys the Luttinger theorem: the average number of electrons is equal to the number of states in $`k`$-space below the chemical potential $`\mu `$: $`n=(1/N)_{k,\sigma }\{\mathrm{exp}[E_k\mu )/T]+1\}^1.`$ ¿From these conditions we have estimations: $`Z(1n)/(1n/2)`$ and $`A_{inc}(\omega )(n/2)^2/(1n/2)(W\mathrm{\Gamma })`$ where we have suggested that the coherent band lies in the range $`\mathrm{\Gamma }\omega \mathrm{\Gamma }`$ while the incoherent band lies below the coherent band in the range $`W\omega \mathrm{\Gamma }`$. By taking into account the renormalization of the coherent part of the spectral weight by $`Z_k`$ we write the equation for the gap in the weak coupling approximation in the form:
$$\mathrm{\Delta }_k=\frac{1}{N}\underset{q}{}K(q,kq)\frac{Z_q^2\mathrm{\Delta }_q}{2\mathrm{\Omega }_q}\mathrm{tanh}\frac{\mathrm{\Omega }_q}{2T}.$$
(3)
where $`\mathrm{\Omega }_q=[(E_q\mu )^2+\mathrm{\Delta }_q^2]^{1/2}`$ and $`E_qt_{eff}[\gamma (q)+\alpha \eta (q)]t_{eff}^{^{}}\gamma ^{}(k)`$ with renormalized due to strong correlations hopping parameters and $`\gamma (q)=(1/2)(\mathrm{cos}q_x+\mathrm{cos}q_y),\eta (q)=(1/2)(\mathrm{cos}q_x\mathrm{cos}q_y),\gamma ^{}(q)=\mathrm{cos}q_x\mathrm{cos}q_y`$ . $`K(q,kq)=\{2g(q,kq)\lambda (q,kq)\}`$ with the vertex $`g(q,kq)=t(q)J(kq)/2`$ and $`\lambda (q,kq)=g{}_{}{}^{2}(q,kq)\chi (kq)`$. The first term in the vertex, $`t(q)`$, is due to the kinematical interaction caused by constraints and the second one, $`J(q)`$, is the exchange coupling. They have different $`q`$-dependence and are effective at different doping. The spin-fluctuation coupling in $`\lambda (q,kq)`$ is defined by the static spin susceptibility $`\chi (kq)`$ for which we used the model $`\chi (q)=\chi _0/[1+\xi ^2(1+\gamma (q))]`$ where the antiferromagnetic (AFM) correlation length $`\xi `$ is a fitting parameter while $`\chi _0`$ is normalized by the condition: $`1/N_i𝐒_i𝐒_i=(3/4)n`$.
We performed numerical solution of Eq. (3) for the gap in the form (2) with $`\mathrm{\Delta }_d(k_x,k_y)=\mathrm{\Delta }\eta (k)`$ and $`\mathrm{\Delta }_s(k_x,k_y)=\mathrm{\Delta }\gamma (k)`$. By taking into account the constraint of no double occupancy, Eq.(1), we estimate the weight $`ϵ`$ of the $`s`$-component. The critical temperature $`T_c(\delta )`$ (in units of $`t`$) is shown on Fig.1 in the tetragonal, $`\alpha =0.0`$, (bold line) and orthorhombic, $`\alpha =0.1`$, (dashed line) phases for $`J=0.4t,\xi =2,t^{}=0.0`$. Suppression of $`T_c`$ in the orthorhombic phase is due to a deformation of the FS resulting in a less favourable electron pairing by the AFM spin fluctuations. Increasing of AFM interaction due to larger $`J`$ or/and $`\xi `$ strongly enhances $`T_c`$ though does not change the shape of the curve. Its maximum at $`\delta 0.35`$ is due to an interplay between the shape of the FS (defined by the quasiparticle spectrum $`E_q`$) and the coherent spectral weight $`Z^2`$ in Eq.(3). Particularly, for $`t^{}/t=0.1(+0.1)`$ we observed a strong enhancement (suppression) of $`T_c(\delta )`$ due to change of the FS. In the orthorhombic phase the $`s`$-wave component with $`ϵ1`$ (depending on the doping) appears in Eq.(2) shifting 4 nodes of the gap at the FS from the diagonals $`k_x=\pm k_y`$ in a tetragonal phase as in Ref. .
To conclude, in the present paper we estimate the role of orthorhombic deformation in cuprates within the $`t`$-$`t^{}`$-$`J`$ model by solving the equation (3) for a gap of general symmetry (2) with constraint of no double occupancy, Eq.(1). The obtained dependence $`T_c(\delta )`$ in Fig.1 can explain a suppression of $`T_c`$ in the orthorhombic phase that observed in LSCO and anisotropic pressure dependence of $`T_c`$ in YBCO . Quite a large ”$`s`$”-wave component in the gap (2)is in accord with a nonzero tunnelling along the $`c`$-axis in YBCO . |
warning/0001/hep-ex0001012.html | ar5iv | text | # 1 Introduction
## 1 Introduction
My assignment at this conference is to assess where we are in high-energy physics and speculate on where we might be going. This frees me from any obligation to summarize all that went on here and allows me to talk just about those topics that interest me the most at this time. I will let my prejudices show and talk a bit about physics in general, CP violation, neutrinos, accelerators, non-accelerator experiments, and even theory.
I start with a bit of my own history in physics and what I see as cycles in science—I think these are relevant. I began research in physics when I was in my third under-graduate year at MIT. I worked with Professor Francis Bitter in his magnet laboratory and learned how to do an experiment. I spent three years with Francis Bitter, but, as I began my second year of graduate school, I found myself becoming less interested in the nuclear magnetic moments that I was trying to measure and more interested in the fundamentals of the protons and neutrons that contributed to the moments. I shifted to particle physics, never regretting it.
When in graduate school I did both theory and experiment. It was not hard to do that back then because, if you could manipulate $`\gamma `$ matrices, set up and solve integral equations and understood what a Green’s function was, you could keep up with theory and even do some of your own. I think that an experimenter in graduate school today would not have an easy time keeping up with the sophisticated mathematics required to understand what’s going on in string theory (but, I think it would be worth a try).
In the 1950s, virtually every major university had an accelerator of its own, including MIT, Harvard, Yale, Columbia, Cornell, Princeton, Chicago, the University of Illinois, Michigan, Minnesota, Berkeley, Caltech, etc. The really big facilities were the 6-GeV Bevatron at Berkeley, and the 3-GeV Cosmotron at the Brookhaven National Laboratory. I learned about accelerators at a time when the experimenters and the accelerator builders were the same people. When I was a graduate student, four of us maintained, modified and ran the MIT synchrotron. Accelerators have become much more sophisticated since then and our field has become more specialized. The accelerator builders have split off from the experimenters just as the experimenters split off from the theorists a hundred years ago.
Particle physics as practiced has changed enormously. Small experiments of two or three people that were the norm when I started out have given way to huge collaborations culminating in the ATLAS and CMS collaborations at the LHC, which have 1500 to 2000 collaborators each. These collaborations are larger than many of our laboratories.
Accelerators are huge and are built by specialists. There are very few of them. The theory has become so mathematically sophisticated that it is very difficult to keep up. An experimental physicist at the height of his or her powers (that means as a graduate student or early post doc) would have a difficult time with topology, knots, branes, and all of today’s machinery. However, the questions really remain the same as they were when I was a graduate student: what are the fundamental entities, what determines their properties, what governs their interactions, how did it begin and how is it going to end?
Our science, like all sciences, repeatedly goes through a three-part cycle. In one phase, experiment leads theory. The discoveries come thick and fast, and no model exists to absorb and interpret them. That was the situation in the 1950s when meson and baryon resonances were proliferating and it became clear that these could not be the fundamental entities from which all matter was built because there were too many of them.
The second and perhaps most enjoyable phase is the time when experiment and theory advance rapidly together. The 1960s and 1970s were like that, when the experimenters and theorists played a kind of ping-pong game when an experiment would lead to a new theoretical interpretation which led to a new experimental test which led to a new modification of the theory. In quick order over a period of ten years, the standard model rose from deep inelastic scattering, scaling, quarks, neutral currents, $`\psi /J`$, the third generation, the GIM mechanism, renormalizable gauge theories, etc.
Now we are in a stressful and frustrating phase. We have an obviously incomplete model and no experiment definitively points the way toward the next step. The standard model stands uncontradicted, yet we know perfectly well that it is wrong: it has too many arbitrary constants, 18 of them if you do not include the seven more that may come from neutrino masses and a neutrino-mixing matrix. It is interesting to note that there are more arbitrary constants in the standard model now than there were resonances in the particle data book when it became the accepted view that there were too many resonances for them to be fundamental. There are other problems as well: there is not enough CP violation in the standard model to create the baryon asymmetry of the universe and there are potential problems with longitudinal $`W`$ scattering.
We are at a point in high-energy physics where we have a wonderful model with its quarks, leptons, force particles, and its $`SU(3)\times SU(2)\times U(1)`$ structure. All that is currently accessible to experiment has been correctly predicted by the theory (except for the 18 to 25 constants, of course). Yet there are things that the standard model cannot deal with and, at all our conferences in recent years, we have hoped to hear the experiments and theoretical ideas that will lead us beyond it.
## 2 CP Violation
Many of the advances in high-energy physics, indeed in all of science, come from the overthrow of unexamined assumptions. Such is the history of CP violation. The first assumption to go was that of parity conservation, a theoretical simplification which had been experimentally shown to be correct for the electromagnetic and nuclear forces. It was then assumed that it was correct for the weak interaction as well.
In the early 1950s, high-energy physics was struggling with the $`\theta `$-$`\tau `$ paradox. Experiments had turned up what seemed to be two spin-zero particles, one of which decayed into two $`\pi `$-mesons while the other decayed into three $`\pi `$-mesons. Both appeared to have the same mass to a precision of a few tenths of a percent. Lee and Yang, in their famous paper , analyzed what we really knew about parity conservation in the weak interactions and concluded that there were no constraints. They gave several examples of experimental tests that could settle the question, and C. S. Wu did one of those, studying nuclear beta decay and finding that parity was not conserved. Very soon thereafter experiments at the Columbia and Chicago cyclotrons found that parity was also not conserved in the $`\pi `$-$`\mu `$-electron decay chain. Fitch and Cronin in their experiment of 1966 found that CP was not conserved in kaon decay, and ever since we have been trying to understand what is going on.
So far, all we know about CP violation comes from the study of the $`K`$-meson system. That should change next year when the asymmetric $`B`$-Factories at SLAC and KEK produce enough data for significant measurement of CP violation in the $`B`$-meson system. Both machines and detectors have started up very well and, at the time of this writing (November 1999), PEP-II has reached a luminosity of $`1.4\times 10^{33}\mathrm{cm}^2s^1`$, about 40% of design. KEK-B is not far behind. It is reasonable to expect that roughly 10 $`fb^1`$ of data will be accumulated by the time of the year 2000 Rochester Conference in Kyoto, which should be enough to achieve an error of about $`\pm 0.15`$ in sin$`2\beta `$ from the $`\psi `$-$`K_s`$ channel.
There are new results from studies of the $`K`$-meson system that address the existence of CP violation outside the CKM matrix. Both CERN (NA-48) and FNAL (K-TeV) groups have new results measuring $`ϵ^{}/ϵ`$. These results are given together with their previous results in Table 1. The weighted average of the results has a $`\chi ^2`$ per degree of freedom of 2.8, which gives a confidence level of less than 3%, an improbable result. The particle data group in such cases scales the errors to get the appropriate $`\chi ^2`$ per degree of freedom, a procedure that some may object to. I do it anyway and, with the error scaled, $`ϵ^{}/ϵ`$ is still not consistent with superweak interactions.
The low confidence level of these results illustrates a problem that we all should think carefully about; are we handling errors properly? The two FNAL experiments, for example, have a probability of only one in 300 of coming from a random set with gaussianly distributed errors. The Fermilab group is led by the most careful physicist I know, Bruce Winstein, and the low probability may simply be bad luck. There is, however, another possibility and that is that we don’t really fully understand the error distribution functions.
All of us in high-energy physics are guilty of treating a collection of systematic errors as if they were random gaussianly distributed errors, a procedure that we know is wrong. However, we don’t know how to do it any better. We also know that even if an individual error is random and gaussianly distributed, ratios of sums and differences of such quantities may not be gaussianly distributed. There are also non-gaussian tails on acceptance functions, tracking functions, etc. We need to understand this problem better but, until we do, complex experiments with many variables, complicated triggers, and many cuts in the analysis process, should perhaps be treated with some caution. Treating errors in these complicated experiments as if they were gaussian may lead us to ascribe a much higher confidence level to a conclusion than is really deserved.
## 3 Neutrinos
Certainly one of the most exciting areas of research at present is neutrino physics. It is fair to say that the results of the last decade on the neutrinos from the sun, from the atmospheric interaction of cosmic rays, and from accelerators, are changing our thinking and challenging the standard model. There are new data from Super Kamiokande (Super-K) on solar neutrinos; new data from Super-K and others on atmospheric neutrinos; and still a problem with the Los Alamos experiment (LSND) which doesn’t seem to fit well with our current prejudices.
Before discussing the present situation, I want to mention two people who are not currently involved in the program, but who played an absolutely critical role in the evolution of neutrino physics. The first of these is Dr. Raymond Davis, Jr., now retired, of Brookhaven National Laboratory. Ray Davis had an idea that at first seemed impossible, to detect neutrinos from the sun as a way to find out about the nuclear physics of the solar cycle. As far as I can tell, the original idea of using the chlorine-argon inverse beta-decay reaction to detect neutrinos goes back to a paper by Bruno Pontecorvo written in 1946. In that paper, Pontecorvo discussed ways to detect neutrinos, including neutrinos from the sun. He dismissed the solar idea because the flux would be too small for a one-cubic-meter detector which was the largest that he could think of.
Ray Davis thought much bigger. In 1955, he began working at Brookhaven with a 1000-gallon chlorine detector. He found no events because of the small size of the detector but, with this apparatus, he perfected the argon-chlorine radio-chemical separation techniques that allowed his later experiments to succeed. In the early 1960s, he began construction of a 100,000-gallon detector at the Homestake Mine in South Dakota and began collecting data in 1967. From the very beginning his results indicated that the flux of neutrinos from the sun was less than that predicted by the solar models. This was greeted at first with considerable skepticism, but over the years Davis’ results began to be taken more seriously and, with the results of the SAGE and GALLEX experiments, the reality of the neutrino deficit has been fully accepted. Davis’ work opened this field.
The second person I want to mention is Professor Masatoshi Koshiba, now retired from the University of Tokyo. In 1979, Koshiba proposed the construction of the huge water Cerenkov counter that became known as the Kamiokande Detector. Construction was completed in 1983 and initial experiments focused on the search for proton decay. By 1985 the simplest $`SU(5)`$ version of grand unified theories had been ruled out. With the addition of an outer veto layer to the detector beginning in 1984, Kamiokande became capable of detecting neutrinos as well. Kamiokande measured the solar neutrino flux, confirming Davis’ result, and also detected neutrinos from supernova SN1997A, opening up a new field of neutrino astrophysics. Even before much in the way of results from Kamiokande had come in, Koshiba proposed, in 1983, the construction of the very much larger Super-K. Under his leadership the project was approved just before his retirement from the University of Tokyo. The results from Super-K are what are generating all the excitement about atmospheric neutrinos.
The results of the solar neutrino experiments are summarized in Table 2. The SAGE and GALLEX experiments with their 235-kV neutrino energy threshold are the only ones that are sensitive to the proton-proton part of the solar cycle where almost all of the energy of the sun is produced. The Homestake experiment (Ray Davis’) with its 800-kV threshold is sensitive to the beryllium-7 line plus the boron-8 continuum. The Kamioka and Super-K experiments are sensitive only to the boron-8 continuum. The surprising, to me, result is the Homestake experiment which gives a ratio to the solar model of one-third, while the other four give a ratio of one-half. Any explanation of these results is complicated by the required energy dependence of the neutrino depletion process needed to account for all of the results. Super-K also reports a statistically marginal distortion of the neutrino spectrum compared to the standard solar model, and a 1-1/2 to 2 standard deviation day-night effect. Neither of these two effects is as yet of the significance required to be taken seriously.
Figure 1 is from Y. Suzuki’s talk at this conference summarizing the solar neutrino situation within the framework of neutrino oscillations. The upper part of the figure gives the allowed regions in the MSW model (a resonant conversion of the electron-neutrinos into another species of neutrinos) while the lower figure gives the allowed regions for pure vacuum oscillations. It is worth noting that if there is a systematic effect in the Homestake experiment, and all of the results were to be consistent with an energy-independent reduction of the solar neutrino flux, the large mixing-angle solution would be allowed with any $`\mathrm{\Delta }m^2`$ from about $`10^9`$ to $`10^4\mathrm{eV}^2`$.
The LSND experiment has been a problem since the first results were presented in 1995. It is the only experiment that purports to show the conversion of muon neutrinos into electron neutrinos with a large $`\mathrm{\Delta }m^2`$. The first results appeared to be in conflict with other experiments, particularly the KARMEN experiment at the Rutherford Laboratory. There is more data now, and as DiLella showed in his talk, there appears now to be a narrow range where the results of all of these experiments are consistent. Figure 2 shows the present situation.
The LSND experiment uses muon antineutrinos arising from the decay of mesons produced at the LAMPF 1-GeV proton accelerator. If these muon antineutrinos oscillate into electron antineutrinos, these can interact with protons in the detector tank (behind a thick shield) to produce positrons (20 to 60 MeV cuts) plus a neutron which is detected after a delay by a neutron-captured gamma ray. Electron neutrinos, which can be produced directly in the beam, can also interact in the detector, but they interact with carbon and there is no delayed neutron-captured gamma ray. The experiment analyzing all the data from 1993 through 1998, claims to see an excess of positrons in coincidence with a delayed neutron capture of 40 $`\pm `$ 9 events. That is certainly statistically significant, even with the statistical skepticism I evidenced earlier, but there are still many doubters and I include myself among the doubters.
Figure 3 is from the LSND paper and shows the distribution of accidental events in their detector tank. There is clearly a large excess of such events at the bottom front of the tank and this can only come from neutron leakage under their shield. This is a “beam on” background that the experimenters eliminate by focusing on a fiducial region toward the back and above the bottom of the detector.
Neutron diffusion is much less of a potential problem in the KARMAN experiment, because their machine produces a 10 $`\mu s`$ beam pulse while LAMPF produces a 600 $`\mu s`$ pulse. Neutron diffusion under the shield is a slow process and so KARMAN can gate most of this out with a relatively narrow time window around the beam pulse.
This experiment will have to be done again and it will be done again at FNAL in the MiniBooNE. The MiniBooNE detection limits are shown in Fig. 4 and the first results should be available in the year 2003. We have to wait.
The most exciting news on the atmospheric neutrino front comes from the Super-K data. Figure 5 shows an apparent disappearance of muon neutrinos as a function of zenith angle. Muon-like events are strongly depleted when they are generated by neutrinos passing through the entire earth compared to those generated by neutrinos coming down through the relatively thin covering of the Super-K detector. The allowed mass difference region assuming $`\nu _\mu `$ to $`\nu _\tau `$ oscillations is shown in Fig. 6. The Macro and Soudan experiments see a similar effect though with much looser mass constraints. Note that the Super-K/Macro/Soudan data alone do not tell us what happens to the muon neutrinos, they only tell us that they disappear. Muon neutrinos could go to sterile neutrinos that don’t interact; $`\tau `$ neutrinos that can only generate neutral current events since all the neutrinos are below $`\tau `$ production threshold; or even to some extent into electron neutrinos since it is hard to tell an enhancement of electron neutrinos from a shortage of muon neutrinos in this data.
The CHOOZ reactor experiment, in conjunction with the Super-K data, imposes the tightest restrictions on muon neutrino electron neutrino mixing. Fogli et al. analyzes Super-K and CHOOZ data together (Fig. 7). The shaded regions give the 90% and 99% allowed regions for the two experiments. The analysis done is done in the “dominant mass” mode and Super-K alone allows a quite large mixing of electron neutrinos. Super-K combined with CHOOZ, however, says that if there are oscillations, the muon neutrino oscillates almost entirely into tau neutrinos. Another reactor experiment, the Palo Verde experiment, has only been running for a relatively short time but already confirms the CHOOZ result in the mass range favored by Super-K.
Before discussing the next generation of experiments, it is useful to pause to summarize what we know:
1. Neutrinos from the sun are below the predictions of the standard solar model. All of the experiments integrate over neutrino energies above some threshold. The SAGE and GALLEX experiments, with a threshold of 235 keV, are a factor of two below the expectations; the Homestake experiment, with a threshold of 817 keV, is a factor of three below expectation; and the Kamiokande and Super-K experiments, with thresholds of 5 to 7 MeV, are a factor of two below expectations.
2. The atmospheric neutrino experiments clearly see an azimuthal dependence of the muon-neutrino/electron-neutrino ratio indicating a decrease of muon neutrinos coming 12,000 km through the earth.
3. The LSND experiment claims to see muon-neutrino to electron-neutrino conversion, but this experiment needs confirmation.
4. The CHOOZ experiment is the most sensitive of the reactor experiments and sees no loss of reactor-generated electron neutrinos over distances in the order of kilometers.
Taking all of this together, the favored explanation is neutrino-flavor oscillations plus the MSW effect. However, this is not the only explanation consistent with the data we have so far. For example, Barger et al. proposed to explain the data with a mixture of neutrino decay and the MSW effect. Gonzales-Garcia et al. hypothesized flavor changing in neutral current interactions of neutrinos. Both of these alternative hypotheses invoke new physics but so does the favored explanation. The job of the next generation of experiments is to sort all of this out.
On the solar neutrino front, we will certainly get more data from Super-K and data from two new experiments, SNO (already operating) and Borexino (due to start up in one to two years). I doubt that more statistics from Super-K will tell us anything new, but the other two experiments certainly will. Of particular interest to me is Borexino, which should be able to resolve the Be<sup>7</sup> line and so pin down any energy dependence with some precision.
The KamLAND reactor experiment in Japan is also important. Its electron neutrinos come from twelve nuclear reactors at an average distance of 150 km. KamLAND is a conversion of the old Kamiokande experiment and it should have a sensitivity to electron-neutrino oscillations down to $`\mathrm{\Delta }m^2`$ less than $`10^5\mathrm{eV}^2`$, covering the large mixing-angle MSW solution to the solar neutrino deficit.
Super-K will get much more data on atmospheric neutrinos. Of particular interest would be data around 90 degrees from the vertical. If oscillation is the answer and the mass difference is around Super-K’s central value, their sub-GeV and multi-GeV samples will show different behavior in this region. The 0.5-1 GeV neutrinos would be near their first oscillation minimum at 200-km distance, while the above 1-GeV neutrinos would show little reduction. Super-K may not have enough flux, consistent with the required angular resolution, to do this analysis.
The K2K experiment takes muon-neutrinos from KEK’s 12 GeV proton synchrotron to the Super Kamiokande detector and has just begun taking data. This is a disappearance experiment using muon neutrinos of less than 2 GeV. If the mass difference is above about 2 $`\times 10^3\mathrm{eV}^2`$ they should clearly see an effect in two to three years of data taking. This would be the first independent check of the Super-K result. In principle it is possible for the K2K experiment to definitively confirm the oscillation hypothesis. At the central value of the Super-K mass difference, 500 MeV neutrinos oscillate away and 250 MeV neutrinos oscillate back. However, the neutrino flux and the Super-K energy resolution are probably not good enough to see this, but it would be wonderful if they were.
I have already mentioned the MiniBooNE experiment. It supplies a definitive check on LSND and we will have to wait a few years for results.
The long baseline experiment, MINOS, is under construction at Fermilab and at the Soudan Mine. It is expected to begin data taking in the year 2003. MINOS is capable of seeing tau neutrino appearances over most of the mass difference range allowed by Super-K. Under MINOS’ conditions, neutrinos of 2 GeV will oscillate away and 1 GeV will return. If I were running the experiment, I would certainly tune the beam to low-energy neutrinos for the first few years.
CERN and the Gran Sasso Laboratory in Italy are discussing a European long-baseline experiment. The beam and the distance are very similar to the MINOS experiment. If this experiment is approved, I hope the apparatus is sufficiently different from MINOS to make the investment worthwhile.
Finally, there is much discussion of the potential of a muon storage ring as a high-intensity neutrino source. Such a facility would be much simpler to build than the muon collider that has been discussed for the past few years, but there are issues that couple the storage-ring design to the experiment. These issues have not yet been fully explored. Figure 8 shows what the excitement is all about; lots of neutrinos and greatly improved sensitivity. Figure 9 shows the problems; the storage ring is a mixed source of, for example, muon-neutrinos and electron-antineutrons.
The muon and electron neutrino spectra from unpolarized muons in the storage ring (much easier for the machine builder) are not that different, and only a detector with really good energy resolution could separate them on a statistical basis. To complicate the situation further, if the muon beam in the storage ring has a transverse momentum comparable to the muon mass (this gives the highest flux), then the two spectra are smeared further making the two types still harder to separate.
Polarization of muons in the storage ring allows the electron neutrinos to be tuned all the way to zero (in principle). However, the best polarization that the machine designers have come up with so far is in the range of 20-30% and even achieving this costs considerably in complexity and somewhat in flux in the machine.
Groups in the U.S. and in Europe are working on the machine and on the physics. They need to work more closely together to better understand the trade-offs between the experiment and the machine so that both can be properly optimized in order to carry out the required physics program.
## 4 Accelerators
At this conference, J. Lykken has spoken of the physics potential of the next generation of proton and electron accelerators, and G. A. Voss has spoken of the state of technologies. I can, therefore, be brief. The next big machine, the LHC, is under construction at CERN and is close to being the first world project. Contributions to the accelerator and detector have been made by many nations outside the group of CERN member states. The LHC will have an energy of 14 TeV in the proton-proton center of mass, a luminosity of 10$`{}_{}{}^{34}\mathrm{cm}_{}^{2}s^1`$ and a mass reach of about 1 TeV. Operations are expected to begin in the year 2005.
The two main experiments ATLAS and CMS will each have 1500 to 2000 collaborators. The size of these collaborations is unprecedented and presents difficult organizational problems in getting ready and new sociological problems in operation. In the 500-strong collaborations of today, we already have a bureaucratic overlay to the science with committees that decide on the trigger, data analysis procedures, error analysis, speakers, paper publications, etc. The participating scientists are imprisoned by golden bars of consensus. While we have survived so far and preserved some opportunity for scientific initiative, this will become more difficult as the collaborations grow to three times the size of today’s largest. This needs thinking about and talking about, but is a topic for another time.
The theorists tell us what they want us to find with the new accelerators, the source of electroweak symmetry breaking. But, that is not the experimenters’ job. Our job is to find out what is really there and, in this case, that means to find if there is any electroweak symmetry making. The experiments are difficult and the detectors are complex, expensive devices. Data rates, particularly at proton colliders, are enormous and there is no way to digest it all. Complex, multi-tiered trigger systems are needed to reduce the flood of data coming from the machine by a factor of 10 million or more so that our computer systems can handle the load. Those events that do not pass the trigger screen are discarded.
There is a danger here. Will we set up the experiments that can only find what we expect to find? Some years ago I asked the spokesperson for a large experiment if they could see a magnetic monopole that decays into a large number of 0.5 to 1 GeV gamma rays (a model that was plausible at the time). The answer was “no”; the trigger screened out such isotropic events that had no jet-like clustering and no transverse energy unbalance. While this particular problem is fixed, there is a lesson here. Be prepared and allow some of the trigger to be dedicated to what may currently be unfashionable.
With the LHC well underway, the next accelerator facility will be an electron-positron linear collider. Lykken, in his presentation, told of the physics potential and how the linear collider seems a necessary compliment to the LHC. Electron-positron machines have one advantage over proton colliders: all of the cross-sections are within a few orders of magnitude of each other (Fig. 10) and it is still possible to record all of the events that occur. Beams in linear colliders can be polarized and this is a great help in untangling the physics. Research and development is nearly done and it is hoped that construction of a new facility can start in 2004 or 2005 and the physics program can begin in 2010. The machine will be costly and will surely have to be built by worldwide collaboration.
A critical issue not yet settled is the initial energy for a new facility and what, if any, expansion capabilities should be built into the initial design (spending some extra funds now can save time and costs later for an energy increase). There seems to be a convergence on 500 GeV as an appropriate initial center-of-mass energy with the potential to expand to something in the 1 to 1.5 TeV region. There is a worldwide physics study currently underway with coordinators in Asia, Europe and North America. Whatever the initial energy, this project will be expensive and will only happen if the entire high-energy physics community gets behind it.
There has been much talk about the possibility of muon colliders. The attraction of the muon collider is also its biggest problem. Synchrotron radiation, which may limit the performance of multi-TeV electron-positron linear colliders, is absent in a muon collider. This absence requires the development of a new method of emittance damping to shrink the phase space in which the muons are produced so that high luminosity can be obtained in the collision. This will require a major R&D effort and one is being planned. It will not be simple.
The latest story on the state of the R&D in systems design is given in Ankenbrandt et al. . The luminosity of such a collider appears not to be interesting until a multi-TeV energy is reached. The indicated low-energy option has luminosity less than that of LEP-II, while the medium-energy example has luminosity less than the linear collider. All of the cross sections for muon colliders are the same as those for an electron-positron collider except for an s-channel process coupled to mass, Higgs-boson production. Here also the muon collider is not as effective as the electron collider. The rate for the production of 100-GeV mass Higgs in the muon collider is given in the reference above to be about 4000 per year. The 500-GeV electron-positron machine produces five times that rate through the Higgs plus $`Z^0`$ channel and the $`W`$-fusion channel.
The muon source itself is expensive. Current studies use a 4-MW proton source to produce the muons, four times the power of the SNS spallation neutron source which is estimated to cost $1.3 billion. Perhaps this can be done for less by the upgrade of one of the existing proton machines, but any real test of the muon collider concept will require a physics program that justifies the large cost. Perhaps the muon storage rings being discussed as neutrino sources can supply the justification for a real trial of the technology.
## 5 Non-Accelerator Experiments
Non-accelerator experiments are playing an ever more important role in high-energy physics; witness the time spent at this conference on Super Kamiokande neutrinos, cosmic microwave background fluctuations, and supernova distributions. These experiments allow tests of our theories not possible with accelerators. More such experiments are coming, such as:
* The Sloan Digital Sky Survey which will map the large-scale structure of luminous matter in the universe to much larger $`z`$ than now. This information is critical to test cold versus warm versus hot dark-matter scenarios.
* The AUGER ultra-high energy cosmic ray experiment which will get considerably more information on cosmic rays of mysterious origin with energies beyond the cutoff from interaction with the $`3^{}K`$ microwave background radiation.
* The Whipple, Hegra, Egret, GLAST experiments, etc., looking for an explanation of the origin of ultra-high-energy cosmic gamma rays.
* Neutrino observatories, under the ice and under the oceans.
* Follow-ons to COBE and the Supernova Cosmology Project that will give clues to the basic structure of the universe.
Some in high-energy physics are concerned that such experiments may drain funds from accelerator-based activities that have dominated high-energy physics for many decades. That may be so, but high-energy physicists should go where the high-energy physics is, whether it be in space, underground, underwater, or with accelerators. There is more to high-energy physics than the Higgs boson, and collaborations with the astrophysicists and cosmologists are more likely to broaden rather than narrow support for the things we are interested in through collaborations with new (to us) areas of science.
Two particularly important programs have been discussed at this conference, the cosmic-microwave-background radiation and the supernova search. The microwave- background radiation experiments measure the fluctuations in temperature (on the order of $`10^5K`$) as a function of the angular scale of the fluctuations. The pattern of these temperature fluctuations is a critical test of inflation, which ten years ago was thought to be an untestable metaphysical concept. It is testable now (see, for example, C. Contaldi, et al. ).
The experiment looks back to about 300,000 years after the big bang when the temperature in the early universe had dropped to the point where electrons could form hydrogen atoms (see Fig. 11). At that time the mean-free path of light changed from very short to very long as free electrons were bound to atoms; metaphorically the universe changed from foggy to clear. That light has been red- shifted by the Hubble expansion of the universe to become the $`3^{}K`$-background microwave-background radiation. Cold spots cooled early and hot spots later. The expansion leads to the fluctuations in temperature in the microwave background.
The available data are summarized in Fig. 12, which plots the power spectrum of the fluctuations versus spherical harmonic number. The information needed for test of cosmology occurs at high-harmonic numbers and the data is not precise enough yet to constrain the models. Over the next five years two new experiments will be launched, the MAP satellite by NASA, and the Planck satellite by the European Space Agency. These experiments are designed to be remarkably precise. The expected errors at large $`\mathrm{}`$ are expected to be comparable to the width of the lines in Fig. 12 up to $`\mathrm{}1000`$ for MAP and 2000 for Planck. Beyond these, the errors blow up very rapidly. It requires extraordinary precision to distinguish between cosmological models and, when both satellites are up, I’ll believe the data from two independent satellites with two independent instruments.
Perlmutter discussed the supernova cosmology project which addresses one of the most fundamental questions in physics, the constancy of the expansion rate of the universe. The results seem to say that it has not been constant, that the matter density of the universe is less than the critical density expected, and that there seems to be a cosmological constant in general relativity.
The experiment uses Type 1a supernovas as standard candles, measuring the apparent brightness (distance by $`R^2`$) and red shift $`z`$ (distance by Hubble expansion). In an unscientific sample of 15 high-energy physicists, I found only one who understood how Type 1a supernovas worked and what generates the light, and so I will briefly digress to tell you about it.
Type 1a supernovas are thought to come from white dwarf stars in binary systems that, over time, accrete matter from their companions until they reach a critical mass (1.4 solar masses). At that point, the white dwarf collapses and explodes as a supernova. Neutrinos escape immediately (SN1987A, for example). The debris cloud is heated by the decay of radioactive nickel (the minimum of the nuclear- binding energy curve), but the light is trapped because the cloud is optically opaque. As the cloud expands, its optical thickness decreases ($`R^2`$ for a uniform sphere, and $`R^1`$ for a shell), and the light intensity rises. It falls with the decay of the radioactive heat source. Type 1a supernovas are rare, occurring about one per 500 years per galaxy. The collaboration has collected fifty of them out to a $`z`$ of 0.8.
The data are shown in Fig. 13 that plots the apparent magnitude versus red shift. The data clearly deviate from the straight line expected for a flat universe with a matter density equal to the critical density. Perlmutter, et al. find that for a flat universe the matter density is 0.28 $`\pm `$ 0.085 (statistical) $`\pm `$ 0.05 (systematic) of that expected for a flat universe. The rest has to be made up by a cosmological constant.
Figure 14 shows the results in the matter-density/cosmological-constant plane. It is very far from being consistent with what we all had been assuming for many years: a matter density equal to one and a cosmological constant of zero. Personally, I would like to see this analysis done in a slightly different fashion. Since the universe obviously has some matter in it, the matter-density should be constrained to be greater than zero. This constraint will rotate the error ellipse and shift it somewhat.
The real question about the analysis has to do with the assumption that Type 1a supernova are really standard candles, i.e., do supernovae that exploded six or seven billion years ago ($`z=1`$) give the same light output as those that explode today? Stars formed recently have more heavy elements in them than did stars formed many billions of years ago, and I don’t know of any studies that look at what effect if any such a systematic difference in composition might have on light output. The collaboration has clearly gone through the detailed analysis of systematic errors, but I am not sure that I have high confidence in the systematic error quoted of 0.05 on the mass density. A systematic shift of one-quarter to one-half magnitude at $`z=1`$ would move the cosmological constant to zero. The data are superb, but the conclusions may not be. The collaboration will be collecting much more data and have proposed a dedicated satellite to extend the data out to a $`z`$ of two to three. If I had the money, I would give it to them—this is really important.
## 6 Theory
I want to talk about theory more philosophically than technically because I think it important that experimenters understand theory better so that they do not become mere technicians for those espousing the latest theoretical fad. Historically, advances in theory have synthesized data, accommodated previous theory as a special case, and simplified our view of the world. We have a bias toward elegance (in the eye of the beholder, of course) in the choice of theories and that has served us well in the past. The more mathematically and conceptually elegant theories are the ones that tended to survive.
All theories are born under-constrained in that there are always alternatives that compete in the contest with experiment, with the losers consigned to the dustbin. The sociologists of science would say that our theories are “socially constructed.” That is certainly true initially but our theoretical models are continually tested and the “social constructs” that don’t pass are discarded; though sometimes it takes a long time. Aristotle and Democritus had different views on whether matter was continuous or atomic in nature. Although there was no evidence to decide between these two views, Aristotle won, and for 2000 years it was believed that matter could be continuously subdivided and that there were no such things as atoms. We all believe that we are much more critical these days in examining assumptions.
The Standard Model when born in the 1970s was thought to be good with only one Higgs boson all the way up to the grand unification scale at 10<sup>15</sup> GeV. Since that time it has accumulated problems. It is hard to suppress strong CP violation. There is not enough CP violation to reproduce the baryon asymmetry of the universe. Longitudinal $`W`$ scattering is a problem if the Higgs mass is too high. There is no way to generate neutrino masses. It has 18 constants without a lepton sector CKM matrix and seven more with it. While the standard model has withstood all experimental tests, we know that it is only a low-energy (a few hundred GeV) approximation to a better model.
The most popular candidate to be the successor to the plain vanilla standard model is supersymmetry. It was introduced to stabilize the Higgs mass which is quadratically divergent in the standard model and only logarithmically divergent in the supersymmetric variants of the standard model. Supersymmetry does reduce to the standard model at “low” energies, but it also introduces 80 real and 44 complex constants. The theorists who are fans of supersymmetry are groping for variants that reduce these 124 new constants to a handful. If the supersymmetric successor to the standard model cannot reduce the total number of constants, it would seem to me to be a step backwards rather than an advance.
To the experimenters I would say that supersymmetry is a pure “social construct” with no supporting evidence despite many years of effort. It is okay to continue looking for supersymmetry as long as it doesn’t seriously interfere with real work (top, Higgs, neutrinos, etc.).
String/brane theory is in a very different situation. It represents an attempt to bring together gravity and quantum mechanics, a problem worth serious effort. It has too many dimensions for those of us living in a four-dimensional world, but these are early days and perhaps they’ll go away in the appropriate limit. There are several alternates that, on further examination, seem to turn into each other through duality transformations. Hidden in the center of all of these alternates are likely to be some kind of phase transitions that may lead to experimental signatures like those found for inflation. String/brane theory may even give the necessary constraints that supersymmetry needs to reduce the number of constants to a believable level. It is still too early to say, but it may be much more than metaphysics.
## 7 Concluding Observations
* Experimenters (and phenomenologists) need to be more concerned about systematic errors and the tails on error-distribution functions.
* Experimenters should learn more theory.
* All theorists should have a required course in statistics before receiving their Ph.D.
* We all hope for new things from LEP-II and the Tevatron, although the chances seem small. The new runs only increase LEP-II’s mass reach by about 8 GeV and the Tevatron’s by 50 GeV.
* We have big hopes for the $`B`$-Factories. The standard model’s CP violation is not enough and new directions may become clear from the factories. The first results should be in next summer.
* Neutrino physics is in ferment. More Super-K data, SNO, Borexino, KamLAND, Minos, K2K, should help to make things clear but it will take four to five years.
* Some redundancy in neutrino experiments is useful; too much is wasteful.
* I would love to see a low-energy muon-neutrino disappearance and reappearance experiment. Can it be done with Minos, or an appropriately designed CERN-Gran Sasso experiment?
* LHC starts up in 2005 and we all hope to find out what is beyond our standard model. The experiments are huge and the sociology will be complex. Beware of too many boards and committees.
* An $`e^+e^{}`$ collider of 0.5 to 1 TeV is a necessary companion to the LHC. It will only come to be if we all get behind it and push it as an international and regional program.
* Muon colliders, VLHC’s and exotic accelerator technologies are machines for 2020 or beyond. Muon colliders need R&D work to demonstrate that damping can work and then we’ll still face formidable problems. VLHC’s are a fantasy now. In addition to cost breakthroughs they need some serious accelerator physics studies (for example, already at LHC some power supplies now need tolerances of one part in a million). Exotic techniques such as plasmas and lasers are still in their infancy. They have achieved accelerating gradients of 1 GeV per meter, but only over a distance of a millimeter.
* Muon storage rings as neutrino sources are interesting and much simpler than are muon colliders. A study of how to use polarization, mixed electron and muon neutrino beams and their antiparticles is needed. The machine and the experiments interact strongly.
* Non-accelerator experiments in space and on the ground will be of increasing importance. The high-energy physics community should not be too parochial.
* The string theorists are doing great things. I hope they justify or eliminate supersymmetry and think up an experimental test.
* There is an exciting future: the work will be difficult, expensive and rewarding. The young generation, with support from governments, can and will do it. This is not “The End of Science.” |
warning/0001/astro-ph0001170.html | ar5iv | text | # Synthesis of 19F in Wolf-Rayet stars
## 1 Introduction
For long, the solar system has been the only location of the Galaxy with a known fluorine ($`{}_{}{}^{19}\mathrm{F}`$) abundance. At the same time, the production site(s) of this element has been a major nucleosynthetic puzzle, even if F is the least abundant (mass fraction of $`\mathrm{4\hspace{0.17em}10}^7`$, following Grevesse & Sauval 1998) of the elements ranging from carbon to calcium.
These last years, the situation has changed quite dramatically, both observationally and theoretically. Fluorine overabundances (with respect to solar) in MS, S and C stars have been reported (Jorissen et al. 1992), and correlate in particular with s-process enrichments. These observations demonstrate that thermally pulsating Asymptotic Giant Branch (AGB) stars are fluorine producers, as predicted by Goriely et al. (1989), and confirmed by calculations conducted in the framework of detailed AGB models (Forestini et al 1992, Mowlavi et al. 1996, 1998). It remains of course to determine the exact level of the contribution of these (mass losing) stars to the solar system and galactic F content.
In direct relation with this question, various calculations have been made in order to estimate the $`{}_{}{}^{19}\mathrm{F}`$ yields from massive stars. The neutrino process operating during supernova explosions has been envisioned as a possible producer of primary $`{}_{}{}^{19}\mathrm{F}`$ (e.g. Woosley & Weaver 1995). On the other hand, Meynet & Arnould (1993a) have investigated on grounds of detailed stellar models the suggestion (Goriely et al. 1989) that the hydrostatically burning He-shell can synthesize $`{}_{}{}^{19}\mathrm{F}`$ of secondary nature. They find that the level of production is relatively modest in $`M<\mathrm{\hspace{0.17em}\hspace{0.17em}20}\mathrm{M}_{}`$. In contrast, they show that stars which are massive enough to become Wolf-Rayet (WR) stars can eject through their winds substantial amounts of fluorine synthesized in the core at the beginning of the He-burning phase.
In the present work, we revisit the question of the galactic contribution of WR stars to $`{}_{}{}^{19}\mathrm{F}`$ with the help of new stellar models that better account for many important observable properties of WR stars. In addition, we extend the range of masses and metallicities considered in our previous study. The broadening of the explored metallicity range may take some additional importance in relation with the recent claim by Timmes et al. (1997) that “positive detection of any fluorine at a sufficiently large redshift ($`z>\mathrm{\hspace{0.17em}\hspace{0.17em}1.5}`$) would suggest strongly a positive detection of the neutrino process operating in massive stars”. The possibility of a significant thermonuclear production of $`{}_{}{}^{19}\mathrm{F}`$ by WR stars of different metallicities might blur this picture, and might at least imply the necessity of establishing observationally the primary or secondary nature of the detected fluorine, if any.
The physical ingredients of the models are discussed in Sect. 2. Section 3 presents our predicted yields from individual WR stars, while Sect. 4 gives a rough estimate of the contribution of WR stars to the galactic $`{}_{}{}^{19}\mathrm{F}`$ content. Some conclusions are drawn in Sect. 5.
## 2 The physical ingredients of the stellar models
The evolutionary models are computed with the same physical ingredients as in Meynet et al. (1994). However, the adopted nuclear reaction network is extended, especially in order to include the reactions involved in the production and destruction of $`{}_{}{}^{19}\mathrm{F}`$ (see below).
The differences with respect to the computations of Meynet & Arnould (1993a) are twofold:
1) A more extended range of initial masses (from 25 to 120 $`\mathrm{M}_{}`$) and metallicities ($`Z=0.008`$, 0.02 and 0.04) is explored;
2) The present grid of models is computed with the mass loss rates adopted by Meynet et al. (1994), which are twice as large as the values of $`\dot{M}`$ recommended by de Jager et al. (1988) and Conti (1988) for the pre-WR and WNL phases. This mass loss rate prescription enables to account for the observed variations of WR populations in different environments (Maeder & Meynet 1994).
The metallicity dependence of the mass loss rates during the pre-WR phases is adopted from previous works (e.g. Meynet et al. 1994). More specifically, $`\dot{M}`$ scales with metallicity $`Z`$ according to $`\dot{M}_Z/\dot{M}_{}=(Z/Z_{})^{0.5}`$, where $`Z_{}`$ is the solar metallicity. This scaling is deduced from stellar wind models (cf. Kudritzki et al. 1987, 1991).
Let us finally add that the models are computed with a moderate core overshooting ($`d/H_p=0.20`$, where $`d`$ is the overshooting distance and $`H_p`$ the pressure scale height at the boundary of the classical core).
### 2.1 The thermonuclear $`{}_{}{}^{19}\mathrm{F}`$ production and destruction paths
The CNO mode of H-burning is responsible for the production and destruction of $`{}_{}{}^{19}\mathrm{F}`$ through the reaction chain
$${}_{}{}^{14}\mathrm{N}(p,\gamma )^{15}\mathrm{O}(\beta ^+)^{15}\mathrm{N}(p,\gamma )^{16}\mathrm{O}(p,\gamma )^{17}\mathrm{F},$$
$${}_{}{}^{17}\mathrm{F}(\beta ^+)^{17}\mathrm{O}(p,\gamma )^{18}\mathrm{F}(\beta ^+)^{18}\mathrm{O}(p,\gamma )^{19}\mathrm{F}(p,\alpha )^{16}\mathrm{O}.$$
The adopted $`{}_{}{}^{19}\mathrm{F}(p,\alpha )^{16}\mathrm{O}`$ rate is the geometrical mean of the lower and upper limits to that rate proposed by Kious (1990).
Fluorine can also be produced and destroyed during He-burning through the chains (see also Meynet & Arnould 1993a)
$`(\beta ^+)^{18}O(p,\alpha )^{15}N(\alpha ,\gamma )^{19}F`$
$``$ $``$
$`{}_{}{}^{14}N(\alpha ,\gamma )^{18}F`$ $``$ $`(n,p)^{18}O(p,\alpha )^{15}N(\alpha ,\gamma )^{19}F(\alpha ,p)^{22}Ne.`$
$``$ $``$
$`(n,\alpha )^{15}N(\alpha ,\gamma )^{19}F`$
The synthesis of $`{}_{}{}^{19}\mathrm{F}`$ thus requires the availability of neutrons and protons. They are mainly produced by the reactions $`{}_{}{}^{13}\mathrm{C}(\alpha ,\mathrm{n}){}_{}{}^{16}\mathrm{O}`$ and $`{}_{}{}^{14}\mathrm{N}(\mathrm{n},\mathrm{p}){}_{}{}^{14}\mathrm{C}`$.
The first chain of transformation of $`{}_{}{}^{14}\mathrm{N}`$ into $`{}_{}{}^{19}\mathrm{F}`$ mentioned above is by far the most important in the conditions of relevance in this work, where the $`\beta ^+`$-decay lifetime $`\tau _\beta `$($`{}_{}{}^{18}\mathrm{F}`$) of $`{}_{}{}^{18}\mathrm{F}`$ is much shorter than its lifetime $`\tau _{n,\alpha }`$($`{}_{}{}^{18}\mathrm{F}`$) or $`\tau _{n,p}`$($`{}_{}{}^{18}\mathrm{F}`$) against ($`n,\alpha `$) or ($`n,p`$) reactions. For example, $`\tau _\beta `$ is a few hours only at the center of a 60 M model at the beginning of core He-burning, while the corresponding $`\tau _{n,\alpha }`$($`{}_{}{}^{18}\mathrm{F}`$) and $`\tau _{n,p}`$($`{}_{}{}^{18}\mathrm{F}`$) amount to about 1 400 and 18 000 years, respectively.
The NACRE compilation of reaction rates (Angulo et al. 1999) was not available yet at the time of completion of the calculations reported here. This is why most of the necessary nuclear data are taken from Caughlan & Fowler (1988). There are some exceptions to this rule. In particular, the $`{}_{}{}^{13}\mathrm{C}`$ $`\alpha `$-capture rate is taken from Descouvemont (1987), whose theoretical prediction of an increase of the astrophysical S-factor at low energies is confirmed experimentally (see NACRE). The $`{}_{}{}^{14}\mathrm{N}(\mathrm{n},\mathrm{p}){}_{}{}^{14}\mathrm{C}`$ rate is taken from Brehm et al. (1988). It is a factor of two lower than the one proposed by Koehler and O’Brien (1989), and leads consequently to a lower limit of the calculated $`{}_{}{}^{19}\mathrm{F}`$ yields.
## 3 Predicted $`{}_{}{}^{\mathrm{𝟏𝟗}}𝐅`$ yields from individual WR stars
As discussed by Arnould et al. (1999) on grounds of the NACRE rates, $`{}_{}{}^{19}\mathrm{F}`$ could be overproduced (with respect to solar) by the CNO cycle only at temperatures around $`15\times 10^6`$ K, the exact level of this overproduction remaining poorly predictable, however, in view of remaining rate uncertainties. This conclusion contradicts the one derived from the use of the rates recommended by Caughlan & Fowler (1988), in which case fluorine can never emerge in significant amounts from the CNO burning. As the latter rates are adopted in our calculations, the CNO zones of the computed model stars are depleted in $`{}_{}{}^{19}\mathrm{F}`$. This translates directly into a decrease of the $`{}_{}{}^{19}\mathrm{F}`$ mass fraction $`X_{19}^\mathrm{s}`$ at the stellar surfaces when the $`{}_{}{}^{19}\mathrm{F}`$-depleted CNO ashes are uncovered by mass loss (with the choice of the ordinate scales, the changes of fluorine abundance at the center and at the surface during the H-burning phase are not visible on Fig. 1). With the NACRE rates, it is expected that more $`{}_{}{}^{19}\mathrm{F}`$ would be present at the surface. However, it is also likely that this change is not able to affect drastically the predicted final yields, as these are dominated by the $`{}_{}{}^{19}\mathrm{F}`$ made during the He-burning phase.
In fact, as seen in Fig. 1, fluorine builds up through $`{}_{}{}^{14}\mathrm{N}(\alpha ,\gamma )^{18}\mathrm{F}(\beta ^+)^{18}\mathrm{O}(p,\alpha )^{15}\mathrm{N}(\alpha ,\gamma )^{19}\mathrm{F}`$ during the early phase of core He-burning. However, at the end of He-burning, $`{}_{}{}^{19}\mathrm{F}(\alpha ,\mathrm{p}){}_{}{}^{22}\mathrm{Ne}`$ is responsible for a significant $`{}_{}{}^{19}\mathrm{F}`$ destruction. Thus, material experiencing the whole He-burning episode cannot be $`{}_{}{}^{19}\mathrm{F}`$-enriched. In contrast, in massive stars going through the WR stage (initial mass $`M_\mathrm{i}>\mathrm{\hspace{0.17em}\hspace{0.17em}25}\mathrm{M}_{}`$ for $`Z=0.02`$, $`M_\mathrm{i}>\mathrm{\hspace{0.17em}\hspace{0.17em}35}\mathrm{M}_{}`$ for $`Z=0.008`$; see Maeder & Meynet 1994), some $`{}_{}{}^{19}\mathrm{F}`$ synthesized early during the core He-burning phase is ejected into the interstellar medium by stellar winds before its destruction. Indeed, Fig. 1 exhibits an increase of $`X_{19}^\mathrm{s}`$ when the He-burning products appear at the surface during the WC phase. As a result, the ratio $`X_{19}^\mathrm{s}(\mathrm{WC})/\mathrm{X}_{19}^{}`$ of the average $`{}_{}{}^{19}\mathrm{F}`$ surface mass fraction during the whole WC phase to the solar system $`{}_{}{}^{19}\mathrm{F}`$ mass fraction takes values as high as about 55, 95 and 60 in the case of the $`60\mathrm{M}_{}`$ model stars with $`Z=0.008`$, 0.02 and 0.04, respectively.
Figure 2 shows the $`{}_{}{}^{19}\mathrm{F}`$ “wind” yields for the computed stars ($`M_\mathrm{i}`$, $`Z`$) with initial mass $`M_\mathrm{i}`$ and metallicity $`Z`$. These yields, noted $`p_{19}^{\mathrm{wind}}(M_\mathrm{i},Z)`$, are equal to
$$p_{19}^{\mathrm{wind}}(M_\mathrm{i},Z)=$$
$$_0^{\tau (M_\mathrm{i},Z)}\dot{M}(M_\mathrm{i},Z,t)[X_{19}^\mathrm{s}(M_\mathrm{i},Z,t)X_{19}^0(Z)]𝑑t,$$
$`(1)`$
where $`\tau (M_\mathrm{i},Z)`$ is the total lifetime of the star ($`M_\mathrm{i}`$, $`Z`$), $`\dot{M}(M_\mathrm{i},Z,t)`$ its mass loss rate at age $`t`$, $`X_{19}^\mathrm{s}(M_\mathrm{i},Z,t)`$ its $`{}_{}{}^{19}\mathrm{F}`$ surface mass fraction at age $`t`$, and $`X_{19}^0(Z)`$ its initial $`{}_{}{}^{19}\mathrm{F}`$ mass fraction, assumed to relate to $`X_{19}^{}`$ by $`X_{19}^0(Z)=(Z/Z_{})X_{19}^{}`$. These yields may be negative in case most of the ejected material has been depleted in fluorine.
Figure 2 demonstrates that the highest yields are obtained for stars with $`Z=Z_{}`$ and $`40<M_\mathrm{i}<\mathrm{\hspace{0.17em}\hspace{0.17em}85}\mathrm{M}_{}`$. At lower metallicities, the winds are indeed weaker, and thus uncover the He-burning core only for the most massive stars and when the <sup>19</sup>F has already been burnt. On the other hand, at higher metallicities and for $`M_\mathrm{i}>\mathrm{\hspace{0.17em}\hspace{0.17em}85}\mathrm{M}_{}`$, the H-burning core mass decreases so rapidly during the main sequence as a consequence of very strong stellar winds that the He-burning core becomes too small for being uncovered by the stellar winds.
The above discussion shows that the most important physical ingredient influencing the WR $`{}_{}{}^{19}\mathrm{F}`$ yields is the metallicity-dependent mass loss rates, quantities like convective core masses being less crucial in this respect. As a numerical example, the value for $`X_{19}^\mathrm{s}(\mathrm{WC})/\mathrm{X}_{19}^{}`$ rises from about 18 in the 60 M low mass loss rate model of Meynet & Arnould (1993a) to about 95 in the same model star computed in this paper with an increased $`\dot{M}`$ value. This high sensitivity to $`\dot{M}`$ might cast doubts on the reliability of the predicted $`{}_{}{}^{19}\mathrm{F}`$ yields. In fact, some confidence in the results presented in this paper may be gained by noting that our present choice of the mass loss rates allows to account for the variation with metallicity of the number ratio of WR to O-type stars in regions of constant star formation rate (Maeder & Meynet 1994).
## 4 Estimate of the contribution of WR stars to the galactic fluorine
In order to evaluate the level of $`{}_{}{}^{19}\mathrm{F}`$ contamination by the winds of WR stars on a galactic scale, we use the $`p_{19}^{\mathrm{wind}}`$ yields \[Eq. (1)\] in a very simple model of galactic chemical evolution making use of the closed box and instantaneous recycling approximations. We also suppose that only WR stars are able to affect the galactic $`{}_{}{}^{19}\mathrm{F}`$ budget through their winds, all other possible production or destruction sites being neglected.
In such conditions, the $`{}_{}{}^{19}\mathrm{F}`$ mass fraction $`X_{19}(t)`$ in the galactic gas at time $`t`$ is equal to (e.g. Tinsley 1980)
$$X_{19}(t)=\stackrel{~}{y_{19}}\mathrm{ln}[1/\sigma (t)],$$
$`(2)`$
where $`\sigma (t)`$ is the mass fraction of the gas in the Galaxy at time $`t`$, and $`\stackrel{~}{y_{19}}`$ is a representative time-independent approximation of the net yield of a stellar generation defined by
$$y_{19}(t)=\frac{1}{1R}_{M_1}^{M_2}p_{19}^{\mathrm{wind}}(M_\mathrm{i},Z(t))\mathrm{\Phi }(M_\mathrm{i})𝑑M_\mathrm{i},$$
$`(3)`$
where $`R`$ is the “returned fraction”, $`M_1`$ and $`M_2`$ the lowest and highest mass of the stars going through the WR phase, and $`\mathrm{\Phi }(M_\mathrm{i})`$ the initial mass function (IMF). It has to be noted that Eq. (2) would break down if the true time- (or $`Z`$-)dependent $`y_{19}(t)`$ yields were used instead of $`\stackrel{~}{y_{19}}`$. In order to evaluate the latter quantity, we notice that $`y_{19}(Z(t))`$ values of about $`10^7`$, $`\mathrm{5.7\; 10}^7`$ and $`\mathrm{2.2\; 10}^7`$ are obtained for $`Z=0.008`$, 0.02 and 0.04 if use is made of the $`p_{19}^{\mathrm{wind}}`$ values reported in Sect. 3 and of the (properly normalized) IMF derived by Kroupa et al. (1993). On such grounds, we just adopt the rough estimate $`\stackrel{~}{y_{19}}\mathrm{3\hspace{0.17em}10}^7`$.
If this approximation is used in conjunction with the value $`\sigma 0.2`$ considered to characterize the solar neighbourhood at the time of the solar system formation 4.5 billion years ago (see Prantzos & Aubert 1995, and references therein), Eq. (2) leads to $`X_{19}\mathrm{5\hspace{0.17em}10}^7`$ in the local $`Z=Z_{}`$ interstellar medium (to be compared with the solar system abundance of 4 10<sup>-7</sup>). Thus, our simple estimate predicts that WR stars might account for most of the solar system $`{}_{}{}^{19}\mathrm{F}`$ content. Even larger $`{}_{}{}^{19}\mathrm{F}`$ quantities would be predicted with the use of the $`{}_{}{}^{14}\mathrm{N}(\mathrm{n},\mathrm{p}){}_{}{}^{14}\mathrm{C}`$ rate of Koehler & O’Brien (1989)! After having faced for long the problem of the underproduction of $`{}_{}{}^{19}\mathrm{F}`$, the theory of nucleosynthesis might now live with the danger of its predicted overabundance. If this is confirmed by further studies, constraints will obviously have to be put on one model or another.
## 5 Implications of $`{}_{}{}^{\mathrm{𝟏𝟗}}𝐅`$ detection at high redshift
Any $`{}_{}{}^{19}\mathrm{F}`$ present at high redshifts has to have been synthesized in massive stars only. Timmes et al. (1997) have argued further that its detection at redshifts $`z>\mathrm{\hspace{0.17em}\hspace{0.17em}1.5}`$ would in fact be a signature of the $`\nu `$-process in massive star explosions. The possibility of $`{}_{}{}^{19}\mathrm{F}`$ production by non-exploding WR stars might in fact weaken this statement, and blur the picture substantially.
Of course, one has to acknowledge that the contribution from WR stars at high redshifts may be reduced as a direct result of the lower metallicities that appear to characterize such regions. According to observations of Damped Lyman $`\alpha `$ systems (Pettini et al. 1997), the metallicity at redshifts between 1.5 and 2 indeed lies around $`0.1Z_{}`$. Such a reduced metallicity lowers the WR $`{}_{}{}^{19}\mathrm{F}`$ yields for two reasons. First, the number of WR stars predicted by non-rotating single star models is considerably reduced as a result of lower mass losses (Maeder & Meynet 1994). Second, the abundances of the CNO seeds that are needed for the secondary WR $`{}_{}{}^{19}\mathrm{F}`$ production are reduced as well.
Even so, it would certainly be premature at this point to completely forget about the role of WR stars in a possible enrichment of high-$`z`$ material with $`{}_{}{}^{19}\mathrm{F}`$, and to relate it strictly with the $`\nu `$-process.This is even more true as the predictions reported in this paper are based on single, non-rotating stellar models only. How binarity and/or rotation would change these results remains to be checked. At present, the published rotating evolutionary models leading to WR stars (Fliegner & Langer 1994, Meynet 1998, 1999) make no predictions concerning the synthesis of fluorine. However they show that rotation favours an early entrance into the WR phase for a given mass, and decreases the minimum initial mass for a star to go through a WR phase at a given metallicity. Moreover, the mixing induced by rotation opens up new nucleosynthetic channels (see Heger 1998) whose importance for the scenario of fluorine production presented in this paper remains to be quantitatively assessed. Finally, let us note that the effects of rotation might be more important at low $`Z`$ if, as suggested by Maeder et al. (1999), the average rotation is faster at low metallicities. In such conditions, and in absence of quantitative calculations, one has to remain alert to the possibility of a significant contamination of low metallicity high redshift regions by the $`{}_{}{}^{19}\mathrm{F}`$-loaded wind of WR stars.
Clearly, observations of $`{}_{}{}^{19}\mathrm{F}`$ at high redshift, if possible at all, would be decisive in order to answer the question of the very production mechanism of this element. An important distinguishing feature would be the primary nature of the observed $`{}_{}{}^{19}\mathrm{F}`$, as predicted by the $`\nu `$-process, or its secondary behaviour, as expected from the thermonuclear model discussed in this paper.
## 6 Conclusion
Detailed stellar model predictions made in the framework of a very rough model for the chemical evolution of the solar neighbourhood leads to the conclusion that non-exploding non-rotating single WR stars alone could account for the solar $`{}_{}{}^{19}\mathrm{F}`$ content. This conclusion remains to be ascertained by the adoption of a more realistic galactic evolution model. Still, it appears likely that the considered WR stars might be significant, and even possibly dominant, galactic $`{}_{}{}^{19}\mathrm{F}`$ contributors. In addition, they might well be responsible for a $`{}_{}{}^{19}\mathrm{F}`$ enrichment, if any, of high-redshift ($`z>\mathrm{\hspace{0.17em}\hspace{0.17em}1.5}`$) low-metallicity regions ($``$ 0.1 $`Z_{}`$). Further predictions are eagerly awaited for rotating, as well as binary, WR stars.
Finally, let us stress that the most direct test of the present model would be the measurement of the abundance of fluorine in the wind of WC stars. It remains to be seen if such observations are really feasible.
###### Acknowledgements.
We thank N. Prantzos for comments on the galactic chemistry aspects of this work. This research has been supported in part by the HCM Programme of the European Union under contract ERBCHRXCT 930339. |
warning/0001/cond-mat0001352.html | ar5iv | text | # Quasiparticle Inelastic Lifetime from Paramagnons in Disordered Superconductors
## I Introduction
In the early seventies, certain rare earth ternary compounds were found to display superconductivity while at the same time showed strong tendencies to be magnetic. A large body of theoretical work has been devoted to the interplay of magnetism and superconductivity. Recently, there is increasing evidence that there is an interplay of magnetism and superconductivity in the boro-carbides as well as RuSr<sub>2</sub>GdCu<sub>2</sub>O<sub>8</sub>. Presently it is widely debated whether these materials are conventional superconductors with sharply peaked density of states (DOS) near the Fermi level (similar to some A-15’s), or unconventional with ground state pairing of lower symmetry than the underlying lattice. For instance, in $`R`$=Lu, Y $`R`$Ni<sub>2</sub>B<sub>2</sub>C, scanning tunneling microscopy (STM) has given evidence for conventional BCS behavior, albeit with a substantially smeared DOS, with a smearing parameter $`\mathrm{\Gamma }/\mathrm{\Delta }=0.2`$. Optical conductivity studies also support moderately strong coupled conventional superconductivity with $`2\mathrm{\Delta }/T_c=3.95.2`$ . At the same time de Haas-van Alphen, magnetic field anisotropy, and electronic Raman scattering experiments have given evidence for at least very small gaps over a portion of the Fermi surface.
Substitutional or positional disorder has played a crucial role in determining whether a material is conventional (i.e., obeys Anderson’s theorem) or not, and recently studies on borocarbides doped with Co have been performed. Heat capacity and magnetic measurements on Y(Ni<sub>2-x</sub>Co<sub>x</sub>)B<sub>2</sub>C have interpreted the drop in T<sub>c</sub> with increasing Co doping as due to the reduction of the DOS at the Fermi level rather than pairbreaking by nonmagnetic impurities. On the other hand, Raman measurements on the same systems have shown an increase in spectral weight below the gap edge as Co is doped in, contrary to conventional BCS behavior.
However it is well known that conventional superconductors which are highly disordered display substantially smeared BCS properties which can mimic unconventional pairing. This can result from vanishing of phase coherence or from the interplay of interactions and disorder. The latter is most manifest in the reduction of quasiparticle (qp) lifetimes. Inelastic collisions broaden qp eigenstates and lead to a smearing of activated or threshold behavior in single- and two-particle correlation functions, measured e.g., by tunneling, optical conductivity, and electronic Raman scattering. While the present status of the superconducting ground state of the borocarbides remains unclear, it is of interest to inspect whether strong inelastic scattering can modify s-wave properties to the point where the ubiquitous exponential behavior of various thermodynamic and transport quantities is obscured. For instance, the absence of a coherence peak in NMR is usually interpreted as a signal of unconventional electronic pairing. However, it is well known that the coherence peak can be suppressed as a consequence of strong inelastic electron-phonon collisions . While the coherence peak can be fully suppressed only for large electron-phonon couplings in clean superconductors , it has been shown that the peak can be further suppressed in disordered superconductors and is absent in the region of strong disorder for only moderate couplings .
The microscopic interplay of disorder, magnetic fluctuations and superconductivity is reflected in the behavior of the qp inelastic lifetime. In this paper we present a calculation for the qp inelastic scattering rate due to spin fluctuations within a formalism developed in previous works . The calculation is undertaken by first obtaining an effective fluctuation propagator in the superconducting state, and then using the exact eigenstate formalism as used in the case of Coulomb scattering with the replacement of the Coulomb propagator and vertex with the derived fluctuation propagators and vertices. It is shown that the rate is qualitatively similar to the rate due to Coulomb interactions with addition of the Stoner enhancement. Finally we discuss our results in terms of STM data on the borocarbide superconductors.
## II Calculations
The scattering rate from paramagnons in clean superconductors on a lattice is well known for the case of $`s`$ or $`d`$wave superconductors. The calculation for the inelastic scattering rate due to paramagnon exchange in disorder metals is also well known. In both cases the results are similar to the scattering rate from long-range Coulomb interactions, with an additional enhancement via the Stoner factor $`1/(1I)`$, where $`I=UN_F`$, $`U`$ is a phenomenological short range interaction, and $`N_F`$ is the DOS per spin at the Fermi level. For dirty metals and superconductors, the electron-phonon interaction is reduced via “collision drag” relative to the clean case, while the electron-electron interaction is enhanced by disorder due to the breakdown of screening by diffusive electrons. The latter enhancement of the scattering rate at the Fermi surface is $`\widehat{\rho }^{3/2}(E_F/T)^{1/2}`$ compared to that of 3D clean materials. Here $`\widehat{\rho }`$ is the dimensionless measure of disorder, with $`\widehat{\rho }=\rho /\rho _M`$, where $`\rho `$ is the extrapolated residual resistivity and $`\rho _M`$ the Mott number, which in a jellium model is given by $`\rho _M=3\pi ^2/e^2k_F`$. We use units such that $`k_B=\mathrm{}=1`$. However, calculations for the scattering rate calculated for superconductors on a lattice have treated impurities and interactions independently and therefore do not capture the disorder enhancement derived for conventional superconductors. Therefore in this paper we investigate the interplay of disorder, superconductivity, and magnetism by revisiting the problem of inelastic scattering.
The spin fluctuation propagator is given by the sum of longitudinal $`K_{}`$ and transverse $`K_{}`$ paramagnons, respectively. They can be expressed in terms of the polarization bubble $`\chi `$ as
$`K_{}=`$ $`U\chi U+U\chi U\chi K_{},`$ (1)
$`K_{}=`$ $`U+U\chi K_{}.`$ (2)
Solving these equations we obtain the fluctuation propagator $`t(𝐪,\omega )`$
$$t(𝐪,\omega )=K_{}+K_{}=\frac{U^2\chi (𝐪,\mathrm{\Omega })}{1U^2\chi ^2(𝐪,\mathrm{\Omega })}+\frac{U}{1U\chi (𝐪,\mathrm{\Omega })}.$$
(3)
However in the superconducting state one must distinguish between charge and spin response couplings due to their different coherence factors. Therefore in the superconducting state the propagator splits into two contributions given by
$`t_c(𝐪,\mathrm{\Omega })={\displaystyle \frac{1}{2}}{\displaystyle \frac{U^2\chi _c(𝐪,\mathrm{\Omega })}{1+U\chi _c(𝐪,\mathrm{\Omega })}},`$ (4)
$`t_s(𝐪,\mathrm{\Omega })={\displaystyle \frac{3}{2}}{\displaystyle \frac{U^2\chi _s(𝐪,\mathrm{\Omega })}{1U\chi _s(𝐪,\mathrm{\Omega })}}U^2\chi _s(𝐪,\mathrm{\Omega }),`$ (5)
with $`\chi _{c,s}`$ the charge, spin susceptibilities, respectively. In the following we perform calculations in the continuum limit and neglect lattice effects. This is certainly important in order to capture strong scattering via qp exchange of antiferromagnetic reciprocal lattice vector momenta $`𝐐`$. However, the incipient magnetic instability via paramagnon exchange nevertheless is reflected via the Stoner criterion. Albeit a naive approach to the borocarbides or other materials with strong antiferromagnetic fluctuations, the results allow us to qualitatively estimate the effects of disorder on qp inelastic scattering from paramagnons.
The gauge-invariant charge polarization $`\chi _c`$ has been calculated in disordered superconductors in Ref. . It has the structure $`\chi (𝐪,\omega )=B(𝐪,\omega )+B_C(𝐪,\omega )`$. Here $`B`$ is the density response function in the pair approximation , while $`B_C`$ contains the collective excitation (the Anderson-Bogolubov mode) which restores gauge invariance. It was shown that for $`k_F\xi 1`$ collective effects can be ignored and that the ”pair approximation” for the polarization is adequate, where $`\xi =\sqrt{1/m\pi \mathrm{\Delta }}\widehat{\rho }^1`$ is the dirty-limit coherence length, For $`T=0`$ the polarization can be written as
$`\chi _c^{\prime \prime }(𝐪,\omega )=\varphi ^{\prime \prime }(𝐪,\sqrt{\omega (\omega 2\mathrm{\Delta })})\mathrm{\Theta }(\omega 2\mathrm{\Delta })`$ (6)
$`\times \left[(\omega +2\mathrm{\Delta })E(\alpha ){\displaystyle \frac{4\mathrm{\Delta }\omega }{\omega +2\mathrm{\Delta }}}K(\alpha )\right],`$ (7)
while the spin susceptibility is given by the Mattis-Bardeen result
$`\chi _s^{\prime \prime }(𝐪,\omega )=\varphi ^{\prime \prime }(𝐪,\sqrt{\omega (\omega 2\mathrm{\Delta })})\mathrm{\Theta }(\omega 2\mathrm{\Delta })`$ (8)
$`\times \left[(\omega +2\mathrm{\Delta })E(\alpha )4\mathrm{\Delta }K(\alpha )\right].`$ (9)
Here, $`\alpha =\frac{\omega 2\mathrm{\Delta }}{\omega +2\mathrm{\Delta }}`$, and $`E`$ and $`K`$ are complete elliptical integrals of the first and second kinds, respectively. $`\varphi ^{\prime \prime }`$ is the spectrum of the density Kubo function for noninteracting electrons. It can calculated by a variety of techniques for various limits of disorder. For clean metals, the spectrum is white,
$$\varphi ^{\prime \prime }(𝐪,ϵ)=\frac{m^2}{4\pi q},\mathrm{clean},$$
(10)
while for diffusive qp dynamics, $`\varphi ^{\prime \prime }`$ is given by a diffusion pole
$$\varphi ^{\prime \prime }(𝐪,ϵ)=N_F\frac{Dq^2}{(Dq^2)^2+ϵ^2},\mathrm{diffusive},$$
(11)
with $`D`$ the diffusion constant. Here we have neglected Cooper propagator renormalization, which can be shown to give a smaller contribution to the scattering rate than Diffusion propagator renormalization by a factor of $`1/k_F\xi `$.
The limiting behavior for finite temperatures with $`T<<\mathrm{\Delta }`$ is given as:
$`\chi _{c,s}^{\prime \prime }(𝐪,\mathrm{\Omega }<<\mathrm{\Delta })`$ $``$ $`\mathrm{\Omega }\varphi ^{\prime \prime }(𝐪,\sqrt{2\mathrm{\Delta }\mathrm{\Omega }})e^{\mathrm{\Delta }/T}`$ (13)
$`\times \{\begin{array}{cc}1,\mathrm{charge},\hfill & \\ (\mathrm{\Delta }/T)\mathrm{ln}(4T/\mathrm{\Omega }),\mathrm{spin},\hfill & \end{array}`$
$`\chi _{c,s}^{\prime \prime }(𝐪,\mathrm{\Omega }2\mathrm{\Delta })`$ $``$ $`\mathrm{\Delta }\varphi ^{\prime \prime }(𝐪,2\sqrt{2}\mathrm{\Delta })\sqrt{\pi T/\mathrm{\Delta }}e^{\mathrm{\Delta }/T}`$ (15)
$`\times \{\begin{array}{cc}1/2,\mathrm{charge},\hfill & \\ 1,\mathrm{spin}.\hfill & \end{array}`$
Thus the behavior of the spin and charge susceptibilities yields different contributions to the paramagnon scattering in the charge and spin channel.
The paramagnon contribution to the self energy can be split in the usual way into an anomalous and even and odd normal pieces. It has been shown for the case of long-range Coulomb interactions that the even part of the normal self energy contribution can be ignored, and can be shown for the spin-fluctuation case as well. Expanding near the qp pole in the BCS Green’s function, we obtain the expression for the on-shell inelastic scattering rate due to paramagnons exchange in the charge channel $`\mathrm{\Gamma }_c`$ and spin channel $`\mathrm{\Gamma }_s`$,
$`\mathrm{\Gamma }_{c,s}(\omega )={\displaystyle \frac{1}{Z^{}}}{\displaystyle \underset{𝐪}{}}{\displaystyle \frac{dϵ}{\pi N_F}\varphi ^{\prime \prime }(𝐪,ϵ\omega )}`$ (16)
$`\times {\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{dx}{\pi }}[f(x+\omega )+n(x)]t_c^{\prime \prime }(𝐪,x)`$ (17)
$`\times \left[G^{\prime \prime }(ϵ,\omega +x)\pm \mathrm{\Delta }/\omega F^{\prime \prime }(ϵ,\omega +x)\right]+(\omega \omega ),`$ (18)
where $`n,f`$ are Bose and Fermi distributions, respectively, $`G^{\prime \prime }`$ and $`F^{\prime \prime }`$ are the imaginary parts of the bare normal and anomalous BCS Green’s functions, respectively, $`Z^{}`$ is the real part of the qp renormalization, and $`(\omega \omega )`$ denotes the addition of terms which differ from the ones written only by the sign of $`\omega `$.
Substituting Eqs.(3-9) into Eq.(10), we obtain the inelastic scattering rate $`\tau _s^1=2\mathrm{\Gamma }_s`$. It can be shown that the contribution to the scattering rate from the charge channel yields a subdominant contribution for all values of disorder and interaction $`I<1`$ compared to the long-range Coulomb contribution calculated in Ref. . Therefore for the remainder of the paper we neglect $`\mathrm{\Gamma }_c`$ and focus on $`\mathrm{\Gamma }_s`$. The scattering rate is dominated by qp population at the gap edge. For $`T=0`$, an injected qp must have enough energy to give up to break a Cooper pair ($`3\mathrm{\Delta }`$) and for $`\frac{\mathrm{\Omega }3\mathrm{\Delta }}{\mathrm{\Omega }+3\mathrm{\Delta }}<<1`$ we obtain,
$`\mathrm{\Gamma }_s^{T=0}(\mathrm{\Omega }3\mathrm{\Delta })`$ $`={\displaystyle \frac{3I^2\pi ^2}{16(1I)^{3/2}}}{\displaystyle \frac{\mathrm{\Delta }}{Z^{}}}{\displaystyle \frac{\mathrm{\Delta }}{\mathrm{\Omega }}}{\displaystyle \frac{\mathrm{\Delta }}{E_F}}F\left({\displaystyle \frac{\mathrm{\Omega }}{\mathrm{\Delta }}}\right)`$ (20)
$`\times \{\begin{array}{cc}\frac{1}{\sqrt{1I}},\mathrm{clean},\hfill & \\ \left(\frac{3\widehat{\rho }}{\pi }\right)^{3/2}\sqrt{\frac{E_F}{\mathrm{\Delta }}},\mathrm{dirty},\hfill & \end{array}`$
with $`F(x)=x(x^2/2x1)\sqrt{(x2)^21}+(x/22)\mathrm{ln}[x2+\sqrt{(x2)^21}]`$. At finite temperatures, the Cooper pair recombination rate is dominated by the kinematic factor $`\mathrm{\Gamma }^Re^{2\mathrm{\Delta }/T}`$ and qp scattering rate $`\mathrm{\Gamma }^Se^{\mathrm{\Delta }/T}`$. For a qp at the gap edge, the dominant contribution to the recombination rate is given by
$`\mathrm{\Gamma }_s^R(\mathrm{\Delta }>>T)=`$ $`{\displaystyle \frac{3\pi ^2I^2}{8\sqrt{2}(1I)^{3/2}}}{\displaystyle \frac{T^2}{Z^{}E_F}}e^{2\mathrm{\Delta }/T}`$ (22)
$`\times \{\begin{array}{cc}\frac{1}{\sqrt{1I}},\mathrm{clean},\hfill & \\ \left(\frac{3\widehat{\rho }}{\pi }\right)^{3/2}\sqrt{\frac{E_F}{\mathrm{\Delta }}},\mathrm{dirty},\hfill & \end{array}`$
while for the scattering rate we obtain to leading order
$`\mathrm{\Gamma }_s^S(\mathrm{\Delta }>>T)=`$ $`{\displaystyle \frac{3\pi I^2\mathrm{ln}(2)}{8(1I)^{3/2}}}{\displaystyle \frac{\mathrm{\Delta }}{Z^{}}}\sqrt{{\displaystyle \frac{\pi T}{2\mathrm{\Delta }}}}e^{\mathrm{\Delta }/T}`$ (24)
$`\times \{\begin{array}{cc}\frac{1}{\sqrt{1I}},\mathrm{clean},\hfill & \\ \left(\frac{3\widehat{\rho }}{\pi }\right)^{3/2}\sqrt{\frac{E_F}{\mathrm{\Delta }}},\mathrm{dirty}.\hfill & \end{array}`$
We see a similar behavior between the paramagnon and long-range Coulomb contributions to the inelastic scattering rate. $`\mathrm{\Gamma }_s`$ possesses the same temperature dependence as the Coulomb contribution, with the exponential temperature dependence reflecting the necessity of two qps per scattering event. Further, we see the same disorder enhancement ($`\widehat{\rho }^{3/2}(E_F/\mathrm{\Delta })^{1/2}`$ ) relative to the clean case as in the long-range Coulomb case. Lastly, we note that the energy gap $`\mathrm{\Delta }`$ acts as a cut off for the divergence of the rate that occurs in the $`2d`$ dirty normal calculation , just as in the long-range Coulomb case.
On top of the disorder enhancement, there is the Stoner enhancement relative to the Coulomb contributions due to the nearness of a magnetic instability. In materials close to the instability, this contribution will be dominant over the Coulomb and phonon terms except for very low temperatures, where the power-law temperature dependence of the phonon contribution takes over . We note that our expression are valid for $`\mathrm{\Delta }/E_F<<1I<<1`$, i.e., provided that one is not too close to the Stoner criterion for magnetism, $`I=1`$. At the instability, the rate saturates as it does in the case of a normal metal near the metal-insulator transition . However in order to accurately describe the dynamics at the magnetic transition one needs to use a more sophisticated spin fluctuation propagator than the one derived here from RPA diagrams only, which tend to overestimate paramagnon effects .
Finally, we can compare the results to the values of the scattering rates inferred from STM data on clean and thin films of LuNi<sub>2</sub>B<sub>2</sub>C. To our knowledge, a temperature dependence of the scattering rate has not yet been published, nor has a reliable estimate of the scattering rate been made from optical (Raman or infrared) or Hall probes as has been done in the high T<sub>c</sub> cuprates. Moreover no systematic study of the effects of impurities and doping have been made concerning the scattering rate. Nevertheless we can estimate if inelastic scattering from paramagnons in an s-wave superconductor is sufficient to explain the broadening observed in STM measurements. As a rough estimate for $`1/\tau _s`$ we take $`I2/3`$, $`Z^{}=1/2`$, Fermi velocity $`v_F3.5\times 10^7`$ cm/s, Fermi energy $`E_F0.3`$ eV given from LDA estimates for LuNi<sub>2</sub>B<sub>2</sub>C from Ref. . STM data taken at low temperatures in Ref. gives $`\mathrm{\Delta }=18`$cm<sup>-1</sup>, which is consistent with Raman measurements. This yields a scattering rate for clean systems at $`T=0.5T_c`$ from Eq. (13) of $`1/\tau _s=1.3\times 10^3`$meV, or $`1/\tau _s\mathrm{\Delta }=6\times 10^4`$, which is clearly too small to match experiments. Either the scattering is most likely due to electron-phonon collisions or perhaps due to large gap anisotropy.
Since $`\rho (T=0)`$ increases quickly as Co is doped in, rising by over an order of magnitude for 15% Co doping, it may be feasible that the disorder enhancement for $`1/\tau _s`$ in an s-wave scenario could lead to spectral weight at low frequencies observed via magnetic field anisotropy or Raman measurements. An estimate for the Mott number is difficult since the parameters $`v_F,k_F`$ and the other parameters entering in Eqs. (12-13) are presumably disorder dependent, and it is not clear where the metal-insulator transition occurs for this compound. A conservative estimate from the Ioffe-Regel criterion in Ref. gives $`\rho _M400\mu \mathrm{\Omega }`$-cm. Therefore taking $`\rho (T=0)100\mu \mathrm{\Omega }`$cm as in Ref. into Eq. (13) only gives $`1/\tau _s\mathrm{\Delta }10^3`$, which is clearly too small to account for the large broadening observed via STM even in relatively clean films nor is it sufficient to account for the substantial spectral weight observed at low frequencies via Raman scattering. It is tempting to therefore conclude that the large broadening comes either from nodal qps in conventional (extended s-) or unconventional (d-) pair states.
However, there are problems in each scenario. Small amount of Co doping (on the few percent level) quickly push these materials into the dirty limit ($`\xi /l<<1`$). If the gap possessed extended $`s`$wave symmetry, the disorder would be sufficient to wash out any remaining anisotropy and necessarily lead to sharp threshold behavior, which is not observed. On the other hand, the disorder would also lead to a sharp drop in $`T_c`$ if the gap possessed $`d`$wave symmetry and unconventional superconductivity would be expected to be completely suppressed for 15% Co doping, which again is not observed. Therefore it is unclear from current data whether superconductivity is conventional or not, and perhaps the situation is clouded by the presence of additional non-superconducting bands, which would also yield a non-vanishing zero bias conductance and low frequency spectral weight. It would thus be extremely useful to study impurity and cation dopings further to determine if the enhanced scattering rates are responsible for the behavior indicative of unconventional pairing as the disorder is increased. Raman scattering measurements would be very useful in this regard, and remains a topic for further investigation. |
warning/0001/astro-ph0001515.html | ar5iv | text | # Toy model of obscurational variability in active galactic nuclei
## 1 Introduction
Radio-quiet active galactic nuclei (AGN) are well known to be variable in the X-ray band since early observations by EXOSAT (e.g. Lawrence et al. 1987, McHardy & Czerny 1987, Green, McHardy & Lehto 1993). The combined study of both the X-ray spectra and variability offer the most direct insight into the structure of accretion flow onto the black hole which powers the AGN (for a review, see Mushotzky, Done & Pounds 1993).
The variability trends have been extensively studied in various spectral bands (e.g. Ulrich, Maraschi & Urry 1997, Peterson et al. 1998). Generally, more luminous objects are less variable both in X-rays (Nandra et al. 1997 and the references therein) and in the optical/UV band (Ptak et al. 1998, Giveon et al. 1999).
Multi-wavelength studies showed that this variability is surprisingly complex. In the X-ray band, most of AGN vary but individual sources display various spectral trends with the brightening of the source (e.g. Ciliegi & Maccacaro 1997, George et al. 1998, Nandra et al. 1997). The correlated variability between different energy bands is also difficult to interpret (e.g. Nandra et al. 1998 for UV and X-ray connection in NGC 7469).
The question appears to be whether this observed variability is entirely intrinsic, i.e. related to strongly non-stationary release of gravitational energy of the in-flowing matter, or is caused, at least partially, by the effect of variable obscuration towards the nucleus.
Partial covering models were popular mostly at the beginning of spectral studies in the X-ray band (e.g. Mushotzky et al. 1978, Matsuoka et al. 1986). Recently, Seyfert 1 galaxies and QSOs are modeled through an accretion disk, with X-ray emission coming either from the disk corona or from the disrupted innermost part of the disk (e.g. Loska & Czerny 1997 and the references therein).
However, occasionally, the problem was revived. An eclipse by a cloud was successfully considered as a model for the faint state of MCG-6-30-15 (McKernan & Yaqoob 1998, Weaver & Yaqoob 1998) and it may be a possible cause of the $`K_\alpha `$ line profile variability in NGC 3516 (Nandra et al. 1999). In the case of the variability of Narrow Line Seyfert 1 galaxies (e.g. Boller et al. 1997 for IRAS 13224-3809; Brandt et al. 1999 for PHL 1092) it was suggested that the observed huge amplitudes are hard to explain directly through a variable energy output, and instead, partial covering and/or relativistic erratic beaming may be necessary. It is interesting to note that the partial covering mechanism may also apply to similarly variable galactic sources (Brandt et al. 1996 for Cir X-1). Obscuration events are also sometimes invoked to explain the variations of the Broad Emission Lines in Seyfert galaxies - for example, a spectacular change of morphological type from Seyfert 2 to Seyfert 1 of the galaxy NGC 7582 (Aretxaga et al. 1999) was interpreted as due to an obscuration event by a distant cloud (Xue et al. 1998). Finally, there is a large amount of partially ionized material lying along the line of sight to the nucleus which manifests itself as a warm absorber (cf. Reynolds 1997). Variability in absorption edges can either be due to a variation of the ionisation state of the absorber or to obscuration by clouds passing through the line of sight.
Although there are now strong suggestions that the accretion flow proceeds predominantly through a disk the physics of the innermost part of the flow is highly uncertain. As the inner regions of a standard Shakura-Sunayev (1973) disk are thermally and viscously unstable in an AGN owing to the large radiation pressure, the disk is frequently supposed to be disrupted. It may form a hot, optically thin, quasi spherical ADAF (advection-dominated accretion flow; Narayan & Yi 1994, and many subsequent papers), or it may proceed in the form of thick clouds (Collin-Souffrin & al 1996, hereafter referred as Paper I). Such clouds might be optically thick, unlike clouds optically thin for electron scattering which form spontaneously as a result of thermal instability in X-ray irradiated cold gas and coexist in equilibrium with the surrounding hot optically thin gas (Krolik 1998, Torricelli-Ciamponi & Courvoisier 1998). Such clouds provide significant mass flux and they are not easily destroyed (see Sect. 4 in Paper I); they may even condensate under favorable conditions but the criterium for evaporation or condensation depends on the assumed heating mechanism (see Różańska & Czerny 1999 and the references therein). In another model (Celotti, Fabian & Ress, 1992, Sivron & Tsuruta 1993, Kuncic, Celotti & Rees 1996) very dense blobs confined by magnetic pressure, optically thin to scattering and thick to free-free absorption, form in a spherical relativistic flow.
Within the frame of the cloud scenario, obscuration events are expected. If clouds have a column density in excess of $`10^{25}`$ cm<sup>-2</sup>, the obscuration events would be observed as variability, whatever the size of the clouds (smaller or larger than that of the central source).
In the present paper we discuss the possibility of explaining, at least partially, the variability phenomenon by variable obscuration. We consider the case of random variability due to the statistical dispersion in location of clouds along the line of sight for a constant covering factor. We provide simple analytical estimates of the mean spectral properties, timescales and variability amplitude of AGN, and we support them with computations of radiative transfer done with the use of the codes titan (Dumont, Abrassart & Collin 1999) and noar (Abrassart 2000).
## 2 Model
### 2.1 Cloud model scenario
Accretion flow onto a black hole proceeds most probably down to a hundred Schwarzschild radiiin the form of an accretion disk. Closer in, the relatively cool accretion disk is disrupted and replaced with a mixture of still cold optically thick clumps and hot gas which is responsible for hard the X-ray emission. The cold clumps may not be constrained to the equatorial plane but, under the influence of radiation pressure, may have a quasi-spherical distribution. We envision this general picture in Fig. 1 (see also Fig. 1 in Karas et al. 1999). Such a scenario offers an interesting possibility for explaining the observed spectra of AGN. Although other geometries of the innermost part of accretion flow are not excluded, we concentrate on exploring this particular model.
The mechanism of disk disruption is unknown and therefore the dynamics of the cloud formation process cannot be described. The formation of the hot medium is also not well understood, although it must proceed in some way through cloud evaporation. Therefore we have to resort to a phenomenological parameterization of the cloud distribution and the hot gas geometry. A realistic approach to the cloud model would require a large number of arbitrary parameters reproducing the radial structure and the departure from spherical symmetry. However, the most essential properties of the model can be studied within the frame of a much simpler model.
In the present paper we follow the geometry adopted by Czerny & Dumont (1998) for their model C (see their Fig. 2).
We assume that the hot medium forms a central spherical cloud of radius $`r_X`$. It is characterized by its Thomson optical depth, $`\tau `$, and by its electron temperature, $`T_e`$. Those two parameters uniquely determine the Compton amplification factor $`A(\tau ,T_e)`$ of the hot medium.
The cold clouds are located at a typical distance $`r_{UV}`$ from the center and their distribution is characterized by the covering factor $`\mathrm{\Omega }/4\pi `$. They reprocess the X-ray emission from the hot medium. A fraction of the incident X-ray photons is backscattered, but most of the photons are absorbed. This absorbed radiation is mostly reemitted in UV by the bright side of the clouds, but a fraction is leaking through the dark side. The photons reemitted by the bright side provide seed photons for Comptonization by the central hot medium.
The bremsstrahlung emission from the central comptonizing cloud (CCC) does not significantly contribute to the energy budget, for all the parameters we considered. In this model the prime movers of the engine are the X-rays, in the sense that they are produced in a region (of radius $`r_X`$) close to the black hole, and they drive the reprocessed UV emission further away. But the intrinsic emission of the hot plasma being small, it basically acts as a reservoir of energy, used to upscatter a fraction of the reprocessed UV photons. As the thick clouds can be highly reflective in the UV, this “recycling” of photons can be very efficient when $`\mathrm{\Omega }/4\pi `$ is large. The efficiency of the subsequent production of X-rays is also high if the scattering probability of UV photons, determined by the optical depth of the hot plasma $`\tau `$ and the relative geometrical cross-section $`(r_X/r_{UV})^2`$, is not too low. This geometry may thus imply a stronger coupling between the two emission regions than the alternative disk/corona configuration, where about half of the comptonized power escapes directly.
### 2.2 Radiative transfer computations
An accurate description of the radiative coupling between the hot plasma and the relatively cold clouds requires complex computations of radiation transfer. In particular, the emission from the bright side of the clouds strongly depends on the ionization state of the gas.
We calculate the emission from the bright and dark sides of the clouds in the optical/UV range using the radiative transfer code titan of Dumont, Abrassart & Collin (1999) for Compton thick media. The same code was used by Czerny & Dumont (1998) over the entire optical/X-ray range. However, in the present version we pay more attention to an accurate description of the hard X-ray transfer within the hot plasma and within the surface layers of clouds, therefore we use the Monte Carlo code noar of Abrassart (2000).
In order to obtain a single spectrum model, both codes - titan and noar \- have to be used in a form of iterative coupling, as described in detail by Abrassart (2000). The method is very time consuming, so for the purpose of the present studies of variability, we develop a simple analytical toy model which allows estimates of the observed trends in a broad range of parameters. Therefore, we use the numerical results as a test of the toy model and as a source of information about physically reasonable mean quantities such as the frequency-averaged albedo, the fraction of radiation lost by the dark sides of the clouds, etc.
### 2.3 Toy model
In numerical simulations both the scattered and the reemitted component form a complex reflected spectrum with a shape mostly determined by the ionization parameter
$$\xi =4\pi F_{inc}/n,$$
(1)
where $`F_{inc}`$ is the incident radiation flux and $`n`$ is the number density of a cloud, although the shape of the observed spectrum is even more influenced by the value of the cloud covering factor. In Fig. 2 we show the ratio of the radiation reflected/reemitted by the bright side of the irradiated cloud for a value of the ionization parameter $`\xi `$ equal 300 and a power law spectrum proportional to $`\nu ^1`$ (the ratio depends little on the density and column density, provided the latter is larger than 10<sup>25</sup> cm<sup>-2</sup>, for a range of density between 10<sup>10</sup> and 10<sup>15</sup> cm<sup>-3</sup>). As we see, the true absorption is very low below the Lyman discontinuity, even for such a relatively low $`\xi `$, the reflectivity is very close to unity, and the absorbed and thermalized X-ray flux emitted in the UV creates an excess seen in UV/soft X-rays.
Therefore, due to reprocessing by cold clouds, the initial almost power law distribution of photons produced by Comptonization is replaced by a still broad hard X-ray component and a much more peaked UV/soft X-ray component. The X-ray component is produced by subsequent scattering of photons within the hot cloud, thus forming basically a power law spectral shape with a high energy extension and a slope determined by the properties of the hot cloud ($`\tau ,T_e`$). Cloud emission is closer to a black body and the maximum of the spectrum is mostly determined by the temperatures of the dark and bright sides of the clouds.
Here, for the purpose of an analytical analysis, we simplify this distribution.
The entire optical/UV/X-ray range is schematized by defining two representative energy ranges: soft (UV) or hard (X-ray), with mean energies $`E_{UV}`$ and $`E_X`$. The exact values of those energies are not essential as they do not enter as model parameters through conservation laws (see Appendix A).
We reduce the description of radiative transfer to coefficients which determine the efficiency of the change of an X-ray photon into a UV-photon, and the reverse.
We reduce the photon-frequency-dependent albedo, describing the efficiency of reflection of X-ray photons, to some energy-integrated value, $`a`$. We introduce a frequency-averaged fraction of the X-ray luminosity, $`\beta _d`$, which leaks through the dark side of the clouds in the form of UV emission instead of being reemitted by the bright side.
We also simplify the description of Comptonization by using a mean Compton amplification factor independent of the energy of the input soft photon.
Variability properties are related to the mean spectral shape of an AGN. Therefore we have to introduce a relation between the cloud distribution and the observed spectral shape before discussing the variability itself.
#### 2.3.1 Mean properties of the cloud distribution
We can formulate the conditions for the stationary situation, in which the escape of the photons from the system should be compensated by the production of new photons at the expense of the energy flux supplied to the hot cloud. Noting $`C=\mathrm{\Omega }/4\pi `$ and $`\gamma =\tau /(1+\tau )(r_X/r_{UV})^2`$, the fate of a UV photon at every travel across the region is determined by the following probabilities:
$$(1\gamma )(1C)=P_{esc}^{UV}$$
(2)
$$(1\gamma )C=P_{refl}^{UV}$$
(3)
$$\gamma =P_{ups}^{UV},$$
(4)
where $`P_{esc}^{UV}`$ is the probability of escape from the system, $`P_{refl}^{UV}`$ is the probability of reflection by the clouds and $`P_{ups}^{UV}`$ describes the probability of conversion into an X-ray photon due to the Compton upscattering.
For an X-ray photon we have the following probabilities:
$$1C=P_{esc}^X$$
(5)
$$aC=P_{refl}^X$$
(6)
$$(1\beta _d)(1a)C=P_{abs}^X$$
(7)
$$(1a)C\beta _d=P_{dark}^{UV},$$
(8)
where $`P_{esc}^X`$ is the escape probability, $`P_{refl}^X`$ is the probability of reflection by the clouds, $`P_{abs}^X`$ describes the probability for an X-ray photon to be absorbed by the clouds and reemitted in the form of UV radiation through the bright side, while $`P_{dark}^{UV}`$ describes the probability of absorption of an X-ray photon and subsequent reemission in the form of UV radiation through the dark side.
We have to compensate for the loss of UV photons with new UV photons created in a sequence of events: upscattering a fraction of UV photons by the hot plasma and creating X-rays, absorption of X-rays by the cold clouds, and reemision of this energy in the form of UV photons (see Appendix A):
$$\frac{(P_{esc}^{UV}+P_{ups}^{UV})(P_{esc}^X+P_{abs}^X)}{P_{ups}^{UV}P_{abs}^X}=A(\tau ,T_e).$$
(9)
The stationarity can be clearly achieved only if $`A`$ is greater than unity.
In the present formulation of the model we assumed that the emission from the dark sides of the clouds is provided by X-ray photons. We could instead assume that it is mostly the leaking of UV photons diffusing through an optically thick cloud which powers the dark side emission. This would mean adding the factor $`(1\beta _d)`$ to the left side of Eq. 3, adding a new equation to the set formulated for a UV photon: $`(1\gamma )C\beta _d=P_{dark}^{UV}`$, while putting $`\beta _d=0`$ in Eq. 8, and dropping Eq. 7 in the set formulated for an X-ray photon. These two approaches are almost equivalent unless $`\gamma `$ is very close to 1, corresponding to a change of all UV photons into X-rays.
The ratio of the intrinsic bolometric luminosities of these two spectral components inside the region is given by (see Appendix A)
$$\left(\frac{L_X}{L_{UV}}\right)_{int}=\frac{P_{esc}^{UV}+P_{ups}^{UV}}{P_{abs}^X}=\frac{A\gamma }{1aC}.$$
(10)
This equation shows that for a given system (specified by C, a, and $`\gamma `$), $`\frac{L_X}{L_{UV}}`$ is upper bounded by the limited heat supply to the central comptonizing cloud.
#### 2.3.2 Mean properties of the observed spectrum
The observed spectrum is not identical to the spectrum inside the production region since the escape probability for UV and X-ray radiation differs by a factor of $`(1\gamma )`$ (compare Eqs. 2 and 5), and X-rays transformed into dark emission additionally change the relative proportions.
The observed ratio of the bolometric luminosities in the two components is therefore given by:
$$\left(\frac{L_X}{L_{UV}}\right)_{obs}=\frac{\left(\frac{L_X}{L_{UV}}\right)_{int}(1C)}{(1C)(1\gamma )+\left(\frac{L_X}{L_{UV}}\right)_{int}\beta _d(1a)C}.$$
(11)
This formula simplifies considerably if the contribution of the dark sides of the clouds is neglected:
$$\left(\frac{L_X}{L_{UV}}\right)_{obs}=\left(\frac{L_X}{L_{UV}}\right)_{int}\frac{1}{(1\gamma )}$$
(12)
This is true if the clouds are very thick, as then their dark side temperature is low and the radiation flux peaks in the optical band.
It is worth noting that the observed luminosity ratio can be expressed independently of $`\gamma `$ or $`\beta _d`$:
$$\left(\frac{L_X}{L_{UV}}\right)_{obs}=\frac{A\left(1+C\right)}{1+a\left(A1\right)CAC}.$$
(13)
### 2.4 Model of stochastic variability
We can consider two variability mechanisms expected within the frame of the cloud model and of the geometry of the inner flow.
Random variability at a certain level is always expected since the cloud distribution is assumed to be uniform (spherically symmetric) in a statistical sense. Even without any systematic evolutionary changes of the covering factor $`C`$, our view of the nucleus will undergo variations due to the instantaneous covering factor of the side of the distribution facing the observer.
Systematical trends may also show up if there is a systematic evolutionary change in the mean covering factor, $`C`$.
#### 2.4.1 Random variability amplitude
The most basic parameter of the variability is its amplitude at a given wavelength. In order to reproduce it within the frame of the cloud model we have to assume that clouds can have some random velocities which do not change the mean covering factor but change randomly the number of clouds which occupy the hemisphere just facing us.
If we have $`N`$ clouds, half of them on average occupy the front hemisphere and the dispersion around that mean value would be given by
$$\delta (N/2)=\sqrt{N/2}$$
(14)
so the covering factor as seen by us also displays variations with an amplitude
$$\delta C=C\sqrt{2/N}$$
(15)
Therefore the amplitude of variability of the observed X-ray emission is given by
$$\left(\frac{\delta L_X}{L_X}\right)_{obs}=\frac{C\sqrt{2/N}}{1C}.$$
(16)
High X-ray variability is achieved if the covering factor is high and the number of clouds is not too large.
The UV variability depends again on the optical depth and the contribution of the dark sides predominantly to the UV or optical band.
If dark sides of the clouds are too cool to radiate in the UV, the variability amplitude is given by the same formula as for X-rays
$$\left(\frac{\delta L_{UV}}{L_{UV}}\right)_{obs}=\frac{C\sqrt{2/N}}{1C}.$$
(17)
However, if the dark sides contribute to the UV we have
$$\left(\frac{\delta L_{UV}}{L_{UV}}\right)_{obs}=\frac{C\sqrt{2/N}[(1\gamma )(1aC)\gamma A(1a)\beta _d]}{(1\gamma )(1C)(1aC)+\gamma A(1a)C\beta _d}$$
(18)
which reduces the UV relative amplitude.
Such normalized variability amplitudes depend on the number of clouds, $`N`$, constituting the UV emitting region. This means that an additional free parameter is involved.
The normalized variability amplitudes can be determined observationally. Studies of variability provide us either with the rms value in a given spectral band, or just with a typical value of the variability factor. Quoted rms values can be directly identified with our normalized amplitude. If instead only a variability factor is given, we can understand it as statistical variations at 2 standard deviation level. This means that we actually observe variability of a factor $`𝒜`$ when the luminosity changes from $`L_X2\delta L_X`$ to $`L_X+2\delta L_X`$.
Therefore this variability factor can be expressed as
$$𝒜=L_X(max)/L_X(min)=\frac{1+2\left(\frac{\delta L_X}{L_X}\right)_{obs}}{12\left(\frac{\delta L_X}{L_X}\right)_{obs}}.$$
(19)
We can also define the ratio $`R`$ as the ratio of the normalized variability amplitudes in the X-ray band and in the UV band:
$$R=\frac{\left(\frac{\delta L_X}{L_X}\right)_{obs}}{\left(\frac{\delta L_{UV}}{L_{UV}}\right)_{obs}}.$$
(20)
Such a ratio does not depend on the number of clouds, $`N`$, constituting the UV emitting medium. Therefore this ratio is fully determined by the parameters of our toy model.
#### 2.4.2 Random variability timescales
The random variability we consider results from the clouds passing through our line of sight to the central source. On the one hand, each passing cloud produces an eclipse phenomenon. On the other hand, the dispersion in the cloud velocities, due to even a small dispersion of the distances from the gravity center, results in variations of the cloud distribution, as described in Sect. 2.4.1.
Let us estimate a representative duration $`t_{dip}`$ of an obscuration event. This is determined by the size and velocity of a cloud as well as the size of the X-ray source and the entire region involved. The cloud velocity is of the order of the keplerian velocity so it is determined by the mass of the black hole (or its gravitational radius) and $`r_{UV}`$. A typical size of a cloud is given simply by the number of clouds, the radius $`r_{UV}`$ and the covering factor, if there is at most one cloud in the line of sight, i.e. clouds are not overlapping.
The mean number of clouds in a line of sight is of the order of unity (cf Appendix B) so we get the typical cloud size:
$$r_{cl}^2\frac{C}{N}r_{UV}^2.$$
(21)
Here we have ignored numerical factors of the order of $`\pi `$, since this expression, and the following formulae, serve only as an order of magnitude estimate.
Since the clouds are ten or more times smaller than the region involved, the fastest variation corresponds to an ingress or an egress from a single eclipse
$$t_{min}\frac{r_{Schw}}{c}\sqrt{\frac{4C}{N}(\frac{r_{UV}}{R_{Schw}})^3}.$$
(22)
Typical random rearrangement of the cloud distribution will proceed on a dynamical timescale connected with Keplerian motion
$$t_d\frac{r_{Schw}}{c}\sqrt{(\frac{r_{UV}}{r_{Schw}})^3}.$$
(23)
If the cloud distribution covers a range of radii, the longest timescale involved will be given by $`t_d`$ at the outer edge of the cloud distribution.
#### 2.4.3 Variations of the covering factor in the toy model
This kind of variability is more complex to study since variations in the covering factor lead not only directly to variations of $`L_X`$ and $`L_{UV}`$ but also result in a change of the mean number of clouds along the line of sight, and, under an assumed constant total luminosity, in a change of the ionization state of the clouds, i.e. albedo. Such trends can only be studied by solving the radiative transfer equation for a sequence of models, which we postpone to a future paper.
### 2.5 Relation between toy model parameters and the ionization parameter
The spectral features in X-ray band allow us to have an insight into the ionization state of the reprocessing medium. This ionization state is mostly determined by the value of the ionization parameter $`\xi `$. Since we frequently have some direct estimates of this parameter, it is interesting to relate this parameter to the parameters of our toy model. It will allow for simple estimates of those parameters through the quantities which can be estimated on the basis of observational data. That way, we can judge to some extent the applicability of the model.
For a bolometric luminosity of an object, $`L`$, the incident flux determining the ionization parameter, $`\xi `$, (Eq.1) in the case of the quasi-spherical distribution of clouds is:
$$F_{inc}=\frac{L}{r_{UV}^2},$$
(24)
where $`n`$ is the number density of the clouds and $`r_{UV}`$ is the representative distance of a cloud from the black hole.
The most convenient form of relating the number of clouds to the ionization parameter and other observables is through the explicit dependence on the covering factor and the column density of the clouds, $`N_H`$. We can derive the appropriate formula knowing that the typical size of the cloud $`r_{cl}`$ is related to $`N`$, the covering factor and the characteristic radius of the cloud distribution (see Eq. 21).
The obtained relation
$$NL^2CN_H^2r_{UV}^2\xi ^2.$$
(25)
indicates that the number of clouds scales with the square of the luminosity to the Eddington luminosity ratio, $`R_{Edd}`$, if $`N_H`$, $`\xi `$ and $`r_{UV}`$ (expressed in Schwarzschild radii) are similar in all objects. It means that high luminosity, low covering factor objects should be the least variable. If this estimate of the number of clouds is inserted into Eqs. 17 and 16 we see that we do not expect a direct scaling with the central mass of the black hole or bolometric luminosity. The very weak trends between the UV variability and the bolometric luminosity observed in the data (Paltani & Courvoisier 1997, Hook et al. 1994) might be caused by some secondary coupling between the mass of the black hole and the model parameters.
The size and the number of clouds is constrained by the dimension of the emission region, as the volume filling factor has to be smaller than unity:
$$\frac{Nr_{cl}^3}{r_{UV}^3}1,$$
(26)
or:
$$N_H\xi r_{UV}L^11.$$
(27)
One can now illustrate the properties of the cloud model with some numbers.
Expressing the mass of the black hole in 10<sup>8</sup>M, $`R_{Edd}`$ in 10<sup>-1</sup>, $`N_H`$ in 10<sup>26</sup> cm<sup>-2</sup>, $`\xi `$ in 10<sup>3</sup> and $`r_{UV}`$ in 10 $`r_{Schw}`$, one gets, using Eqs. 21 and 25:
$$r_{cl}2\times 10^{12}N_{26}\xi _3\left(\frac{r_{UV}}{10r_{Schw}}\right)^2M_8\left(\frac{R_{Edd}}{0.1}\right)^1\mathrm{cm}$$
(28)
$$n_H5\times 10^{13}\xi _3^1\left(\frac{r_{UV}}{10r_{Schw}}\right)^2M_8^1\frac{R_{Edd}}{0.1}\mathrm{cm}^3$$
(29)
$$N10^4CN_{26}^2\xi _3^2\left(\frac{r_{UV}}{10r_{Schw}}\right)^2\left(\frac{R_{Edd}}{0.1}\right)^2,$$
(30)
and the condition on the filling factor:
$$10^3CN_{26}\xi _3\frac{r_{UV}}{10r_{Schw}}\left(\frac{R_{Edd}}{0.1}\right)^11,$$
(31)
Clouds with such properties could form for instance as broken pieces of the inner accretion disk.
Using Eqs. 25 and 16 the X-ray variability amplitude is:
$$\left(\frac{\delta L_X}{L_X}\right)_{obs}\frac{\sqrt{C}}{(1C)}L^1N_Hr_{UV}\xi $$
(32)
It is maximum for a covering factor equal to 0.5, corresponding to two clouds along the line of sight, on average.
Using the same notation as in Sect. 2.5, the X-ray variability amplitude is:
$`\left({\displaystyle \frac{\delta L_X}{L_X}}\right)_{obs}`$ $``$ $`\sqrt{2}\times 10^2{\displaystyle \frac{\sqrt{C}}{(1C)}}N_{26}\xi _3`$
$`\left({\displaystyle \frac{R_{Edd}}{0.1}}\right)^1{\displaystyle \frac{r_{UV}}{10r_{Schw}}}.`$
and the time scale of the variability, using Eq. 22:
$$t_{dip}10^2N_{26}\xi _3\left(\frac{R_{Edd}}{0.1}\right)^1M_8\mathrm{sec}$$
(34)
We see that a large variability amplitude in a very small time scale is easily obtained in this model.
We cannot expect, however, that the overall variability properties of AGN may be explained by eclipsing clouds all being at the same distance and having the same radius. Actually, any more realistic picture would require clouds to occupy a range of radii.
## 3 Results
The complete toy model of the stationary distribution of clouds depends on four parameters: the covering factor, $`C`$, the probability of upscattering a $`UV`$ photon into an X-ray photon, $`\gamma `$, the X-ray albedo of the bright side of the clouds, $`a`$, and the fraction of the X-ray radiation leaking through the dark side of the clouds in the form of $`UV`$ radiation, $`\beta _d`$. These four values determine the Compton amplification factor of the hot plasma, $`A`$, through the conservation law given by Eq. 9. In full radiative transfer numerical models the number of free parameters is the same: $`C`$, $`r_X/r_{UV}`$, $`r_X`$, and the hot plasma parameters $`\tau `$ and $`T_e`$. Fortunately the number density and the column density of the cloud system play a small role. However the freedom in modeling is considerable. Therefore, in order to study the cases which are of direct interest to observed AGN, we consider first in some detail the mean Seyfert spectrum in order to fix some of these parameters to be used in most of our further considerations.
### 3.1 Toy model of the mean properties of Seyfert 1 galaxies
The broad band spectra of Seyfert 1 galaxies are relatively well known. Most objects have an intrinsic hard X-ray slope of about $`0.9`$ if the spectra are corrected for the reflection component (Nandra & Pounds 1994). However, determination of the shape of the high energy tail of the spectrum of a single object poses a considerable problem. Relatively strong constraints on the spectrum extension for typical Seyfert galaxies can be only achieved through the analysis of composite spectra, i.e. summed spectra of a number of objects. Such a composite spectrum was published by Gondek et al. (1996) on the basis of combined Ginga/OSSE data for 7 Seyfert 1 galaxies. We use this spectrum to reduce the degrees of freedom in our models.
#### 3.1.1 Compton amplification factor
A relatively low dispersion in the observed X-ray slopes of Seyfert galaxies around the mean value 0.9 allows us to estimate the characteristic Compton amplification factor. We use for this purpose the Monte Carlo results of Janiuk, Życki & Czerny (2000).
The Comptonization process by a hot spherical cloud depends on two parameters: optical depth of the cloud, $`\tau `$, and electron temperature, $`T_e`$, and leads to a broad range of X-ray slopes. However, for each value of $`T_e`$ we can find an optical depth which gives a spectral slope of 0.9. Therefore we finally obtain the value of the Compton amplification factor as a function of $`T_e`$ only. The corresponding plot is shown in Fig. 3.
Determination of the value of $`T_e`$ from the observed high energy cutoff in the X-ray spectra is rather uncertain. Gondek et al. (1996) give the value of 260 keV for the combined Seyfert 1 spectrum from the Ginga/OSSE data. Fortunately, the dependence of $`A`$ on $`T_e`$ is weak. This is not surprising since both the spectral slope and the Compton amplification factor are mostly determined by the value of the Compton parameter $`y`$.
Adopting Te = 260 keV, we fix the value of the Compton amplification factor $`A=4.0`$ in further considerations (Fig. 3). It reduces the number of free parameters of our toy model to 3.
#### 3.1.2 Expected trends in stationary parameters
Fixing the value of the Compton amplification factor leads to a relation between the original model parameters: $`C`$, $`a`$, $`\beta _d`$ and $`\gamma `$. It leaves a choice of 3 parameters out of 4, as basic independent model parameters. We analyse the relation in order to make the most convenient choice. We also study the dependence of the predicted luminosity ratios on our selected set of parameters.
Negligible dark side contribution
We consider first the case of negligible contribution of the dark sides of the clouds to the UV spectrum, due to their large optical depth ($`\beta _d`$ = 0).
The stationarity condition allows to express the probability of scattering by the hot cloud, $`\gamma `$, as the function of the X-ray albedo, $`a`$, and covering factor, $`C`$. We show this dependence in Fig. 4.
Factor $`\gamma `$ depends mostly on the covering factor but it also shows a trend to increase with the albedo (i.e. the ionization parameter). Only large values of the covering factor are allowed, as the optical depth of the hot cloud is small and the physical size of the X-ray region $`r_X`$ cannot be larger than $`r_{UV}`$ within the frame of our model. So $`\gamma `$ cannot be larger than $`\tau `$. It means that $`C`$ cannot be smaller than about 0.7 for $`\tau =0.1`$ although larger optical depth allow for smaller covering factors.
The role of the dark side contribution
The contribution from the dark sides of clouds may not, however, be negligible. Therefore, it is more appropriate to present the results by constraining the albedo. For a broad range of the ionization parameter $`\xi `$, the radiative transfer computations (Dumont & Abrassart 2000) show that the X-ray albedo, $`a`$, for partially ionized gas is about 0.5. If we use this constraint, it leaves us a choice of two independent parameters among $`C`$, $`\gamma `$ and $`\beta _d`$.
We plot the dependence of $`\beta _d`$ on the covering factor $`C`$ and the scattering probability $`\gamma `$ in Fig. 5. We see that $`\beta _d`$ should be relatively important for large covering factors and large $`\gamma `$. There is also a range of parameters (small $`C`$ and $`\gamma `$) for which the formally computed value of $`\beta _d`$ is negative (negative values are not plotted on the figure). Using $`C`$ and $`\gamma `$ as independent variables, we have to be aware of this unphysical range of parameters which depends to some extent on the adopted albedo. If we choose another set of two arbitrary parameters the unphysical range does not disappear.
The cross-section of Fig. 5 for constant $`\beta _d`$ produces $`\gamma `$ closely correlated to $`C`$. Such a relation for $`\beta _d=0.2`$ (it is a high value for the high optical thickness of our clouds, according to the transfer computations of Dumont & Abrassart, 2000) is rather similar to results for $`\beta _d=0`$.
#### 3.1.3 Expected trends in the luminosity ratios
The observed ratio of the X-ray to the UV luminosity depends only on the covering factor if the X-ray albedo and Compton amplification factor are fixed; it is not influenced by the second model parameter ($`\gamma `$ or $`\beta _d`$) within the frame of our toy model (see Eq. 13), in opposite to other quantities.
This quantity is one of the important observables, although in practice the determination of the $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$ ratio may not be easy because of extinction. It is also of major importance in all accretion models, so we have to study its dependence on the parameters which we usually fix, namely Compton amplification factor $`A`$ and X-ray albedo $`a`$.
In Fig. 6 we show the $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$ ratio. We see that this ratio mostly depends on the covering factor. For an albedo $`a=0.15`$ (cold matter) the X-ray luminosity starts to dominate the bolometric output only if the covering factor is smaller than 0.6. In the case of a larger albedo of partially ionized gas, $`a=0.5`$, the X-ray emission dominates the UV luminosity up to about $`C=0.75`$.
We see from this plot that $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$ is smaller than unity for the range of covering factors we consider. This is a very important characteristic of this model: it is able to account for a low X-ray to UV luminosity ratio without any ad-hoc hypothesis.
Fig. 7 shows the dependence of $`\left(\frac{L_X}{L_{UV}}\right)_{obs}/\left(\frac{L_X}{L_{UV}}\right)_{int}`$ on the covering factor $`C`$ and $`\gamma `$ for $`a=0.5`$. This ratio, in contrast to $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$, depends on $`\gamma `$ as well, not just on the covering factor.
Over most of the allowed parameter space, $`\left(\frac{L_X}{L_{UV}}\right)`$ is slightly lower inside the clouds system than what is observed. It is also interesting to note that the X-ray to UV luminosity ratio seen outside the cloud system can be larger than this ratio determined inside the medium, or in other words that inside the medium the radiation field is softer than the observed spectrum. The effect is stronger towards small values of $`\gamma `$, but the parameter space where this happens is rather narrow, because the highest values happen where there are no physical solutions. Indeed, the energy losses through the dark sides of the clouds take negative values for a large range of covering factor and probability of scattering by the hot plasma (see Fig. 5).
The reason for this can be more clearly seen on Fig. 8, for which we relaxed the A=4 assumption, but fixed $`\beta _d`$. This last parameter does not qualitatively change the behavior of the equilibrium amplification factor but simply forces it to higher values when it raises. One sees that fixing A defines a unique relation in the C,$`\gamma `$ plane, but that there are no solutions on either side of this line. By raising $`\beta _d`$, it is possible to find an equilibrium on the side where A is too high, but the side where A is not sufficient is clearly forbidden.
The accuracy of the determination of the Compton amplification factor does not influence much the X-ray to UV luminosity ratio if the covering factor is large. This dependence is plotted in Fig. 9, for an X-ray albedo of one half. Higher albedos (i.e. ionization state) require higher covering factors, whatever the amplification factor, in order to keep a reasonably low X-ray to UV luminosity ratio.
#### 3.1.4 Toy model for the Gondek et al. spectrum
We can now find model parameters which reproduce the shape of the mean Seyfert spectrum of Gondek et al. (1996). We need, for this purpose, the observed $`(\frac{L_X}{L_{UV}})`$ ratio, the X-ray albedo and the contribution from the dark sides $`\beta _d`$. Also the optical depth of the hot plasma (or its temperature) is required in order to obtain the relative size of the cloud distribution to the radius of the hot plasma, $`r_X/r_{UV}`$.
We have no accurate measurements of the X-ray to the UV bolometric luminosity ratio for this sample since the UV data are very sensitive to even minor absorption, and the extension of this spectral component into the EUV is difficult to incorporate into the bolometric luminosity of the soft component. The $`\nu F_\nu `$ fluxes measured at 1375 Å and 2 keV given by Walter & Fink (1993) can be used to make a rough estimate. Of the sources included in the combined Seyfert 1 spectrum by Gondek et al. (1996), two are heavily absorbed in the UV: MCG 8-11-11 and MCG-6-30-15 so they have to be rejected. The remaining five sources have the average ratio $`\nu F_\nu `$ at 1375 Å to $`\nu F_\nu `$ at 2 keV about 14, with the highest value being $`20`$. Since the X-ray component is much broader, we have to apply a bolometric correction $`K`$ to it in order to roughly reproduce the broad band luminosity ratio. We assume $`K=4`$ so we obtain $`(L_X/L_{UV})_{obs}=K/14=0.3`$. This is only a rough estimate; the result of detailed analysis would depend significantly on the temperature of the $`UV`$ component.
Such an observed luminosity ratio can be used to estimate the parameters of the cloud distribution. Assuming the value of albedo for partially ionized matter $`a=0.5`$, $`A=4`$, $`(L_X/L_{UV})=0.3`$ and neglecting the contribution from the dark sides ($`\beta _d=0)`$ we obtain the covering factor $`C=0.87`$ and $`\gamma `$ equal 0.052 (for $`a=0.3`$, $`\gamma `$ is the same, and $`C=0.91`$). The value of $`\gamma `$ can be used to determine the $`r_X/r_{UV}`$ if we know the optical depth of the hot cloud. Taking the optical depth $`\tau =0.1`$ after Gondek et al. (1996) we obtain $`r_X/r_{UV}=0.72`$. However, observational determination of $`\tau `$ strongly depends on the (rather poor) accuracy of the determination of the high energy cut-off. Assuming a larger value of $`\tau =0.3`$ we obtain $`r_X/r_{UV}=0.42`$. We therefore see that we can determine reliably the parameter $`\gamma `$, but the factorization of its value between the physical size of the hot cloud and the optical depth is only weakly constrained by the details of the spectral shape.
Allowing for a relatively high ratio of energy leaking through the dark sides, $`\beta _d=0.1`$, we obtain: $`C=0.93`$, $`\gamma =0.03`$ which translates to $`r_X/r_{UV}`$ equal 0.55 for $`\tau =0.1`$ and 0.32 for $`\tau =0.3`$.
### 3.2 Random variability
Detailed observations of the variability pattern are available only for a few sources and the observed trends are possibly not characteristic of all AGN as a whole. Therefore we first outline the most basic trends and later on we apply the toy model to the best monitored sources in order to check whether cloud obscuration may account for their specific behavior.
#### 3.2.1 X-ray variability amplitude
The X-ray variability amplitudes predicted by the cloud model can be easily estimated from Eq. 16 in the case of no leakage from the dark side.
A large covering factor leads to large variability even if the number of clouds is large. For example, 1000 clouds covering 0.9 of the source produce a maximum to minimum flux ratio 9 (see Sect. 2.4.1). This factor reduces to only 1.2 for covering factor 0.5; only considering significantly smaller number of clouds would increase the variability.
The mean Seyfert 1 spectra analysis suggested that the typical covering factor is about 0.9 (see Sect. 3.1). Those objects display an X-ray variability by a factor 2 on average. This means that the required number of clouds is of the order of a few thousands.
#### 3.2.2 X-ray/UV relative amplitude
The ratio $`R`$ of the normalized variability amplitudes in the X-ray band to that in the UV band is fully determined by the parameters of the toy model (see Eq. 20).
This ratio is equal to unity in the toy model as long as the contribution of the dark sides of the clouds is negligible (see Eqs. 16, 17, 18).
When the contribution from the dark sides of clouds is allowed, the relative amplitude in X-ray and $`UV`$ is reduced. We show this trend in Fig. 10 plotting $`R`$ against $`\beta _d`$ and $`C`$, assuming $`a=0.5`$, $`A=4`$. The value of $`R`$ is always greater than 1 and even very high values are allowed if $`\beta _d`$ is large.
#### 3.2.3 Application of the toy model to monitored objects
A few AGN were recently monitored both in the UV and X-ray bands, and the results of these campaigns can be used directly to determine the properties of the cloud distribution in those sources within the frame of our random variability picture.
Observational data provides us either directly with the rms values both in UV and X-rays or with the value of variability factor $`𝒜`$ in those bands. Data were taken from Goad et al. (1999) and Edelson & Nandra (1999) for NGC 3516, from Nandra et al. (1998) for NGC 7469, from Edelson et al. (1996) for NGC 4151 and from Clavel et al. (1992) for NGC 5548.
Observational estimation of the ratio of the X-ray to UV luminosity is unfortunately rather complex. Most sources have red optical/UV spectra, with negative slopes on a $`\nu F_\nu `$ plot. It most probably means that the objects are considerably reddened, with a significant fraction of the energy reemitted in the IR (e.g. Wilkes et al. 1999). A conservative discussion of this problem for NGC 4151 by Edelson et al. (1996) concluded that X-ray luminosity is 3 times smaller than the UV/optical/IR luminosity. The situation for the other objects looks qualitatively similar. Also this ratio determined for the Gondek et al. (1996) composite is of the same order (i.e. 0.3). We therefore adopted $`L_X/L_{UV}`$ ratio equal 1/3 for all objects.
We now fix two of the toy model parameters: the Compton amplification factor $`A=4`$ and X-ray albedo $`a=0.5`$. Now, two observed quantities (normalized variability amplitudes in the UV and X-ray bands) allow us to calculate the covering factor $`C`$, the number of clouds, $`N`$, the probability of upscatering for a UV photon $`\gamma `$, and the contribution of the dark sides of clouds, $`\beta _d`$. The results for the four selected objects are given in Table 1.
The covering factor obtained is the same for all objects, as it is determined purely by the X-ray to UV luminosity ratio. The assumption of $`L_X/L_{UV}=1/3`$ influences the obtained value of the covering factor, but not strongly. A value of 0.5 would give $`C=0.86`$ and a value of 0.1 would give $`C=0.96`$. The last value is more appropriate for quasars than for Seyfert 1 galaxies. Also the adopted value of the albedo does not influence the results significantly: $`a=0.3`$ would give a covering factor of 0.86 for our objects.
The value of $`\beta _d`$ is equal to zero in two of our objects, since the normalized variability in the UV and X-rays are equal and no contribution from the dark sides is required.
### 3.3 Mean spectrum from radiative transfer codes
Since the complete solution of radiative transfer equation within the cloud system are extremely time consuming, we used the toy model results to guide our choice of the model parameters. Now we can test our toy model parameterization against numerical results.
The present version of the radiative transfer computations does not include yet the finite size of the hot Comptonizing medium located close to the black hole, so the effect of Comptonization is replaced by a central point source emitting a power law continuum of an arbitrary slope and extension.
The mean spectrum from our computation is shown in Fig. 11, along with the parameters used.
Spectral features in the UV and X-ray band are clearly visible, including iron $`K_\alpha `$ line and a number of emission lines in soft X-rays. These features are due to ’reflection’ of X-rays by partially ionized clouds.
We can now estimate from the computations some of the toy model parameters.
The ratio of the X-ray to UV luminosity can be now calculated from the model, assuming a certain division between these two energy bands. Looking at the Fig. 11 we clearly see the two-component character of the spectrum, with a big blue bump extending roughly to 50 eV, and an X-ray component of roughly a power law shape. Therefore we adopt 50 eV as a division point between the UV and X-ray bands.
The observed $`L_X/L_{UV}`$ ratio depends on the adopted size of the central cloud. If this size is neglected, as in this computation, the $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$ is equal to 0.24. The assumption of $`r_X/r_{UV}=0.7`$ reduces it to 0.22, so the effect of this uncertainty on the spectral appearance is not essential.
The frequency-averaged albedo in the full numerical computations, calculated as a ratio of the X-ray luminosity to the total incident luminosity, is equal to 0.58, in agreement with the values used to construct Fig. 6-8.
We can also calculate the effect of leaking through the dark sides of the clouds. For the presented computations $`\beta _d=0.26`$.
This value is much higher than in three out of the four monitored objects. It means that in Seyfert galaxies like NGC 5548, the column density of clouds should be larger than $`N_H=10^{26}`$ cm<sup>-2</sup> adopted in the computations. Therefore we compare the mean observed spectrum of NGC 5548 with the computed one, neglecting the contribution from the dark side of clouds, i.e. $`\beta _d=0`$ (see Fig. 12). The predicted spectrum may still be too bright in far UV range, which suggests that the ionization parameter adopted in the numerical computations was slightly too low. Also, the high energy cut-off adopted in the computations (100 keV) should be increased in order to fit this particular data. When the dark side contribution is neglected, the obtained $`\left(\frac{L_X}{L_{UV}}\right)_{obs}`$ ratio is higher, equal to 0.43.
From the toy model, if we assume $`C=0.9`$, $`a=0.58`$ and $`A=4`$ we obtain $`\left(\frac{L_X}{L_{UV}}\right)_{obs}=0.38`$. Therefore, the toy model reproduces surprisingly well the complex solutions of the radiative transfer if supplemented with appropriate values for the X-ray albedo and dark side emission efficiency.
### 3.4 Random variations from the radiative transfer codes
Random rearrangement of the clouds leads to minor changes of the effective covering factor in the direction towards the observer. We illustrate here these changes using the results of the numerical computations described in detail in the previous section.
Assuming the number of clouds, $`N=1000`$, we show two representative examples of the observed spectra at any moment of time (see Fig. 11). Since the contribution of the dark sides was non-negligible for the adopted parameter set, the variations seen in UV are of much lower amplitude than the variations of the X-ray emission.
## 4 Discussion
### 4.1 Quasars, Seyfert 1 galaxies and Narrow Line Seyfert 1 galaxies
Various classes of AGN differ systematically although it may not necessarily mean that they form truly separate classes. Instead, they rather represent various parts of the same continuous multidimensional distribution in some parameter space. Interesting approach to this problem was formulated within the frame of the method of principal component analysis (PCA) by Boroson & Green (1992) and Brandt & Boller (1998).
Our cloud model can well reproduce the typical trends. Stronger Big Blue Bump component characteristic of quasars and narrow line Seyfert 1 galaxies (NLS1) corresponds to a larger covering factor, of order of 0.98, instead of 0.9 obtained for Seyferts (see Sect. 3.2.3). Stronger variability of NLSy1 galaxies in comparison with quasars (see Leighly 1999) means that the number of clouds in these objects is of order of $`10^4`$, not much higher than for Seyfert 1 galaxies but an order of magnitude lower than in quasars.
These trends are not surprising. If we imagine for simplicity that all accreting clouds are identical, the number of clouds present would depend on the accretion rate and the travel time of clouds scales with the mass of the black hole, $`N\dot{M}M`$. The covering factor would be given by the number of clouds and the size of the typical radius $`r_{UV}`$, scaling again with the mass, i.e. $`CN/M^2\dot{m}/M`$. This would mean that the covering factor $`C`$ is mostly determined by the luminosity to the Eddington luminosity ratio and both quantities are high in quasars and NLSy1 galaxies, while additionally the number of clouds depends on the mass and accretion rate, causing fainter objects to vary more than brighter objects with the same luminosity ratio.
### 4.2 Relative normalized amplitudes in X-rays and UV
Our model of variability predicts that the normalized amplitude of variations in X-rays should be equal to that in the UV if the contribution of the dark sides of the clouds to the UV emission can be neglected. Therefore, equality of these two amplitudes in a number of objects is expected. In those objects which show stronger variability in the X-ray band than in the UV band, a certain level of contribution from the dark sides is required.
This kind of behavior is quite in contrast with the frequently adopted picture in which the UV variability is caused by reprocessing of the variable X-ray flux. In this case the basic effect is in the change of the temperature of the illuminated cool gas. Equality of the normalized variability amplitudes in X-rays and in the UV requires fine tuning. We can estimate the effect in the following way.
In Fig. 13 we show the result of complete thermalization of the X-ray flux varying by a factor 2, described using a simple black body approach. The variability in the UV (1315 Å) depends on the mean value of the cold matter temperature $`T_{cl}`$. $`R(T_{cl})`$ is close to 1 only for a single specific value of the mean cloud temperature ($`27000`$ K). For colder matter, the UV variability is stronger since we are at the exponential part of the Planck function. Hotter clouds exhibit variations with an amplitude four times smaller than the amplitude in the X-rays which caused variations.
Fine resolution monitoring of the Seyfert galaxy NGC 7469 in the UV and X-ray bands show that, at least in this source, the relative amplitudes in the UV and X-ray bands are the same. It is easily explained within the frame of the cloud model by adopting a relatively low contribution of the dark sides of the clouds to the UV band, i.e. rather high column density of the clouds. In the case of the simple reprocessing picture, the same value for both normalized amplitudes is difficult to obtain (see Fig. 13). The typical behavior of the other three sources considered in Sect. 3.2.3 also conforms to the cloud model expectations, with one more source showing the same amplitude ratio (NGC 5548). Unfortunately, the number of objects monitored both in the UV and X-rays is still low.
Within the frame of our model, large amplitude of variations should be unavoidably accompanied by occasional complete coverage of the X-ray source due to large statistical deviations. Such events are actually observed; the best documented case was seen in NGC4051 (Guainazzi et al. 1998, Utley et al. 1999).
Surprisingly low amplitude in the optical band, or no variations seen in some sources strongly variable in X-ray band (see Miller 1999,: also Boller, privite communication), can be also accommodated by our model as it may correspond to a very large contribution of the dark sides of clouds. However, this effect may be also due to the domination of the optical band by the starlight of the host galaxy.
Our obscuration model of variability cannot explain those events which are characterized by larger amplitude variations in the UV than in the X-ray band. Such events are occasionally observed in some sources. For example, NGC 5548 showed a spectacular brightening in the UV in 1984 which was accompanied by a very moderate brightening in X-ray band (see Clavel et al. 1992). The nature of such long timescale events should be definitively different from random variability of the covering factor discussed in this paper.
### 4.3 Hard X-ray slope and high energy cut-off
The prediction of the obscuration model of variability is that nothing actually changes within the source. Random rearrangements of clouds blocking our line of sight to the X-ray emitting region do not change either the optical depth, or the temperature of the hot plasma. The amount of reflection also should in principle not vary with respect to the primary emission.
Observational constraints on the variability of the slope of the X-ray ’primary’ component are not conclusive. A number of studies suggest that strong variations in the luminosity are not accompanied by spectral changes (e.g. Turner, George & Netzer 1999 for Akn 564, George et al. 1998 for NGC 3227; see also Gierliński et al. 1997 for Cyg X-1). However, Zdziarski, Lubiński & Smith (1999) and Done, Madejski & Życki (1999) suggest the presence of a correlation between the slope of the primary component, amount of reflection and the source luminosity for NGC 5548. The problem of decomposition of the spectrum into two components - primary and reflected - is difficult for the short time sequences of the data used in variability studies. Also, a fraction of the reflection comes from very large distances and does not respond to the nuclear emission within the observed time, as suggested by the remnant emission observed during the off state in NGC 4051 (Guainazzi et al. 1998)
The prediction of the model that no spectral variability in X-rays is expected is also not firm. In a real situation, if there are any temperature gradients within the hot cloud and the cold clouds are located on a range of distances with a complex overlapping pattern, we can expect some weak spectral variations caused by cloud eclipses. However, the predictions of the trends would require either an ad hoc parameterization, or a dynamical study of cloud formation and disruption which is beyond the scope of the present paper.
## 5 Conclusion
The cloud scenario offers an interesting quasi-spherical model of accretion onto a central massive black hole. It explains the observed large ratio of the Big Blue Bump luminosity to the X-ray luminosity. The second prediction of the model is the variable obscuration of the X-ray source. It can explain the following observed trends:
* $`\mathrm{}`$ significant X-ray variability not accompanied by the change of the spectral slope; this is due to random cloud redistribution without changes of the covering factor,
* $`\mathrm{}`$ amplitude of the variability in X-ray band larger or equal to amplitude in UV band: this is due to the contribution from the dark sides of clouds,
* $`\mathrm{}`$ variability timescales ranging from $`10^2`$ s to $`10^6`$ s; this is due to the size of the optically thick clumps, the size of the X-ray medium and Keplerian motion.
Variability analysis indicates that in many objects the column density of the clouds is very large, $`N_H>>10^{26}`$cm<sup>-2</sup> suggesting their origin in violent disk disruptions.
Further progress is required in order to incorporate the finite size of the hot central cloud into the numerical computations of the radiative transfer, and to address the problem of intrinsic variability.
###### Acknowledgements.
We are grateful to Suzy Collin for many extensive discussions and suggestions, to Anne-Marie Dumont for participation in computing the numerical model and for helpful comments to the manuscript, and to Katrina Exter for her help with English. We also thank Dirk Grupe, our referee, for his help in improving the presentation of the paper and for pointing out valuable references. Part of this work was supported by grant 2P03D01816 of the Polish State Committee for Scientific Research and by Jumelage/CNRS No. 16 “Astronomie France/Pologne”.
## Appendix A : Energy conservation and Compton amplification factor
We denote the number of UV photons inside the radius $`r_{UV}`$ by $`N_{UV}`$, the number of X-ray photons by $`N_X`$, and their mean energies by $`E_{UV}`$ and $`E_X`$, correspondingly. We also introduce the efficiency $`\eta _X`$ of creating an X-ray photon of an energy $`E_X`$ during subsequent scatterings within the hot plasma and the efficiency $`\eta _{UV}`$ of creating UV photons from an absorbed X-ray photon of energy $`E_X`$. The value of $`E_X`$ and $`\eta _X`$ depend sensitively on the hot cloud optical depth $`\tau `$ and its electron temperature $`T_e`$ since they are related to the spectral shape of produced X-rays. $`\eta _{UV}`$ is simply given by
$$\eta _{UV}=E_X/E_{UV}.$$
(35)
The conservation law for the number of UV photons within the system is given by
$$(P_{esc}^{UV}+P_{ups}^{UV})N_{UV}=\eta _{UV}P_{abs}^XN_X$$
(36)
where all the probabilities are given by Eqs. (1) - (7).
The conservation law for a number of X-ray photons is given by
$$(P_{esc}^X+P_{abs}^X+P_{dark}^{UV})N_X=\eta _XP_{ups}^{UV}N_{UV}$$
(37)
Those two equations can be combined into the condition of the stationarity
$$(P_{esc}^{UV}+P_{ups}^{UV})(P_{esc}^X+P_{abs}^X+P_{dark}^{UV})=\eta _X\eta _{UV}P_{ups}^{UV}P_{abs}^X$$
(38)
If true losses from the system (i.e. leak through the dark sides and escape) are equal zero, $`\eta _{UV}\eta _X`$ must be equal 1, i.e. no Compton amplification is possible (or required) and the hot and cold plasma achieve thermal equilibrium. Only stationary losses and energy supply to the hot plasma support the existence of the two strongly different media.
The product $`\eta _{UV}\eta _X`$ is directly related to the Compton amplification factor $`A`$ of the hot plasma
$$\eta _{UV}\eta _X=A(\tau ,T_E)$$
(39)
Substituting this relation into Eq. 38 we obtain the stationarity condition 9.
We can also determine the ratio of the bolometric luminosity in X- ray and UV spectral components.
$$\left(\frac{L_X}{L_{UV}}\right)_{int}=\frac{N_XE_X}{N_{UV}E_{UV}}.$$
(40)
The $`N_X/N_{UV}`$ ratio can be determined from Eq. 36
$$\left(\frac{L_X}{L_{UV}}\right)_{int}=\frac{(P_{esc}^{UV}+P_{ups}^{UV})}{P_{abs}^X(1\beta _d)}\frac{E_X}{E_{UV}\eta _{UV}}$$
(41)
where the last term is equal 1 on the basis of Eq. 35.
## Appendix B : Mean number of clouds on the line of sight
We assume that $`N`$ clouds of radius $`r_{cl}`$ are homogeneously distributed in a shell of radius $`r_{UV}`$ whose thickness is small with respect to its radius. The coverage factor of the system of clouds is $`\mathrm{\Omega }/4\pi =C`$. The mean number of clouds on the line of sight, $`N_{ls}`$, is given by:
$$N_{ls}N\frac{X}{4C}$$
(42)
where $`X=(r_{cl}/R)^2`$.
The total number of clouds is related to the coverage factor by:
$$X=1(1C)^{1/N}$$
(43)
This expression can be expanded to the first order:
$$X=\frac{1}{N}Ln\frac{1}{1C},$$
(44)
provided that $`Ln\frac{1}{1C}N`$. This is easily achieved except for extremely small values of $`1C`$ (in our case $`1C0.1`$). One deduces thus that $`X2/N`$, and therefore $`N_{ls}1/C`$. |
warning/0001/hep-th0001056.html | ar5iv | text | # References
1. Introduction
In this paper we will consider the world volume theory of a set of $`N`$ D3 branes, in the limit of large $`N`$, at finite $`\alpha ^{}`$. When embedded in flat 10-d space-time, the theory reduces at low energies to $`𝒩=4`$ supersymmetric Yang-Mills theory. More generally, we can deform the model by turning on gauge invariant couplings $`\varphi ^i`$ and consider the quantum partition function as a function of the $`\varphi ^i`$
$$\mathrm{exp}\left(i\mathrm{\Gamma }(\varphi )\right)=\mathrm{exp}(i\varphi ^iO_i).$$
(1)
Via the famous AdS/CFT correspondence , this partition function has, for $`\alpha ^{}0`$ and for large ’t Hooft coupling $`\lambda =Ng_{ym}^2`$, a dual represention as that of IIB string theory on (a deformation of) $`AdS_5\times S^5`$. In this correspondence, the couplings $`\varphi ^i`$ – which among others include the dilaton, 4-d space-time metric, as well as RR fields – specify the non-dynamical asymptotic values of the corresponding set of closed string fields . When the $`\varphi ^i`$ represent finite expectation values, they generally break the conformal invariance and part or all of the supersymmetry of the low energy theory.
An interesting generalization of this duality arises when all directions perpendicular to the D3-branes are taken to be compact . In this case one can not take the decoupling limit and interactions between open and closed strings will remain relevant. Instead, the target space of the IIB string theory is described by a warped compactification similar to the Randall-Sundrum geometry , and variations of the couplings $`\varphi ^i`$ correspond to normalizable, and thus dynamical, fluctuations of closed string modes. The low energy dynamics of these modes is described by an effective action $`S(\varphi )`$, that includes besides the non-local induced action $`\mathrm{\Gamma }(\varphi )`$ also a local contribution $`S_E(\varphi )`$ arising from the KK reduction of the 10-d effective action of the IIB string theory
$$S(\varphi )=S_E(\varphi )+\mathrm{\Gamma }(\varphi ).$$
(2)
In some properties of this effective action where studied for large $`\lambda `$ with the help of the AdS/CFT correspondence. Using the dual supergravity approximation and elementary results of Hamilton-Jacobi theory, it was found that in this regime $`S(\varphi )`$ satisfies an RG flow equation of the following schematic form
$$\beta ^i(\varphi )\frac{S}{\varphi ^i}\frac{1}{2}G^{ij}\frac{S}{\varphi ^i}\frac{S}{\varphi ^j}=\mathrm{\hspace{0.17em}0}.$$
(3)
Here $`G^{ij}`$ denotes some apropriate metric on the space of couplings, and $`\beta ^i(\varphi )`$ are ‘beta-functions’ that describe the classical flow velocities of the $`\varphi `$-fields as a function of the holographic extra dimension. They are expressed in terms of the local action $`S_E(\varphi )`$ as
$$\beta ^i(\varphi )=G^{ij}\frac{S_E}{\varphi ^i}.$$
(4)
For a more detailed explanation of these relations we refer to .
In eqn (3) was interpreted as an RG invariance of the total gravitational effective action $`S`$, in the sense that any classical extremum of $`S`$ automatically lies on a complete RG trajectory of classical solutions connected by the flow relation
$$\dot{\varphi }^i=\beta ^i(\varphi ).$$
(5)
This property of $`S`$, if true in general, could be of particular importance for the cosmological constant problem. The derivation of (3) as given in , however, is valid only in the limit of large $`N`$ and large ’t Hooft coupling $`\lambda 1`$. It is important therefore to investigate whether these relations can be generalized and extended to other regimes as well. In the following we will concentrate on the dual regime of small $`\lambda ,`$ but still infinite $`N.`$ We will find the positive result that the exact same equations (2), (3) and (4) remain valid also in this regime, except that the quantities $`\mathrm{\Gamma }`$, $`S_E`$, $`G^{ij}`$ and $`\beta ^i`$ are now all defined via their weak coupling descriptions.
The main idea behind our derivation is that the world sheet of a planar multi-loop diagram in open string theory is conformally equivalent to a closed string tree diagram. Indeed, all holes in the open string diagram can be represented in the dual channel by means of external closed string states, equal to the appropriate D-brane boundary state $`|B`$. Via this dual representation all potential UV divergences of the open string diagram become equivalent to potential IR divergences due to on-shell closed string states in the dual channel. It is not too surprising therefore that the RG structure of the large $`N`$ open string theory can be made to look identical to a classical evolution equation of closed string theory.
Several elements in our reasoning have appeared in earlier works. We mention here:
* The Fischler-Susskind mechanism for cancelling string loop divergences ; for a recent discussion of the FS-mechanism in relation with D3-brane physics, see .
* The interpretation given in of the BV symmetry of closed string field theory as an ‘exact’ RG invariance à la Polchinski .
* The non-linear flow relation, proposed by Polyakov in , satisfied by the tree level partition function in a general non-critical string theory.
In a pure field theory context, the structure described below also seems closely related to the old recursive formula for QFT counterterms due to Bogolyubov , that forms the basis for the classic BPHZ renormalization method. We will comment on this correspondence (which may prove useful for taking the $`\alpha ^{}0`$ limit of our result) in section 3. In section 4 we summarize our conclusions.
2. Derivation of the flow equation
For small ’t Hooft coupling $`\lambda `$, the total low energy effective action $`S(\varphi )`$ is obtained by summing over all n-loop planar open string diagrams in the closed string background specified by $`\varphi .`$ Schematically<sup>3</sup><sup>3</sup>3Here we are using the same notation $`O_i`$ for the closed string vertex operators dual to the fields $`\varphi ^i`$, as used in eqn (1) for the gauge theory operators dual to $`\varphi ^i`$. In principle, one can make a one-to-one correspondence between the two, by comparing the 2-d world-sheet action with the D3-brane world volume action in a given closed string background.
$$S(\varphi )=\mathrm{\Gamma }_0(\varphi )+\underset{\mathrm{n}1}{}\lambda ^n\mathrm{\Gamma }_\mathrm{n}(\varphi )$$
(6)
$$\mathrm{\Gamma }_\mathrm{n}(\varphi )=\mathrm{exp}\left(i\varphi ^iO_i\right)_\mathrm{n}$$
(7)
This term $`\mathrm{\Gamma }_n(\varphi )`$ is the $`\mathrm{n}1`$-loop open string contribution, given by the partition function of the world-sheet sigma-model (parametrized by $`\varphi `$) on a sphere with n holes, integrated over all moduli parameterizing the relative sizes and locations of these holes. Both the sigma-model expectation value and the integral over these moduli may produce potentially infinite answers, which both can be regularized by introducing an explicit cut-off scale $`ϵ`$. We will give a concrete prescription for this cut-off momentarily. In any case, in the end all physical answers should, when expressed in terms of renormalized couplings, be independent of this cut-off. In particular, writing
$$S(\varphi (ϵ);ϵ),$$
(8)
where $`\varphi (ϵ)`$ is the renormalized sigma-model background satisfying the RG equation
$$ϵ\frac{\varphi ^i}{ϵ}=\beta ^i(\varphi ),$$
(9)
with $`\beta ^i(\varphi )`$ the sigma model beta-functions, we must require that the total $`ϵ`$ dependence cancels
$$ϵ\frac{dS}{dϵ}=\beta ^i(\varphi )\frac{S}{\varphi ^i}+ϵ\frac{S}{ϵ}=\mathrm{\hspace{0.17em}0}.$$
(10)
This requirement reduces to the usual condition of conformal invariance in the limit $`\lambda 0`$, when the holes of the open string loops are absent. The condition (10) with $`\lambda >0`$ relates divergences arising from deviation from world-sheet conformal invariance to divergences coming from the shrinking of open string loops; the cancellation of these two different types of divergences is known as the Fischler-Susskind mechanism .
To implement this cancellation mechanism, we need a sufficiently precise definition of the cut-off $`ϵ`$, both for the divergences of the sigma-model as well as for regulating the moduli integral. For this we will make use of some technology from closed string field theory . To any Riemann surface (with holes) with given conformal structure, we can assign a unique minimal area metric $`g_{\alpha \beta }.`$ For a given point on the moduli space of the open string loop diagram, we can thus measure the minimal geodesic length $`\mathrm{}(C)`$, as defined using this minimal area metric, of all non-contractible contours $`C`$ surrounding a non-zero number of holes. (See fig 1a). The UV divergences of the loop diagram arise when one or more of these geodesic lengths $`\mathrm{}(C)`$ tends to zero. We will therefore introduce a UV regulator $`ϵ`$ by requiring that the moduli integral is restricted to those conformal structures for which
$$\mathrm{}(C)ϵ$$
(11)
for all non-contractible contours $`C.`$ Hence the boundary of the regulated moduli space are degenerate surfaces for which the above bound (11) is saturated for one or more contours $`C`$.
Since in the end we need to compare this type of degeneration of the open string loop diagram with the sigma-model divergences, it seems most practical to regulate the sigma-model expectation values in an analogous fashion. To this end, we explicitly expand the exponential in eqn (7)
$$\mathrm{\Gamma }_\mathrm{n}(\varphi ;ϵ)=\underset{k0}{}\frac{1}{k!}\genfrac{}{}{0pt}{}{\underset{}{\mathrm{\Phi }\mathrm{}\mathrm{\Phi }}}{k\times }_\mathrm{n}=\underset{k0}{}\frac{1}{k!}(\mathrm{\Phi })^k_\mathrm{n}$$
(12)
with
$$\mathrm{\Phi }=\underset{i}{}\varphi ^iO_i.$$
(13)
The $`k`$-th order term on the right-hand side is a correlator, defined in the $`\varphi =0`$ sigma-model, of $`k`$ operators $`\mathrm{\Phi }`$ on an $`\mathrm{n}1`$ loop open string diagram. The resulting amplitude is therefore an integral over the moduli space of a sphere with $`n`$ holes and $`k`$ punctures. (See fig 2a.) We can now apply the same construction as above, and use the unique minimal area metric on this punctured surface to assign a given minimal geodesic length to all closed contours surrounding a non-zero number of holes and/or punctures, and require that all such lengths must be larger than the cut-off $`ϵ.`$ In this way we have indeed introduced one uniform cut-off procedure for both types of divergences.<sup>4</sup><sup>4</sup>4Given the limited available tools for dealing with sigma-models with RR backgrounds, the procedure outlined here seems at present the only precise method for extracting the cut-off dependence of the sigma-model expectation values. We should further mention that, in case the string theory under consideration is an orientifold compactification, one also needs to include world-sheets with an arbitrary number of cross-caps. These can be treated in a similar fashion as the open string holes and $`\mathrm{\Phi }`$ operator insertions.
The above equation (12) in fact needs some extra specification for the case $`n=0`$. On the sphere without holes one needs, due to the invariance under the global conformal group $`SL(2,C)`$, at least three operator insertions $`\mathrm{\Phi }`$ to obtain a well-defined expectation value. To write the $`n=0`$ contribution to the effective action, we must thus include a separate kinetic term via
$$\mathrm{\Gamma }_0(\varphi ;ϵ)=\frac{1}{2}\mathrm{\Phi }|Q|\mathrm{\Phi }+\underset{k3}{}\frac{1}{k!}(\mathrm{\Phi })^k.$$
(14)
Here $`Q`$ denotes the nil-potent world-sheet BRST charge of the $`\varphi =0`$ sigma-model and $`|\mathrm{\Phi }=_i\varphi ^i|O_i`$ is the state corresponding to the sigma-model background $`\mathrm{\Phi }`$. The above expression for $`\mathrm{\Gamma }_0(\varphi )`$ is of the same form as the standard classical action of closed string field theory (for a detailed discussion of its definition and properties, see ). Note, however, that in the present context the classical equations of motion of $`\varphi `$ must be derived from the total action $`S(\varphi )`$ given in (6), and not from just $`\mathrm{\Gamma }_0(\varphi )`$. Indeed, since $`\mathrm{\Gamma }_0(\varphi )`$ has no contribution from surfaces with holes, it does not have any direct knowledge of the presence of the D3 branes.
It is now straightforward to determine the $`ϵ`$ dependence of the total effective action. Each term $`(\mathrm{\Phi })^k_\mathrm{n}`$ is given by an integral over the corresponding moduli space, whose only dependence on $`ϵ`$ is via the restriction (11) on the geodesic lengths. Hence if we differentiate $`(\mathrm{\Phi })^k_\mathrm{n}`$ with respect to $`ϵ`$, the result is an integral over the boundary sub-space for which (at least) one of the contours $`C`$ has reached its minimal length $`\mathrm{}(C)=ϵ`$. Now it’s a well known fact that such a pinched surface is conformally equivalent to a surface for which the closed string propagator in the dual channel has acquired a large length proportional to $`1/ϵ.`$ (See figs 1b and 2b.) The partition function for this degenerate surface factorizes into a sum of products of two one-point functions defined on each half of the surface on each side of this long propagator. Concretely, writing the evolution operator along this long tube as
$$ϵ^{L_0+\overline{L}_0}=\frac{1}{2}\underset{i}{}|O_iG^{ij}O_j|,$$
(15)
(here we assume that the set of states $`|O_i`$ forms a complete basis of closed string states) we can express the explicit $`ϵ`$-dependence of the $`\mathrm{n}`$ loop partition function $`\mathrm{\Gamma }_\mathrm{n}`$ as follows
$$ϵ\frac{\mathrm{\Gamma }_\mathrm{n}}{ϵ}=\frac{1}{2}\underset{0\mathrm{m}\mathrm{n}}{}G^{ij}\frac{\mathrm{\Gamma }_\mathrm{m}}{\varphi ^i}\frac{\mathrm{\Gamma }_{\mathrm{n}\mathrm{m}}}{\varphi ^j}$$
(16)
The summation here runs over all possible ways of dividing the surface with $`\mathrm{n}`$ holes into two parts, as indicated in figs 1b and 2b.
The $`\mathrm{m}=0`$ and $`\mathrm{m}=\mathrm{n}`$ terms in eqn (16), the ones containing a factor $`\mathrm{\Gamma }_0/\varphi ^i`$, describe the ‘splitting off’ of a sphere without holes as indicated fig 2b, and represent the scale dependence due to the sigma model divergences.<sup>5</sup><sup>5</sup>5Note that (via the presence of the kinetic term in (14)) $`\mathrm{\Gamma }_0/\varphi ^i`$ starts out for small $`\varphi `$ with a linear term proportional to $`\mathrm{\Delta }_j^i\varphi ^j`$, with $`\mathrm{\Delta }_j^i`$ the matrix of scaling dimensions of the $`O_i`$ (defined via $`L_0|O_i=\mathrm{\Delta }_i^j|O_j`$). The corresponding explicit cut-off dependence of the correlators $`(\mathrm{\Phi })^k_\mathrm{n}`$ arises from the local regularization of the individual $`\mathrm{\Phi }`$ operator insertions. Our method of renormalization of the sigma model is to absorb this particular type of divergence by means of the transition to renormalized couplings $`\varphi ^i(ϵ)`$. In other words, the $`ϵ`$ dependence of $`\varphi ^i(ϵ)`$ should be such that the terms proportional to $`\mathrm{\Gamma }_0/\varphi ^i`$ cancel in the total $`ϵ`$-variation of $`\mathrm{\Gamma }_n`$. From this we deduce that the renormalized couplings must satisfy (9) with
$$\beta ^i(\varphi )=G^{ij}\frac{\mathrm{\Gamma }_0}{\varphi ^j}.$$
(17)
Note that a relation of this form is indeed expected from the identification (14) of $`\mathrm{\Gamma }_0(\varphi )`$ with the classical closed string field theory action.
\[ This renormalization procedure (17) is further motivated by the requirement that the eventual renormalized form of the RG-equation should be independent of the cut-off $`ϵ`$. In particular the factorization metric $`G^{ij}`$ introduced in (15) should be cut-off independent. It seems reasonable to assume that the only divergences, that induce a renormalization of $`G^{ij}`$, are those of the sigma-model. In our regularization method, these arise from ‘contact term’ contributions produced when one or more of the $`\mathrm{\Phi }`$-operator insertions get ‘trapped’ somewhere on the long tube in fig 2b. The effect of these terms is to replace the $`L_0`$-operators in (15) by those in the $`\mathrm{\Phi }`$-background. The resulting condition of $`ϵ`$-independence of $`G^{ij}`$ takes the form
$$ϵ\frac{G^{ij}}{ϵ}+\beta ^k\frac{G^{ij}}{\varphi ^k}\mathrm{\Delta }_k^iG^{kj}\mathrm{\Delta }_k^jG^{ik}=0.$$
(18)
with $`\beta ^i`$ as in (17) and where $`\mathrm{\Delta }_i^k=_i\beta ^k`$ denotes the matrix of conformal dimensions of the $`O_i`$ in the $`\mathrm{\Phi }`$-background. \]
Via the definition of the total action $`S`$ as the sum of all $`\mathrm{\Gamma }_n`$’s, we can summarize the infinite set of relations (16) as a single non-linear equation for $`S`$.
$$ϵ\frac{S}{ϵ}=\frac{1}{2}G^{ij}\frac{S}{\varphi ^i}\frac{S}{\varphi ^j},$$
(19)
and we thus indeed obtain the announced form for the total FS scale invariance condition (10)
$$\beta ^i(\varphi )\frac{S}{\varphi ^i}\frac{1}{2}G^{ij}\frac{S}{\varphi ^i}\frac{S}{\varphi ^j}=\mathrm{\hspace{0.17em}0}.$$
(20)
This equation, as well as our derivation, are strongly reminiscent of Polchinski’s version of the Wilsonian renormalization group . Indeed our UV-regulator (11) looks just like an IR cut-off on the (proper length of the) closed string propagator in the dual channel. It is not surprising therefore that the explicit $`ϵ`$-dependence of $`S`$ looks just like (the classical limit of) the RG-equation derived in .
Finally, we remark that the total equation of motion for $`S`$ implies that on-shell
$$G^{ij}\frac{S}{\varphi ^j}=\beta ^i(\varphi )+\underset{\mathrm{n}1}{}\lambda ^\mathrm{n}G^{ij}\frac{\mathrm{\Gamma }_\mathrm{n}}{\varphi ^j}=0,$$
(21)
which tells us that the deviation from scale-invariance coming from the sigma-model beta-functions cancels against that coming from the source terms due to the presence of the D3-branes, and/or from the open string loop contributions.<sup>6</sup><sup>6</sup>6The term with n=1 is a pure D3-brane source term; the other terms with n$`>`$1 can either be viewed as higher order closed string source terms or as quantum corrections due to open string loops.
3. Equation of Motion and Bogolyubov’s Recursion Formula
Although the above line of argument did not make any use of the AdS/CFT correspondence, it seems still quite essential that the planar diagrams were in fact string world sheets, so that one can easily visualize the transition from the open to the closed string channel. In principle, however, it should be straightforward (though presumably quite tedious) to explicitly take the $`\alpha ^{}0`$ limit and translate all steps into purely field theoretic language. Here we limit ourselves to just a couple of remarks.
As indicated in fig 3, our regulatization method for the open string diagrams carries over quite directly to ordinary planar diagrams. Feynman diagrams of the low energy large $`N`$ gauge theory can also be written as integral over a “moduli space” of Schwinger parameters, parametrizing the proper length $`t_i`$ of the propagators. The UV regulator (11) thus amounts to the restriction that for any closed path $`C`$ on the graph (as in fig 3)
$$\underset{iC}{}t_i>ϵ.$$
(22)
where the sum is over all propagators that make up the contour $`C`$. This restriction indeed renders the integral UV finite.
We can now use a similar reasoning as above to try and extract the $`ϵ`$ dependence, by explicitly differentiating the total integral over all Schwinger parameters with respect to the UV cut-off (22). The analog of the formula (15) should now be extracted from analysing the UV limit of the one-loop gauge theory amplitude in a background large $`N`$ gauge-field $`A`$, with couplings $`\varphi ^i`$ turned on; equation (15) then corresponds to the fact that, to leading order in $`1/N`$, this amplitude factorizes into a sum over gauge invariant single trace-operators $`O_i`$.
Useful insight into how one should interpret the sigma model data contained in $`\mathrm{\Phi }`$ is obtained by considering the equation of motion of the total effective action (16). It is possible to write it in the form of a recursion relation, by expanding the closed string background $`\mathrm{\Phi }`$ in powers of the string coupling $`\lambda `$
$$\mathrm{\Phi }=\underset{\mathrm{n}1}{}\lambda ^\mathrm{n}\mathrm{\Phi }_\mathrm{n}$$
(23)
where $`\mathrm{\Phi }_\mathrm{n}`$ is assumed to be independent of $`\lambda `$. The equation of motion of $`\mathrm{\Phi }_\mathrm{n}`$
$$\frac{\delta S}{\delta \mathrm{\Phi }_\mathrm{n}}=\mathrm{\hspace{0.17em}0}$$
(24)
now takes the following form
$$Q|\mathrm{\Phi }_\mathrm{n}=[\mathrm{\hspace{0.33em}1}]_\mathrm{n}+\underset{1mn1}{}\underset{_\mathrm{j}\mathrm{j}k_\mathrm{j}=\mathrm{m}}{}\frac{1}{k_1!\mathrm{}k_\mathrm{m}!}[(\mathrm{\Phi }_1)^{k_1}\mathrm{}(\mathrm{\Phi }_\mathrm{m})^{k_\mathrm{m}}]_{\mathrm{n}\mathrm{m}}.$$
(25)
Here $`[\mathrm{}]_\mathrm{n}`$ denotes the state associated to a surface as indicated in fig 4: the sphere with $`n`$ holes at the end of a tube with length $`1/ϵ`$, and with operator insertions specified by the $`(\mathrm{})`$. The above formula can be used to recursively construct $`\mathrm{\Phi }_\mathrm{n}`$ from the previous $`\mathrm{\Phi }_\mathrm{m}`$’s with $`\mathrm{m}<\mathrm{n}`$.
The field theoretic meaning of the above recursive form of the equation of motion is as follows: the term $`\mathrm{\Phi }_\mathrm{n}`$ is associated with the n-th order counter term designed to cancel the divergent contribution $`[\mathrm{\hspace{0.33em}1}]_\mathrm{n}`$ of the n-loop vacuum graph, that is left after subtracting all sub-divergences of lower order. There is indeed a strong resemblance between eqn (25) and the famous recursion relation due to Bogolyubov for constructing the successive counterterms in the BPHZ renormalization method of QFT. We suspect that, by carefully taking the $`\alpha ^{}0`$ limit of (25), this match can be made even more precise.<sup>7</sup><sup>7</sup>7It is also not entirely coincidental that the general solution to Bogolyubov’s recursion relation, due to Zimmerman , is known as the forest formula. In the present context, a forest (a nested or disjoint set of sub-diagrams associated with nested or disjoint sub-divergences) in essence corresponds to a set of sub-trees of a closed string tree diagram.
4. Conclusion: World-sheet vs Space-Time RG
Our main result is equation (20) and the fact that it reproduces (at least in form and with the identification $`\mathrm{\Gamma }_0(\varphi )=S_E(\varphi )`$) the flow equation (3) derived from classical 5-d supergravity. Both derivations have their limited range of validity: the AdS/CFT duality used in is accurate only the regime of $`\lambda >>1`$, whereas the perturbative reasoning of this paper requires $`\lambda <<1`$. The correspondence between the two, however, provides strong evidence that the same structure persists for all values of the ’t Hooft coupling $`\lambda `$.
Equation (20) has an interpretation as both a world-sheet and space-time RG relation.<sup>8</sup><sup>8</sup>8Suggestions as well as concrete proposals towards a correspondence between world-sheet and space-time RG were made by many authors, including in , , , and . The first interpretation is manifest from our derivation, since the $`\beta ^i(\varphi )`$ represent the world-sheet scale dependence of the sigma-model couplings $`\varphi ^i`$. The second interpretation arises via the identification of space-time RG transformations with Weyl rescalings of the 4-d target-space metric $`g^{\mu \nu }`$. This motivates the following translation code between the world-sheet and space-time beta-functions (cf. )
$`\beta _{_{\mathrm{WS}}}^i`$ $`=`$ $`\gamma \beta _{_{\mathrm{ST}}}^i`$ (26)
$`\beta _{_{\mathrm{WS}}}^{\mu \nu }`$ $`=`$ $`2\gamma g^{\mu \nu }+\gamma \beta _{_{\mathrm{ST}}}^{\mu \nu }`$ (27)
where $`\beta _{_{WS}}^{\mu \nu }`$ denotes the world-sheet beta-function of the target-space metric $`g^{\mu \nu }`$. Via the above translation, the invariance condition (20) takes the form of a quite conventional RG-flow equation in space-time.
Note that in this correspondence the world-sheet and space-time scales are proportional to each other, so UV and UV are mapped to each other as well as IR to IR. We further see from (27) that the prefactor $`\gamma `$, which determines the relative normalization of the world-sheet and space-time beta-functions, is (via eqn (17)) proportional to the cosmological term in $`\mathrm{\Gamma }_0(\varphi )`$. This term, generated by the RR background as well as by other expectation values, is typically of order the fundamental string scale and can be even larger. Although, as seen from , the interpretation of the 4-d scale as a holographic extra dimension is expected to break down in this regime, we have shown here that the RG-invariance (20) of $`S`$ is nonetheless preserved.
Acknowledgements
This work is supported by NSF-grant 98-02484. We would like to thank M. Berkooz, R. Dijkgraaf, G. Lifschytz, V. Periwal, A. Polyakov and E. Verlinde for useful discussions. J.K. is supported by the Natural Sciences and Engineering Research Council of Canada. |
warning/0001/math0001029.html | ar5iv | text | # Some Generalized Kac-Moody Algebras With Known Root Multiplicities
(Date: 1 July, 1997, revised 20 July 1999)
## Abstract.
Starting from Borcherds’ fake monster Lie algebra we construct a sequence of six generalized Kac-Moody algebras whose denominator formulas, root systems and all root multiplicities can be described explicitly. The root systems decompose space into convex holes, of finite and affine type, similar to the situation in the case of the Leech lattice. As a corollary, we obtain strong upper bounds for the root multiplicities of a number of hyperbolic Lie algebras, including $`AE_3`$.
The author was supported by the Science and Engineering Research Council (UK), and Peterhouse, Cambridge. |
warning/0001/hep-ex0001020.html | ar5iv | text | # 1 Introduction
## 1 Introduction
In calorimeter simulation different tasks can be distinguished: calorimeter studies, physics analysis, and feasibility studies. A detailed simulation, where all secondary particles are tracked individually down to some minimum energy and where the response is predicted from “first principles”, is required for accurate calorimeter studies. For physics analysis and feasibility studies large number of Monte Carlo events may have to be produced. Using individual particle tracking, the computing time needed for such kind of simulations increases approximately linear with the energy absorbed in the detector and can easily become prohibitive. Using parameterizations for electromagnetic (sub)showers can speed up the simulations considerably, without sacrificing precision. The high particle multiplicity in electromagnetic showers as well as their compactness and the good understanding of the underlying physics makes their parameterization advantageous.
Using an Ansatz by Longo and Sestili , a simple algorithm for the description of longitudinal shower profiles has been used successfully for the simulation of the UA1 calorimeter . Later, this Ansatz has been extended to the simulation of individual showers, taking their shower-to-shower fluctuation and correlations consistently into account . For the parameterized simulation of radial energy profiles no conclusive procedure has been established until now.
In homogeneous media, a scaling of the longitudinal and radial profiles in radiation lengths and Molière radii respectively does not lead to a material independent description of electromagnetic shower development. In sampling calorimeters, the shower shapes depend in addition on the sampling structure. We have extended the above Ansatz for parameterized simulation of longitudinal profiles by taking the material and geometry dependence of the parameters into account and developed a new algorithm to simulate radial energy distributions . Correlations between the longitudinal and radial shower development have been included.
## 2 Procedure
To arrive at a general description of electromagnetic shower development, we performed detailed Monte Carlo simulations, on a grid of 1.0 $`X_0`$ in depth and 0.2 Molière radii laterally, for various homogeneous media and sampling calorimeters, using the GEANT package . The materials used were Cu, Fe, W, Pb, U, and scintillator and liquid argon. In a first step only average shower profiles in homogeneous media were analyzed, from which scaling laws for the material and energy dependence of the parameters have been extracted. Starting from the relations which describe the average behavior of the parameterized quantities, we developed parameterizations for individual electromagnetic showers in homogeneous calorimeters, taking fluctuations and correlations into account.
The parameterizations in homogeneous media are a first approximation for electromagnetic shower development in sampling calorimeters, which are viewed as consisting of one single effective medium. The inhomogeneous material distribution in sampling calorimeters influences however the exact behavior of the shower shapes. This is explained mainly by the transition effect which depends on the shower depth . These effects have been taken into account by adding geometry dependent terms to the parameterizations for homogeneous media, which can be easily calculated from the sampling geometry.
## 3 Parameterization Ansatz
The spatial energy distribution of electromagnetic showers is given by three probability density functions (pdf),
$$dE(\stackrel{}{r})=Ef(t)dtf(r)drf(\varphi )d\varphi ,$$
(1)
describing the longitudinal, radial, and azimuthal energy distributions. Here $`t`$ denotes the longitudinal shower depth in units of radiation length, $`r`$ measures the radial distance from the shower axis in Molière units, and $`\varphi `$ is the azimuthal angle. The start of the shower is defined by the space point, where the first electron or positron bremsstrahlung process occurs. A gamma distribution is used for the parameterization of the longitudinal shower profile, f(t). The radial distribution, f(r), is described by a two-component Ansatz. In $`\varphi `$, it is assumed that the energy is distributed uniformly: $`f(\varphi )=1/2\pi `$.
### 3.1 Longitudinal shower profiles – homogeneous media
It is well known that average longitudinal shower profiles can be described by a gamma distribution :
$$\frac{1}{E}\frac{dE(t)}{dt}=f(t)=\frac{(\beta t)^{\alpha 1}\beta \mathrm{exp}(\beta t)}{\mathrm{\Gamma }(\alpha )}.$$
(2)
The center of gravity, $`t`$, and the depth of the maximum, $`T`$, can be calculated from the shape parameter $`\alpha `$ and the scaling parameter $`\beta `$ according to
$`t`$ $`=`$ $`{\displaystyle \frac{\alpha }{\beta }}`$ (3)
$`T`$ $`=`$ $`{\displaystyle \frac{\alpha 1}{\beta }}.`$ (4)
Longitudinal electromagnetic shower development in homogeneous media had been studied analytically by Rossi . An important result of the calculations using “Rossi Approximation B“ is that longitudinal shower moments are equal in different materials, provided one measures all lengths in units of radiation length $`(X_0)`$ and energies in units of the critical energy ($`E_c`$). Numerically, $`E_c`$ can be calculated according to
$$E_c=\mathrm{\hspace{0.17em}2.66}\left(X_0\frac{Z}{A}\right)^{1.1}.$$
(5)
For the depth of the shower maximum
$`T\mathrm{ln}y=\mathrm{ln}{\displaystyle \frac{E}{E_c}}`$ (6)
is predicted .
It is therefore desirable to use $`T`$ in the parameterization. This is demonstrated in Fig.1, where the average depth of the shower maximum for various homogeneous media<sup>2</sup><sup>2</sup>2 The index “hom“ in the following formulae indicates the validity for homogeneous media. For sampling calorimeters the index “sam“ will be used., $`T_{hom}`$, is plotted versus $`y`$, in the energy range from 1 to 100 GeV. As a second variable $`\alpha `$ is used. In this case the parameterization depends on the charge number $`Z`$ of the medium, as can be seen in Fig.2. The lines in both figures correspond to fits to GEANT simulations according to
$`T_{hom}`$ $`=`$ $`\mathrm{ln}y+t_1`$ (7)
$`\alpha _{hom}`$ $`=`$ $`a_1+(a_2+a_3/Z)\mathrm{ln}y.`$ (8)
The values of the coefficients are given in Appendix, where all formulae and numbers, which will be given in the following, are summarized.
Assuming that also individual profiles can be approximated by a gamma distribution, the fluctuations and correlations can be taken into account consistently (for details refer to ). For each single GEANT-simulated shower, $`T`$ and $`\alpha `$ are determined by fitting a gamma distribution. The logarithms of $`T`$ and $`\alpha `$ are used for the parameterization since they are found to be approximately normal distributed. For the parameterization of $`\mathrm{ln}T_{hom}`$ and $`\mathrm{ln}\alpha _{hom}`$ the logarithms of equations 7 and 8 are used. The $`y`$-dependence of the fluctuations can be described by
$$\sigma =(s_1+s_2\mathrm{ln}y)^1.$$
(9)
The correlation between $`\mathrm{ln}T_{hom}`$ and $`\mathrm{ln}\alpha _{hom}`$ is given by
$$\rho (\mathrm{ln}T_{hom},\mathrm{ln}\alpha _{hom})\rho =r_1+r_2\mathrm{ln}y.$$
(10)
The dependence of these quantities on $`y`$ is shown in Fig.3 for various materials together with the parameterizations (see Appendix A.1.2).
From these formulae, correlated and varying parameters $`\alpha _i`$ and $`\beta _i`$ are generated according to
$$\left(\begin{array}{c}\mathrm{ln}T_i\\ \mathrm{ln}\alpha _i\end{array}\right)=\left(\begin{array}{c}\mathrm{ln}T\\ \mathrm{ln}\alpha \end{array}\right)+C\left(\begin{array}{c}z_1\\ z_2\end{array}\right)$$
(11)
with
$$C=\left(\begin{array}{cc}\sigma (\mathrm{ln}T)& 0\\ 0& \sigma (\mathrm{ln}\alpha )\end{array}\right)\left(\begin{array}{cc}\sqrt{\frac{1+\rho }{2}}& \sqrt{\frac{1\rho }{2}}\\ \sqrt{\frac{1+\rho }{2}}& \sqrt{\frac{1\rho }{2}}\end{array}\right)$$
and $`\beta _i=(\alpha _i1)/T_i`$ and $`z_1`$ and $`z_2`$ are standard normal distributed random numbers. The longitudinal energy distribution is evaluated<sup>3</sup><sup>3</sup>3 The GAMDIS function of the CERN computer library is used. by integration in steps of $`\mathrm{\Delta }t=t_jt_{j1}=1X_0`$,
$$dE(t)=E_{t_{j1}}^{t_j}\frac{(\beta _it)^{\alpha _i1}\beta _i\mathrm{exp}(\beta _it)}{\mathrm{\Gamma }(\alpha _i)}𝑑t.$$
It is worthwhile to mention that only one of the five quantities needed, $`\mathrm{ln}\alpha _{hom}`$, depends explicitly on the material, while for the other four this dependence is absorbed by using $`y`$ instead of $`E`$.
In Fig.4 longitudinal profiles of GEANT and parameterized simulations for a lead glass calorimeter (SF5) are compared. Shown are the mean profiles and the mean + 1 RMS in each $`X_0`$ interval. While the means are in perfect agreement, the fluctuations are underestimated by the parameterized simulations at low energies, indicating that the description of individual profiles by gamma distributions becomes a worse approximation with decreasing shower energy. Comparisons for other materials (Fe, Cu, W, Pb, U) are of comparable quality as those in Fig.4 . In the next sections we will show, how the sampling fluctuations in sampling calorimeters can be used to improve the shape fluctuations at low energies.
### 3.2 Sampling fluctuations
In fast simulations, sampling calorimeters consisting of a complicated but repetitive sampling structure are usually described by one single effective medium (the formulae to compute effective material parameters are summarized in Appendix A.2.1). The sampling fluctuations, the scaling of the deposited energy to the visible energy using an appropriate sampling fraction, and the effects of the sampling structure have to be considered in parameterized simulations explicitly.
The simulation of sampling fluctuations are done conveniently with a gamma distribution:
$$G(a,b)=\frac{x^{a1}be^{bx}}{\mathrm{\Gamma }(a)}$$
(12)
with
$$x=\frac{a}{b},\sigma ^2(x)=\frac{a}{b^2}.$$
(13)
The energy in each longitudinal integration step, $`dE(t)`$, is fluctuated<sup>4</sup><sup>4</sup>4 The RANGAM function of the CERN computer library is used. according to equation 12 choosing
$$a=\frac{dE(t)}{c^2}\text{ and }b=\frac{1}{c^2}.$$
(14)
It is then easy to show that the central limit theorem will ensure the total energy to be normal distributed obeying the usual formula for the sampling fluctuations:
$$\frac{\sigma }{E}=\frac{c}{\sqrt{E}}.$$
(15)
Using this procedure the occurrence of negative energies is automatically avoided. Additional fluctuations of the longitudinal shape are introduced, leading to a better agreement in the shape fluctuations. This method is also used to fluctuate energy depositions of real particles (electrons, hadrons), when they are tracked individually through an effective homogeneous volume.
### 3.3 Longitudinal shower profiles – sampling calorimeters
The inhomogeneous material distribution in sampling calorimeters influences the exact behavior of the shower shapes. In the first stages of electromagnetic shower development the signal is dominated by electrons and positrons. Behind the shower maximum low energetic photons become more and more important. The transition effect, being explained mainly by the absorption properties of low energetic photons, must in turn depend on the shower depth. Consequently, the signal ratio of electrons to minimum ionizing particles, $`e/mip`$, decreases continuously as the shower propagates longitudinally. Thus the signal maximum in a sampling calorimeter occurs at an earlier depth than expected for a homogeneous calorimeter with the same effective material properties. This can be seen from Fig.5 (left upper corner), where $`\mathrm{ln}T`$ for homogeneous media is compared to the values in five different sampling calorimeters. In addition, the amount of the shift of $`\mathrm{ln}T`$ depends on the exact geometrical arrangement.
The parameterization of the longitudinal shape as given in section 3.1 for homogeneous media can therefore not be used for sampling calorimeters directly. Instead it may be understood as a first approximation to which geometry dependent corrections have to be added. We use the sampling frequency
$$F_S=\frac{X_{0,eff}}{d_a+d_p}$$
(16)
and the value of $`e/mip`$ (averaged over the shower depth) to account for the shower depth dependence of the transition effect. $`d_a`$ and $`d_p`$ denote the thickness of the active and passive layers, respectively. If $`e/mip`$ is not known, a sufficiently good approximation for many calorimeters with charge numbers $`Z_p`$ and $`Z_a`$ is given by
$$\widehat{e}=\frac{1}{1+0.007(Z_pZ_a)}\frac{e}{mip}.$$
(17)
Averaged over the whole shower, $`e/mip`$ remains energy independent for $`E\stackrel{>}{}1`$ GeV.
The average longitudinal profiles can now be parameterized according to
$`T_{sam}`$ $`=`$ $`T_{hom}+t_1F_S^1+t_2(1\widehat{e})`$ (18)
$`\alpha _{sam}`$ $`=`$ $`\alpha _{hom}+a_1F_S^1,`$ (19)
and the quantities used for the simulation of individual showers are given by
$`\mathrm{ln}T_{sam}`$ $`=`$ $`\mathrm{ln}\left(\mathrm{exp}(\mathrm{ln}T_{hom})+t_1F_S^1+t_2(1\widehat{e})\right)`$ (20)
$`\mathrm{ln}\alpha _{sam}`$ $`=`$ $`\mathrm{ln}\left(\mathrm{exp}(\mathrm{ln}\alpha _{hom})+a_1F_S^1\right).`$ (21)
The fluctuations, $`\sigma (\mathrm{ln}T_{sam})`$, $`\sigma (\mathrm{ln}\alpha _{sam})`$ and the correlation, $`\rho (\mathrm{ln}T_{sam},\mathrm{ln}\alpha _{sam})`$, are described with the help of the same formulae as in the case of homogeneous media (see Appendix A.2.2 and A.2.3).
Fig.5 summarizes the parameterization for sampling calorimeters. The expectation value of $`\mathrm{ln}T`$ no longer scales with $`y`$. The expectation value of $`\mathrm{ln}\alpha `$ depends on the material and the sampling geometry. The fluctuations and correlations of the parameters can still be approximated without any explicit material or geometry dependence.
In Figs.6 to 8 GEANT and parameterized simulations of the lead liquid argon calorimeter (IFE) of the H1 experiment are compared. The GEANT simulations were performed with low energy cuts ($`e`$-cut$`=200`$ keV, $`\gamma `$-cut$`=10`$ keV) and a detailed geometry description, including for example copper pads and G10 layers. These simulations were not used to tune the parameterizations. Both, average longitudinal profiles and their fluctuations (including sampling fluctuations) are in very good agreement (see Fig.6). The energy containment (see Fig.7) and the energy resolution (see Fig.8) as a function of the longitudinal calorimeter length are also well predicted. Comparisons with detailed simulations of other calorimeters (Fe-LAr, Cu-Sc, W-LAr, Pb-LAr, U-Sc) show a comparably good performance .
### 3.4 Radial shower profiles – homogeneous media
Average radial energy profiles,
$$f(r)=\frac{1}{dE(t)}\frac{dE(t,r)}{dr},$$
(22)
at different shower depths in pure uranium are presented in Fig.9. These profiles show a distinct maximum in the core of the shower which vanishes with increasing shower depth. In the tail ($`r\stackrel{>}{}1R_M`$) the distribution looks nearly flat at the beginning ($`12X_0`$), becomes steeper at moderate depths ($`56X_0`$, $`1314X_0`$), and becomes flat again ($`2223X_0`$). A variety of different functions can be found in the literature to describe radial profiles . We use the following two component Ansatz, an extension of :
$`f(r)`$ $`=`$ $`pf_C(r)+(1p)f_T(r)`$
$`=`$ $`p{\displaystyle \frac{2rR_C^2}{(r^2+R_C^2)^2}}+(1p){\displaystyle \frac{2rR_T^2}{(r^2+R_T^2)^2}}`$
with
$$0p1.$$
Here $`R_C`$ ($`R_T`$) is the median of the core (tail) component and $`p`$ is a probability giving the relative weight of the core component. For the shower depth $`12X_0`$ the distributions $`f(r)`$, $`pf_C(r)`$, and $`(1p)f_T(r)`$ are also indicated in Fig.9.
The evolution of $`R_C`$, $`R_T`$, and $`p`$ with increasing shower depth is shown in Fig.10 for 100 GeV showers in iron and uranium. We use the variable $`\tau =t/T`$, which measures the shower depth in units of the depth of the shower maximum, to generalize the radial profiles. This makes the parameterization more convenient and separates the energy and material dependence of various parameters. The median of the core distribution, $`R_C`$, increases linearly with $`\tau `$. The weight of the core, $`p`$, is maximal around the shower maximum, and the width of the tail, $`R_T`$, is minimal at $`\tau 1`$. This behavior can be traced back to the radial profiles shown in Fig.9.
The following formulae are used to parameterize the radial energy density distribution for a given energy and material:
$`R_{C,hom}(\tau )`$ $`=`$ $`z_1+z_2\tau `$ (24)
$`R_{T,hom}(\tau )`$ $`=`$ $`k_1\{\mathrm{exp}(k_3(\tau k_2))+\mathrm{exp}(k_4(\tau k_2))\}`$ (25)
$`p_{hom}(\tau )`$ $`=`$ $`p_1\mathrm{exp}\left\{{\displaystyle \frac{p_2\tau }{p_3}}\mathrm{exp}\left({\displaystyle \frac{p_2\tau }{p_3}}\right)\right\}`$ (26)
The parameters $`z_1\mathrm{}p_3`$ are either constant or simple functions of $`\mathrm{ln}E`$ or $`Z`$ (see Appendix A.1.3 for details). The complicated evolution of $`R_T`$ and $`p`$ with the shower depth and the dependence on the material can be explained mainly with the propagation of low energetic photons . The offset in $`R_T`$ between iron and uranium (Fig.10) for example, indicating a wider distribution in iron, reflects the difference in the mean free path, which for 1 MeV photons is approximately twice as long in iron as in uranium, if lengths are measured in Molière units.
We found a good agreement of mean radial profiles between parameterized and detailed simulations in Fe, Cu, W, Pb, and U absorbers for energies between 0.4 and 400 GeV. This is demonstrated in Fig.12, where radial profiles in various shower depths are compared for 40 GeV showers in lead and 100 GeV showers in uranium.
The introduction of radial shape fluctuations has to be considered with some care. Even if no fluctuations of $`f(r)`$ are simulated explicitly, the radial energy profile at a given shower depth will fluctuate, because the shower maximum $`T`$ and thus $`\tau `$ varies from shower to shower. Another source of radial fluctuations arises from the method, which we have adopted for the simulation of radial distributions. The energy content of a longitudinal interval of length 1 $`X_0`$, $`dE(t)`$, is calculated from the actual longitudinal energy density distribution as described in section 3.1. This energy is divided into $`N_S(t)`$ discrete spots of energy $`E_S=dE(t)/N_S(t)`$, which are distributed radially according to $`f(r)`$ using a Monte Carlo method. This can be done easily since the pdfs, $`f_C(r)`$ and $`f_T(r)`$, can be integrated and inverted:
$`F(r)`$ $`=`$ $`{\displaystyle _0^r}{\displaystyle \frac{2r^{}R^2}{\left(r^2+R^2\right)^2}}𝑑r^{}={\displaystyle \frac{r^2}{r^2+R^2}}`$ (27)
$`F^1(u)`$ $`=`$ $`R\sqrt{{\displaystyle \frac{u}{1u}}}.`$ (28)
Random radii are generated according to $`f(r)`$ in the following way, using two normal distributed random numbers $`v_i`$ and $`w_i`$:
$$r_i=\{\begin{array}{cc}R_C\sqrt{\frac{v_i}{1v_i}},\hfill & \text{if }p<w_i\hfill \\ R_T\sqrt{\frac{v_i}{1v_i}},\hfill & \text{else.}\hfill \end{array}$$
This method leads to additional fluctuations in the energy content of every radial interval which follow a binomial distribution. Thus, the relation
$$\frac{\sigma ^2(ϵ)}{ϵ(1ϵ)}=const=N_S^1$$
(29)
describes the contribution to radial shape fluctuations produced by the Monte Carlo method in each longitudinal integration interval. Here $`ϵ`$ denotes the energy in a given radial interval at a given shower depth:
$$ϵ_{r_1}^{r_2}f(r)𝑑r=\frac{dE(t,r)}{dE(t)}.$$
(30)
We investigated the possibility to tune $`N_S(t)`$ in each longitudinal interval to match the radial shape fluctuations observed in detailed GEANT simulations<sup>5</sup><sup>5</sup>5De Angelis et al. have used a similar method to reproduce shape fluctuations .. As an example, the quantity $`\sigma ^2(ϵ)/(ϵ(1ϵ))`$ at $`t`$ = $`56X_0`$ is displayed in Fig.11 for detailed simulations and parameterized ones without any radial shape fluctuations. The difference of these curves, which is also shown in Fig.11, is approximately constant and determines $`N_S^1`$ in equation 29 (note that the variance is additive). We found that a constant contribution to $`\sigma ^2(ϵ)/(ϵ(1ϵ))`$ can be used to match the total radial shape fluctuations to a good approximation at all shower depths.
Summing $`N_S(t)`$ over all shower depth, the total number of spots, $`N_{Spot}`$, needed for one shower can be obtained and parameterized according to
$$N_{Spot}=\mathrm{\hspace{0.17em}93}\mathrm{ln}(Z)E^{0.876}.$$
(31)
To find the number of spots for each longitudinal integration interval, the density distribution $`1/N_{Spot}dN_S(t)/dt`$ in Fig.11 is parameterized. It is described by a gamma distribution with parameters, which are given by the corresponding longitudinal energy profile:
$`T_{Spot}`$ $`=`$ $`T_{hom}(0.698+0.00212Z)\text{and}`$ (32)
$`\alpha _{Spot}`$ $`=`$ $`\alpha _{hom}(0.639+0.00334Z).`$ (33)
The total fluctuations obtained with this method are compared in Fig.12 by adding 1 RMS to the mean profiles.
Additional correlations between longitudinal and radial shower development are taken into account by introducing a correlation between the radial pdfs and the actual center of gravity,
$$t_i=\frac{\alpha _i}{\beta _i}=T_i\frac{\alpha _i}{\alpha _i1},$$
of an individual shower. This is done by replacing $`\tau `$ in equations 24, 25, and 26 by $`\tau _i`$:
$$\tau =\frac{t}{T}\tau _i=\frac{t}{t_i}\frac{\mathrm{exp}(\mathrm{ln}\alpha )}{\mathrm{exp}(\mathrm{ln}\alpha )1}.$$
(34)
The need to introduce these correlations is demonstrated in Fig.13, where integrated radial profiles are shown, which were calculated by summing over all longitudinal layers. Note that the mean integrated profiles,
$$\frac{1}{E}\frac{dE(r)}{dr},$$
are independent of energy, which is well reproduced by the parameterized simulation. The relative fluctuations of these distributions,
$$\widehat{\sigma }(r)\frac{\sigma _{RMS}}{\frac{1}{E}\frac{dE(r)}{dr}},$$
are shown using both, $`\tau `$ and $`\tau _i`$, in calculating the radial profiles. Only the simulations using $`\tau _i`$ are able to predict the fluctuations observed with GEANT correctly.
For clarity, we summarize the steps of our algorithm as follows: Determine the energy $`dE(t)`$ within one longitudinal integration interval as described in section 3.1. In case of sampling calorimeters apply sampling fluctuations on $`dE(t)`$. Evaluate the number of spots needed to reproduce radial shape fluctuations in this interval according to
$$N_S(t)=N_{Spot}_{t_{j1}}^{t_j}\frac{(\beta _{Spot}t)^{\alpha _{Spot}1}\beta _{Spot}\mathrm{exp}(\beta _{Spot}t)}{\mathrm{\Gamma }(\alpha _{Spot})}𝑑t.$$
Distribute the spots with energy $`E_S=dE(t)/N_S(t)`$ radially according to $`f(r)`$ as described above and uniformly in $`\varphi `$ and in the longitudinal interval $`\mathrm{\Delta }t`$. Finally transform the spot coordinates $`(E_S,t[X_0],r[R_M],\varphi )`$ into the detector reference system $`(E_S,x,y,z)`$.
### 3.5 Radial shower profiles – sampling calorimeters
The influence of the exact geometry on radial energy profiles is rather small. At the start of the shower the profiles look a bit smoother than in homogeneous media. With increasing shower depth they approach the shapes that are expected for homogeneous media with the appropriate effective material. These small deviations have been taken into account by the following corrections to the mean profiles:
$`R_{C,sam}`$ $`=`$ $`R_{C,hom}+z_1(1\widehat{e})+z_2F_S^1\mathrm{exp}(\tau _i)`$ (35)
$`R_{T,sam}`$ $`=`$ $`R_{T,hom}+k_1(1\widehat{e})+k_2F_S^1\mathrm{exp}(\tau _i)`$ (36)
$`p_{sam}`$ $`=`$ $`p_{hom}+(1\widehat{e})(p_1+p_2F_S^1\mathrm{exp}((\tau _i1)^2))`$ (37)
using again the sampling frequency $`F_S`$ and $`e/mip`$ (see Appendix A.2.4).
The total number of spots needed to simulate the radial shape fluctuations is much smaller than in the case of homogeneous media and no longer depends sensitively on the materials used. Instead, the spot number can be parameterized by
$$N_{Spot}=\frac{10.3}{c}E^{0.959},$$
(38)
where $`c`$ measures the sampling fluctuations according to
$$\frac{\sigma }{E}=\frac{c}{\sqrt{E}}$$
(see Fig.14). The density distribution of the spot numbers is given in analogy to the homogeneous media by:
$`T_{Spot}`$ $`=`$ $`T_{sam}(0.831+0.0019Z)\text{and}`$ (39)
$`\alpha _{Spot}`$ $`=`$ $`\alpha _{sam}(0.844+0.0026Z)`$ (40)
GEANT and parameterized simulations of mean radial profiles and their relative fluctuations,
$$\widehat{\sigma }(t,r)=\frac{\sigma _{RMS}}{\frac{1}{dE(t)}\frac{dE(t,r)}{dr}},$$
(41)
are compared in Fig.15 and Fig.16 for the H1 liquid argon calorimeter (IFE) for various energies. The influence of radial leakage on containment and energy resolution is demonstrated in Fig.17 and Fig.18. The energy independence of the energy contained in a cylinder of radius $`r`$ is well reproduced by the parameterized simulations. The energy resolution as defined in Fig.18 does not depend on radial leakage. As can be seen, this is correctly predicted by the parameterized simulation, when the correlation between the longitudinal and radial shower development is taken into account (by using $`\tau _i`$).
## 4 Comparison with data
We have compared parameterized simulations with test beam data from the H1 calorimeter, which is made of lead and liquid argon in the electromagnetic sections . Modules of the inner forward (IFE), the forward barrel (FB1), and the central barrel (CB2/CB3) calorimeters have been studied. Electron beams in the energy range between 5 and 80 GeV entered the stacks under angles of $`11^{}`$ in the IFE and CB3, and under $`35^{}`$ in the FB1 calorimeter in a test set-up at CERN.
The energy resolution of the data can be described by
$$\frac{\sigma (E)}{E}=\sqrt{\frac{c^2}{E}+\frac{b^2}{E^2}+\left(\frac{\sigma (p)}{p}\right)^2}.$$
(42)
Here $`c`$ refers to the sampling fluctuations, $`b`$ considers the noise, and $`\sigma (p)/p`$ denotes the momentum resolution of the beam. In the Monte Carlo the momentum resolution was simulated explicitly. The electronic noise was taken into account by adding random trigger events to the simulated cell energies. The constant $`c`$, which is approximately $`11\%`$ for all modules, was used to simulate the sampling fluctuations.
The simulations were carried out with the H1 detector simulation program H1FAST . The algorithms described so far are part of this program, which is used for the mass production of Monte Carlo events in the H1 detector at the HERA collider at DESY. To keep the required high precision of the parameterization also in complicated detector regions (cracks for example), the following has to be considered. If a shower develops partly inside cracks between adjacent modules, which in general cannot be approximated by a single effective medium, parameterizations will in general fail to reproduce measured signals. In H1FAST<sup>6</sup><sup>6</sup>6A stand alone version (called GFLASH 1.4) running with GEANT and covering the same functionality is available for distribution. Please contact one of the authors. we therefore do not parameterize showers, if they cross such boundaries. Only electromagnetic showers and sub-showers from hadronic interactions are parameterized which fit into one single stack.
During analysis, a $`3\sigma `$ noise cut was applied to both the experimental data and H1FAST data at the cell level, and energy clusters were built from cells containing energies above threshold. Energy distributions of the clusters with maximum channel numbers are compared in Fig.19 for all three modules considered. In addition, the energy in all other cells, not belonging to the selected clusters, are also shown.
Longitudinal profiles are shown in Fig.20 for various energies in the IFE calorimeter. The mean profiles as well as the fluctuations are nearly indistinguishable between data and H1FAST. Energy distributions in individual longitudinal layers of CB3 are compared in Fig.21 for 30 GeV incident electron energy, showing that not only the means and fluctuations but also the shape of the distributions are predicted correctly by the parameterized simulations.
Fig.22 compares lateral profiles in different shower depths in the IFE calorimeter at one energy, and in Fig.23 lateral profiles in the FB1 calorimeter, summed over all longitudinal sections, are shown for various energies. There is good agreement in the peak distributions. The tails of the profiles are dominated by electronic noise.
As shown so far, parameterized simulations can predict measured calorimeter signals very precisely, if the shower development is confined within one single calorimeter stack. Using the concept of partial parameterization as described above, the influence of cracks on the measured signal can be reproduced as shown in Fig.24. We have used test beam data scanning the crack between CB2 and CB3, which consist of two electromagnetic (CB2E, CB3E) and two hadronic stacks (CB2H, CB3H). The width of the crack is approximately 1 cm. Shown are the energies in the electromagnetic modules ($`E_{CB2E}`$, $`E_{CB3E}`$), the sum of both ($`E_{CBE}=E_{CB2E}+E_{CB3E}`$), and the sum measured in the electromagnetic and hadronic modules ($`E_{CB}=E_{CBE}+E_{CBH}`$) as a function of the beam impact position. All energies are normalized to $`E_{+20}E_{CB}(x_{calo}=20cm`$). The energy lost while scanning the crack with a 30 GeV test beam extends to about $`40\%`$, if only the electromagnetic sections are considered, and is still around $`20\%`$ if the hadronic modules are added. The agreement between data and partial parameterization is quite satisfactory.
Of the various comparisons which were made , only a limited number is presented here. Other properties, which are relevant for physics analysis with the H1 detector, like $`e/\pi `$ separation, were studied and confirm the applicability of the fast simulation.
## 5 Timing
The CPU time reduction depends on the complexity of the geometry description and the cut off parameters in the detailed simulation as well as on the type of simulated event. Fully parameterized simulations of electromagnetic showers in a simple (box) geometry are about 7000 times faster at 100 GeV (900 at 1 GeV) compared with GEANT simulations of a detailed geometry and with low energy cuts ($`e`$-cut$`=200`$ keV, $`\gamma `$-cut$`=10`$ keV).
In the framework of the H1 simulation program, partial parameterization of electromagnetic showers is performed as described above, together with individual tracking of hadrons and termination of low energy particles (see also ). The gain factors for 30 GeV showers in the H1 detector (including detailed simulations of tracker volumes) are 200 for electrons and 25 in case of hadronic showers. Medium energy cuts ($`e`$-cut$`=1`$ MeV, $`\gamma `$-cut$`=200`$ keV) were used in the corresponding detailed simulations. Complete detector simulations of HERA events (ep scattering at $`\sqrt{s}`$ = 314 GeV) require at least 10 times less CPU time using partial parameterization.
## 6 Conclusions
We have developed parameterizations of electromagnetic showers for different materials and sampling geometries. Shower to shower fluctuations and correlations are taken into account consistently, as well as correlations between the longitudinal and radial shower development. Comparisons with data have shown that parameterized simulations are able to predict measured calorimeter signals with an acceptable precision. Using the methods described above, the energy resolution is reproduced at the level of $`\pm 0.5\%`$. The energy deposited in longitudinal and lateral layers is predicted with a precision of typically $`\pm 1.5\%`$ for both, the means and the fluctuations. Using partial parameterizations, the energy measured in electromagnetic (and hadronic) modules differs by an amount of $`1.7\%`$ ($`9\%`$), if the beam enters directly into a crack. The parameterizations presented here provide a fast and precise algorithms for large scale Monte Carlo production of events for physics analysis.
## Appendix A Summary of formulae
### A.1 Homogeneous Media
#### A.1.1 Average longitudinal profiles
$`T_{hom}`$ $`=`$ $`\mathrm{ln}y0.858`$
$`\alpha _{hom}`$ $`=`$ $`0.21+(0.492+2.38/Z)\mathrm{ln}y`$
#### A.1.2 Fluctuated longitudinal profiles
$`\mathrm{ln}T_{hom}`$ $`=`$ $`\mathrm{ln}(\mathrm{ln}y0.812)`$
$`\sigma (\mathrm{ln}T_{hom})`$ $`=`$ $`(1.4+1.26\mathrm{ln}y)^1`$
$`\mathrm{ln}\alpha _{hom}`$ $`=`$ $`\mathrm{ln}\left(0.81+(0.458+2.26/Z)\mathrm{ln}y\right)`$
$`\sigma (\mathrm{ln}\alpha _{hom})`$ $`=`$ $`(0.58+0.86\mathrm{ln}y)^1`$
$`\rho (\mathrm{ln}T_{hom},\mathrm{ln}\alpha _{hom})`$ $`=`$ $`0.7050.023\mathrm{ln}y`$
#### A.1.3 Average radial profiles
$`R_{C,hom}(\tau )`$ $`=`$ $`z_1+z_2\tau `$
$`R_{T,hom}(\tau )`$ $`=`$ $`k_1\{\mathrm{exp}(k_3(\tau k_2))+\mathrm{exp}(k_4(\tau k_2))\}`$
$`p_{hom}(\tau )`$ $`=`$ $`p_1\mathrm{exp}\left\{{\displaystyle \frac{p_2\tau }{p_3}}\mathrm{exp}\left({\displaystyle \frac{p_2\tau }{p_3}}\right)\right\}`$
with
$`z_1`$ $`=`$ $`0.0251+0.00319\mathrm{ln}E`$
$`z_2`$ $`=`$ $`0.1162+0.000381Z`$
$`k_1`$ $`=`$ $`0.659+0.00309Z`$
$`k_2`$ $`=`$ $`0.645`$
$`k_3`$ $`=`$ $`2.59`$
$`k_4`$ $`=`$ $`0.3585+0.0421\mathrm{ln}E`$
$`p_1`$ $`=`$ $`2.632+0.00094Z`$
$`p_2`$ $`=`$ $`0.401+0.00187Z`$
$`p_3`$ $`=`$ $`1.313+0.0686\mathrm{ln}E`$
#### A.1.4 Fluctuated radial profiles
$`\tau _i`$ $`=`$ $`{\displaystyle \frac{t}{t_i}}{\displaystyle \frac{\mathrm{exp}(\mathrm{ln}\alpha )}{\mathrm{exp}(\mathrm{ln}\alpha )1}}`$
$`N_{Spot}`$ $`=`$ $`93\mathrm{ln}(Z)E^{0.876}`$
$`T_{Spot}`$ $`=`$ $`T_{hom}(0.698+0.00212Z)`$
$`\alpha _{Spot}`$ $`=`$ $`\alpha _{hom}(0.639+0.00334Z)`$
### A.2 Sampling Calorimeters
#### A.2.1 Material and geometry parameters
$`w_i`$ $`=`$ $`{\displaystyle \frac{\rho _id_i}{_j\rho _jd_j}}\text{ }(\rho =\text{ density})`$
$`Z_{eff}`$ $`=`$ $`{\displaystyle \underset{i}{}}w_iZ_i`$
$`A_{eff}`$ $`=`$ $`{\displaystyle \underset{i}{}}w_iA_i`$
$`{\displaystyle \frac{1}{X_{0,eff}}}`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{w_i}{X_{0,i}}}`$
$`{\displaystyle \frac{1}{R_{M,eff}}}`$ $`=`$ $`{\displaystyle \frac{1}{E_s}}{\displaystyle \underset{i}{}}{\displaystyle \frac{w_iE_{c,i}}{X_{0,i}}}\text{ }(E_s=21.2\text{ MeV})`$
$`E_{c,eff}`$ $`=`$ $`X_{0,eff}{\displaystyle \underset{i}{}}{\displaystyle \frac{w_iE_{c,i}}{X_{0,i}}}`$
$`F_S`$ $`=`$ $`{\displaystyle \frac{X_{0,eff}}{d_a+d_p}}`$
$`\widehat{e}`$ $`=`$ $`{\displaystyle \frac{1}{1+0.007(Z_pZ_a)}}`$
#### A.2.2 Average longitudinal profiles
$`T_{sam}`$ $`=`$ $`T_{hom}0.59F_S^10.53(1\widehat{e})`$
$`\alpha _{sam}`$ $`=`$ $`\alpha _{hom}0.444F_S^1`$
#### A.2.3 Fluctuated longitudinal profiles
$`\mathrm{ln}T_{sam}`$ $`=`$ $`\mathrm{ln}\left(\mathrm{exp}(\mathrm{ln}T_{hom})0.55F_S^10.69(1\widehat{e})\right)`$
$`\sigma (\mathrm{ln}T_{sam})`$ $`=`$ $`(2.5+1.25\mathrm{ln}y)^1`$
$`\mathrm{ln}\alpha _{sam}`$ $`=`$ $`\mathrm{ln}\left(\mathrm{exp}(\mathrm{ln}\alpha _{hom})0.476F_S^1\right)`$
$`\sigma (\mathrm{ln}\alpha _{sam})`$ $`=`$ $`(0.82+0.79\mathrm{ln}y)^1`$
$`\rho (\mathrm{ln}T_{sam},\mathrm{ln}\alpha _{sam})`$ $`=`$ $`0.7840.023\mathrm{ln}y`$
#### A.2.4 Average radial profiles
$`R_{C,sam}`$ $`=`$ $`R_{C,hom}0.0203(1\widehat{e})+0.0397F_S^1\mathrm{exp}(\tau )`$
$`R_{T,sam}`$ $`=`$ $`R_{T,hom}0.14(1\widehat{e})0.495F_S^1\mathrm{exp}(\tau )`$
$`p_{sam}`$ $`=`$ $`p_{hom}+(1\widehat{e})(0.3480.642F_S^1\mathrm{exp}((\tau 1)^2))`$
#### A.2.5 Fluctuated radial profiles
$`\tau _i`$ $`=`$ $`{\displaystyle \frac{t}{t_i}}{\displaystyle \frac{\mathrm{exp}(\mathrm{ln}\alpha )}{\mathrm{exp}(\mathrm{ln}\alpha )1}}`$
$`N_{Spot}`$ $`=`$ $`{\displaystyle \frac{10.3}{c}}E^{0.959}({\displaystyle \frac{\sigma }{E}}={\displaystyle \frac{c}{\sqrt{E}}})`$
$`T_{Spot}`$ $`=`$ $`T_{hom}(0.813+0.0019Z)`$
$`\alpha _{Spot}`$ $`=`$ $`\alpha _{hom}(0.844+0.0026Z)`$ |
warning/0001/physics0001056.html | ar5iv | text | # EXACT SOLUTION OF THE RESTRICTED THREE-BODY SANTILLI-SHILLADY MODEL OF 𝐻₂ MOLECULE
## 1 Introduction
In this paper, we study isochemical model of the $`H_2`$ molecule recently introduced by R. M. Santilli and D. D. Shillady , which is characterized by the conventional $`H_2`$ model set up plus a short-range attractive Hulten potential interaction between the two electrons originating from the deep overlapping of their wave functions at mutual distances of the order of 1 fm; see also . If one assumes that this attractive potential is strong enough to overcome Coloumb repulsion between the two electrons, they can form electron-electron system called isoelectronium. The isoelectronium is characterized by ”bare” mass $`M=2m_e`$, as a sum of masses of two constituent electrons, charge $`2e`$, radius about $`10^{11}`$ cm, and null magnetic moment. The used Hulten potential contains two real parameters, one of which is the isoelectronium correlation length parameter $`r_c`$, which can be treated as an effective radius of isoelectronium.
The main structural difference between the Santilli-Shillady isochemical model and the conventional quantum chemical model of the $`H_2`$ molecule, is that the former admits additional nonlinear, nonlocal, and nonpotential, thus nonunitary effects due to the deep overlapping of the wavepackets of valence electrons at short distances, which are responsible for the strong molecular bond. In a first nonrelativistic approximation, Santilli and Shillady derived the following characteristics of the isoelectronium: total rest mass $`M=2m_e`$, charge $`2e`$, magnetic moment zero, and radius $`6.84323\times 10^{11}cm`$. The value $`M=2m_e`$ of the rest mass was derived via the assumption of a contact, nonpotential interactions due to the mutual wave-overlapping sufficiently strong to overcome the repulsive Coulomb force. The nonpotential character of the bond was then responsible for the essential lack of binding energy in the isoelectronium, and the resulting value $`M=2m_e`$. However, the authors stressed in that the isoelectronium is expected to have a non-null binding energy, and, therefore, a rest mass smaller than $`2m_e`$. One argument presented in is that, when coupled in singlet at very short distances, the two electrons eventually experience very strong attractive forces of magnetic type, due to the two pairs of opposing magnetic polarities, resulting in a bond. The potential origin of the bond then implies the existence of a binding energy, resulting in a rest mass of the isoelectronium smaller than $`2m_e`$. Also, in the subsequent paper , Santilli pointed out that the isoelectronium can at most admit a small instability.
As a result of a correlation/bonding between the two electrons, Santilli and Shillady were able to reach, for the first time, representations of the binding energy and other characteristics of $`H_2`$ molecule which are accurate to the seventh digit, within the framework of numerical Hartree-Fock approach to $`H_2`$ molecule viewed as a four-body system with fixed nuclei, and with the use of Gaussian screened Coloumb potential taken as an approximation to the Hulten potential .
On the other hand, the above mentioned strong short-range character of the electron-electron interaction suggests the use of approximation of stable isoelectronium of ignorably small size, in comparison to the internuclear distance . Indeed, under these two assumptions one can reduce the conventional four-body structure of the $`H_2`$ molecule to a three-body system (the two electrons are viewed as a single point-like particle). Furthermore, in the Born-Oppenheimer approximation, i.e. at fixed nuclei, we have a restricted three-body system, the Shrödinger equation for which admits exact analytic solution.
So, we have the original four-body Santilli-Shillady model of $`H_2`$ molecule, and the three-body Santilli-Shillady model of $`H_2`$, which is an approximation to it. The former is characterized by, in general, unstable isoelectronium and, thus, sensitivity to details of the electron-electron interaction, while the latter deals with a single point-like particle (stable isoelectronium of ignorable size) moving around two fixed nuclei.
Clearly, the three-body Santilli-Shillady model of $`H_2`$ molecule can be viewed as $`H_2^+`$ ion like system. For the sake of brevity and to avoid confusion with the $`H_2^+`$ ion itself, we denote $`H_2`$ molecule, viewed as the restricted three-body system, as $`\widehat{H}_2`$. Note that $`\widehat{H}_2`$ is a neutral $`H_2^+`$ ion like system.
The quantum mechanical problem of the restricted $`H_2^+`$ ion like systems, associated differential equation, and its exact analytic solution have been studied in the literature by various authors -.
In this paper we present the exact analytic solution of the above indicated restricted three-body Santilli-Shillady isochemical model of the hydrogen molecule, study its asymptotic behavior, and analyze the ground state energy, presenting numerical results in the form of tables and plots. Our analysis is based on the analytical results obtained for thoroughly studied $`H_2^+`$ ion.
In Sec. 2, we review some features of the four-body Santilli-Shillady model of $`H_2`$ necessary for our study, and introduce our separation of variables in the Schrödinger equation under the assumption that the isoelectronium is a stable quasiparticle of ignorable size.
In Sec. 3, we review the exact analytic solution of the $`H_2^+`$ ion like systems (which includes the $`\widehat{H}_2`$ system), and study their asymptotic behavior at large and small distances between the two nuclei.
In Sec. 4, we use the preceding solution to find the binding energy of $`\widehat{H}_2`$ system. We then develop a scaling method and use Ritz’s variational approach to check the results. Both the cases of the isoelectronium ”bare” mass $`M=2m_e`$ and of variable mass parameter, $`M=\eta m_e`$ have been studied. All the data and basic results of this Section have been collected in Table 1.
In Sec. 5, we introduce a preliminary study on the application of Ritz’s variational approach to the general four-body Santilli-Shillady model of $`H_2`$, where the isoelectronium is an unstable composite particle, in which case the model re-acquires its four-body structure, yet preserves a strong bonding/correlation between the electrons.
In the Appendix, we present the results of our numerical calculations of the ground state energy of $`H_2^+`$ ion and of $`\widehat{H}_2`$ system, for different values of the isoelectronium mass parameter $`M`$, based on their respective exact solutions, in the form of tables and plots.
Our main result is that the restricted three-body Santilli-Shillady isochemical model of the hydrogen molecule does admit exact analytic solution capable of an essentially exact representation of the binding energy, although under internuclear distance about 19.6% bigger than the conventional experimental value. The mass parameter $`M`$ of isoelectronium has been used here to fit the experimental value of the binding energy, with the result $`M=0.308381m_e`$ (i.e. about 7 times less than the ”bare” mass $`M=2m_e`$). In this paper, we assume that some defect of mass effect may have place leading to decrease of the ”bare” mass $`M=2m_e`$.
We also note that the value $`M=0.308381m_e`$ implies a binding energy of about $`1.7`$ MeV, which is admittedly rather large. Recent studies by Y. Rui on the correct force law among spinning charges have indicated the existence of a critical distance below which particles with the same charge attract each others. If confirmed, these studies imply that the repulsive Coulomb force itself between two electrons in singlet coupling can be attractive at a sufficiently small distance, thus eliminating the need to postulate an attractive force sufficiently strong to overcome the repulsive Coulomb force. As a result, a binding energy in the isoelectronium structure of the order of 1.7 MeV cannot be excluded on grounds of our knowledge at this time.
Clearly, however, that due to the current lack of dynamical description of the above mentioned defect of mass, and the obtained result that the predicted internuclear distance is about 19.6% bigger than the experimental value, our study is insufficient to conclude that the isoelectronium is permanently stable, and one needs for additional study on the four-body Santilli-Shillady isochemical model of $`H_2`$, which is conducted in a subsequent paper by one of the authors .
## 2 Santilli-Shillady model of $`H_2`$ molecule
### 2.1 General equation
The Santilli-Shillady iso-Shrödinger’s equation for $`H_2`$ molecule with short-range attractive Hulten potential between the two electrons can be reduced to the following form :
$`({\displaystyle \frac{\mathrm{}^2}{2m_1}}_1^2{\displaystyle \frac{\mathrm{}^2}{2m_2}}_2^2V_0{\displaystyle \frac{e^{r_{12}/r_c}}{1e^{r_{12}/r_c}}}+{\displaystyle \frac{e^2}{r_{12}}}`$ (2.1)
$`{\displaystyle \frac{e^2}{r_{1a}}}{\displaystyle \frac{e^2}{r_{2a}}}{\displaystyle \frac{e^2}{r_{1b}}}{\displaystyle \frac{e^2}{r_{2b}}}+{\displaystyle \frac{e^2}{R}})|\varphi =E|\varphi ,`$
where $`V_0`$ and $`r_c`$ are positive constants, and $`R`$ is distance between nuclei $`a`$ and $`b`$. By using vectors of center-of-mass system of electrons 1 and 2, $`\stackrel{}{r}_a`$ and $`\stackrel{}{r}_b`$, originated at nuclei $`a`$ and $`b`$, respectively, we have
$$r_{1a}=\left|\stackrel{}{r}_a\frac{m_2}{m_1+m_2}\stackrel{}{r}_{12}\right|,r_{2a}=\left|\stackrel{}{r}_a+\frac{m_1}{m_1+m_2}\stackrel{}{r}_{12}\right|.$$
(2.2)
$$r_{1b}=\left|\stackrel{}{r}_b\frac{m_2}{m_1+m_2}\stackrel{}{r}_{12}\right|,r_{2b}=\left|\stackrel{}{r}_b+\frac{m_1}{m_1+m_2}\stackrel{}{r}_{12}\right|,$$
(2.3)
(for electrons we have $`m_1=m_2=m_e`$). The Lagrangian of the system can be written
$$=\frac{m_1\dot{r}_1^2}{2}+\frac{m_2\dot{r}_2^2}{2}V(r_{12})W(r_{1a},r_{1b},r_{2a},r_{2b},R),$$
(2.4)
Here, $`V`$ is the potential energy of interaction between the electrons 1 and 2,
$$V(r_{12})=\frac{e^2}{r_{12}}V_0\frac{e^{r_{12}/r_c}}{1e^{r_{12}/r_c}},$$
(2.5)
and $`W`$ is the potential energy of interaction between electrons and nuclei, and between two nuclei,
$$W(r_{1a},r_{1b},r_{2a},r_{2b},R)=\frac{e^2}{r_{1a}}\frac{e^2}{r_{2a}}\frac{e^2}{r_{1b}}\frac{e^2}{r_{2b}}+\frac{e^2}{R}.$$
(2.6)
Notice that $`\dot{r}_1=\dot{r}_{1a}=\dot{r}_{1b}`$, and $`\dot{r}_2=\dot{r}_{2a}=\dot{r}_{2b}`$, because $`\stackrel{}{r}_{1a}=\stackrel{}{r}_{1b}+\stackrel{}{R}`$ and $`\stackrel{}{r}_{2a}=\stackrel{}{r}_{2b}+\stackrel{}{R}`$, where $`\stackrel{}{R}`$ is constant vector. Similarly,
$$\stackrel{}{r}_a=\stackrel{}{r}_b+\stackrel{}{R},\stackrel{}{r}_a=\frac{m_1\stackrel{}{r}_{1a}+m_2\stackrel{}{r}_{2a}}{m_1+m_2},\stackrel{}{r}_b=\frac{m_1\stackrel{}{r}_{1b}+m_2\stackrel{}{r}_{2b}}{m_1+m_2}.$$
(2.7)
Then, Lagrangian (2.4) can be rewritten as $`=(r_a,r_b,r_{12})`$,
$$=\frac{M\dot{r}_a^2}{2}+\frac{m\dot{r}_{12}^2}{2}V(r_{12})W(r_a,r_b,r_{12},R).$$
(2.8)
Here, $`M=m_1+m_2`$ is the total mass of the electrons, and $`m=m_1m_2/(m_1+m_2)`$ is the reduced mass. Corresponding generalized momenta take the form
$$\stackrel{}{P}_M=\frac{}{\dot{\stackrel{}{r}_A}}=M\dot{\stackrel{}{r}_A}.\stackrel{}{p}_m=\frac{}{\dot{\stackrel{}{r}_{12}}}=m\dot{\stackrel{}{r}_{12}}.$$
(2.9)
The system reveals axial symmetry, with the axis connecting two nuclei. Also, for identical nuclei we have reflection symmetry in respect to the plane perpendicular to the above axis and lying on equal distances from the two nuclei.
### 2.2 Separation of variables
Santilli and Shillady then assume that, as a particular case under study in this paper (not to be confused with the general four-body case), the two valence electrons of the $`H_2`$ molecule can form a stable quasi-particle of small size due to short-range attractive Hulten potential, such that
$$r_{12}r_a,r_{12}r_b.$$
(2.10)
Therefore, we can ignore $`r_{12}`$ in Eqs.(2.2) and (2.3),
$$\stackrel{}{r}_{1a}\stackrel{}{r}_{2a}\stackrel{}{r}_a,\stackrel{}{r}_{1b}\stackrel{}{r}_{2b}\stackrel{}{r}_b.$$
(2.11)
The Hamiltonian of the system then becomes
$$\widehat{H}=\frac{\widehat{P}_M^2}{2M}+\frac{\widehat{p}_m^2}{2m}+V(r_{12})+W(r_a,r_b,R),$$
(2.12)
where
$$W(r_a,r_b,R)=\frac{2e^2}{r_a}\frac{2e^2}{r_b}+\frac{2e^2}{R}.$$
(2.13)
In this approximation, it is possible to separate the variables $`r_{a,b}`$ and $`r_{12}`$. Namely, inserting $`|\varphi =\psi (r_a,r_b,R)\chi (r_{12})`$ into the equation
$$\left(\frac{\mathrm{}^2}{2M}_{ab}^2\frac{\mathrm{}^2}{2m}_{12}^2V_0\frac{e^{r_{12}/r_c}}{1e^{r_{12}/r_c}}+\frac{e^2}{r_{12}}\frac{2e^2}{r_a}\frac{2e^2}{r_b}+\frac{e^2}{R}\right)|\varphi =E|\varphi $$
(2.14)
we obtain
$$\frac{\mathrm{}^2}{2M}\frac{_{ab}^2\psi }{\psi }\frac{\mathrm{}^2}{2m}\frac{_{12}^2\chi }{\chi }+V(r_{12})+W(r_a,r_b,R)E=0.$$
(2.15)
By separating the variables, we have the following two equations:
$$\frac{\mathrm{}^2}{2m}_{12}^2\chi +V(r_{12})\chi =\epsilon \chi .$$
(2.16)
and
$$\frac{\mathrm{}^2}{2M}_{ab}^2\psi +W(r_a,r_b,R)\psi =(E\epsilon )\psi .$$
(2.17)
In this way, under approximation (2.10), the original four-body problem is reduced to a three-body problem characterized by two differential equations:
1) Equation (2.16), which describes the electron-electron system forming the bound quasi-particle state called isoelectronium, with ”bare” total mass $`M=2m_e`$ and charge $`2e`$. This equation will not be studied in this paper, since we assume that the isoelectronium is permanently stable.
2) Equation (2.17), which is the structural equation of the restricted three-body Santilli-Shillady isochemical model $`\widehat{H}_2`$, in which the stable isoelectronium with ”bare” mass $`M=2m_e`$, charge $`2e`$, null magnetic moment and ignorable size orbits around the two nuclei, hereon assumed to have infinite mass (the Born-Oppenheimer approximation).
This paper is devoted to the study of the exact analytic solution of the latter equation, and its capability to represent the experimental data on the binding energy, bond length, and other characteristics of the hydrogen molecule.
## 3 Exact solution for $`H_2^+`$ ion like system
In this Section, we present analytical solution of the Schrödinger equation for $`H_2^+`$ ion-like systems in Born-Oppenheimer approximation, we analyze the associated recurrence relations, and asymptotic behavior of the solutions at large and small distances between the two nuclei. As it was indicated , this problem arises when Santilli-Shillady model of $`H_2`$ is reduced to the restricted three-body problem characterized by Eq. (2.17), which possesses exact solution under appropriate separation of variables.
### 3.1 Differential equations
In Born-Oppenheimer approximation, i.e., at fixed nuclei, the equation for $`H_2^+`$ ion-like system for a particle of mass $`M`$ and charge $`q`$ is
$$^2\psi +2M(E+\frac{q}{r_a}+\frac{q}{r_b})\psi =0.$$
(3.1)
In spheroidal coordinates,
$$x=\frac{r_a+r_b}{R},1<x<\mathrm{},$$
(3.2)
$$y=\frac{r_ar_b}{R},1<y<1,$$
(3.3)
$$\phi ,0<\phi <2\pi ,$$
(3.4)
where $`R`$ is a fixed separation distance between the nuclei $`a`$ and $`b`$, and
$`^2={\displaystyle \frac{4}{R^2(x^2y^2)}}\left({\displaystyle \frac{}{x}}(x^21){\displaystyle \frac{}{x}}+{\displaystyle \frac{}{y}}(1y^2){\displaystyle \frac{}{y}}\right)`$ (3.5)
$`+{\displaystyle \frac{1}{R^2(x^21)(1y^2)}}{\displaystyle \frac{^2}{\phi ^2}}.`$
We then have from Eq.(3.1)
$`[{\displaystyle \frac{}{x}}(x^21){\displaystyle \frac{}{x}}+{\displaystyle \frac{}{y}}(1y^2){\displaystyle \frac{}{y}}+{\displaystyle \frac{x^2y^2}{4(x^21)(1y^2)}}{\displaystyle \frac{^2}{\phi ^2}}`$ (3.6)
$`+{\displaystyle \frac{MER^2}{2}}(x^2y^2)+2MqRx]\psi =0.`$
Here, we have used
$$\frac{1}{r_a}+\frac{1}{r_b}=\frac{4}{R}\frac{x}{x^2y^2}.$$
(3.7)
Obviously, Equation (3.6) can be separated by the use of the representation
$$\psi =f(x)g(y)e^{im\phi },$$
(3.8)
under which we have two second-order ordinary differential equations,
$$\frac{d}{dx}\left((x^21)\frac{d}{dx}f\right)\left(\lambda 2MqRx\frac{MER^2}{2}x^2+\frac{m^2}{x^21}\right)f=0,$$
(3.9)
$$\frac{d}{dy}\left((1y^2)\frac{d}{dy}g\right)+\left(\lambda \frac{MER^2}{2}y^2\frac{m^2}{1y^2}\right)g=0,$$
(3.10)
where $`\lambda `$ is a separation constant (cf. ). So, the problem is to identify solutions for $`f`$ and $`g`$.
### 3.2 Recurrence relations
By introducing the re-formulations
$$f(x^21)^{m/2}f,$$
(3.11)
$$g(1y^2)^{m/2}g,$$
(3.12)
to handle singularities at $`x=\pm 1`$ and $`y=\pm 1`$ in Eqs.(3.9) and (3.10), respectively, we reach the following final form of the equations to be solved:
$$(x^21)f^{\prime \prime }+2(m+1)f^{}(\lambda +m(m+1)\stackrel{~}{a}xc^2x^2)f=0$$
(3.13)
and
$$(1y^2)g^{\prime \prime }2(m+1)g^{}+(\lambda m(m+1)c^2y^2)g=0,$$
(3.14)
where we have denoted
$$c^2=\frac{MER^2}{2},\stackrel{~}{a}=2MqR.$$
(3.15)
We shall look for solutions in the form of power series. Substituting the power series
$$f=f_kx^k,$$
(3.16)
$$g=g_ky^k,$$
(3.17)
into Eqs. (3.13) and (3.14), we obtain the recurrence relations,
$$c^2f_{n2}+\stackrel{~}{a}f_{n1}(\lambda (m+n)(m+n+1))f_n$$
(3.18)
$$(n+1)(n+2)f_{n+2}=0$$
and
$$c^2g_{n2}(\lambda (m+n)(m+n+1))g_n(n+1)(n+2)g_{n+2}=0,$$
(3.19)
from which coefficients $`f_k`$ and $`g_k`$ must be found. Here, $`f_0`$ and $`g_0`$ are fixed by normalization of the general solution. Note that the recurrence relation (3.18) contains term $`2MqRf_{n1}`$ raised from the linear term $`2MqRx`$ in Eq.(3.9).
In the next two Sections we consider some particular cases of interest prior to going into details of the general solution. These particular solutions are important for the study of the general case.
### 3.3 The particular case $`R=0`$
In the particular case $`R=0`$, the two nuclei are superimposed, so that the system is reduced to a helium-like system,
$$\frac{d}{dx}\left((x^21)\frac{d}{dx}f\right)\left(\lambda +\frac{m^2}{x^21}\right)f=0,$$
(3.20)
$$\frac{d}{dy}\left((1y^2)\frac{d}{dy}g\right)+\left(\lambda \frac{m^2}{1y^2}\right)g=0.$$
(3.21)
From recurrence relations (3.18) and (3.19) we obtain the following particular recurrence sequences,
$$(\lambda (m+n)(m+n+1))f_n(n+1)(n+2)f_{n+2}=0$$
(3.22)
and
$$(\lambda (m+n)(m+n+1))g_n(n+1)(n+2)g_{n+2}=0,$$
(3.23)
which are equivalent to each other, and can be stopped by putting the separation constant
$$\lambda =(m+n)(m+n+1)=l(l+1),$$
(3.24)
with $`m=l,\mathrm{},l.`$ This gives us well known solution for $`g`$ in terms of Legendre polynomials,
$$g=(1y^2)^{m/2}\frac{d}{dy^m}P_l(y),$$
(3.25)
where $`m=|m|`$, and
$$P_l=\frac{1}{2^ll!}\frac{d^l}{dy^l}(y^21)^l.$$
(3.26)
The solution is the well known spherical harmonic function
$$Y_{lm}=N_{lm}P_l^m(y)e^{im\phi },$$
(3.27)
with normalization constant
$$N_{lm}=\sqrt{\frac{(lm)!(2l+1)}{(l+m)!4\pi }}.$$
(3.28)
This solution corresponds to the case of an ellipsoid degenerated into a sphere, and we can put $`y=\mathrm{cos}\theta `$ for identification with the angular spherical coordinates $`(\theta ,\phi )`$. Equation in $`x`$ corresponds to the radial part of the well known solution expressed in terms of Laguerre polynomials.
### 3.4 The particular case $`q=0`$
In the particular case of zero charge, $`q=0`$, we have from Eqs.(3.9) and (3.10)
$$\frac{d}{dx}\left((x^21)\frac{d}{dx}f\right)\left(\lambda c^2x^2+\frac{m^2}{x^21}\right)f=0,$$
(3.29)
$$\frac{d}{dy}\left((1y^2)\frac{d}{dy}g\right)+\left(\lambda c^2y^2\frac{m^2}{1y^2}\right)g=0.$$
(3.30)
One can see that these equations originate straightforwardly also from the standard wave equation $`^2\psi +k^2\psi =0`$, in the spheroidal coordinates $`(x,y,\phi )`$. Recurrence relations (3.18) and (3.19) then become
$$c^2f_{n2}(\lambda (m+n)(m+n+1))f_n(n+1)(n+2)f_{n+2}=0$$
(3.31)
and
$$c^2g_{n2}(\lambda (m+n)(m+n+1))g_n(n+1)(n+2)g_{n+2}=0,$$
(3.32)
which are equivalent to each other.
A general solution for $`f`$ is given by linear combinations of radial spheroidal functions $`R_{mn}^{(p)}(c,x)`$ of first, $`p=1`$, and second, $`p=2`$, kind ,
$$R_{mn}^{(p)}(c,x)=\left\{\underset{r=0,1}{\overset{\mathrm{}^{}}{}}\frac{(2m+r)!}{r!}d_r^{mn}\right\}^1\left(\frac{x^21}{x^2}\right)^{m/2}\times $$
(3.33)
$$\times \underset{r=0,1}{\overset{\mathrm{}^{}}{}}i^{r+mn}\frac{(2m+r)!}{r!}d_r^{mn}Z_{m+r}^{(p)}(cx),$$
where,
$$Z_n^{(1)}(z)=\sqrt{\frac{\pi }{2z}}J_{n+1/2}(z),$$
(3.34)
$$Z_n^{(2)}(z)=\sqrt{\frac{\pi }{2z}}Y_{n+1/2}(z),$$
(3.35)
and $`J_{n+1/2}(z)`$ and $`Y_{n+1/2}(z)`$ are Bessel functions of first and second kind, respectively. The sum in (3.33) is made over either even or odd values of $`r`$ depending on the parity of $`nm`$. Asymptotics of $`R_{mn}^{(1)}(c,x)`$ and $`R_{mn}^{(2)}(c,x)`$ are
$$R_{mn}^{(1)}(c,x)\stackrel{cx\mathrm{}}{}\frac{1}{cx}\mathrm{cos}\left[cx\frac{1}{2}(n+1)\pi \right],$$
(3.36)
$$R_{mn}^{(2)}(c,x)\stackrel{cx\mathrm{}}{}\frac{1}{cx}\mathrm{sin}\left[cx\frac{1}{2}(n+1)\pi \right].$$
(3.37)
Particularly, to have well defined limit at $`x=0`$ we should use only spheroidal function of first kind, $`R_{mn}^{(1)}(c,x)`$, because Bessel function of second kind, $`Y_n(z)`$, has logarithmic divergence at $`z=0`$.
General solution for $`g`$ is given by linear combination of angular spheroidal functions of first and second kind ,
$$S_{mn}^{(1)}(c,y)=\underset{r=0,1}{\overset{\mathrm{}^{}}{}}d_r^{mn}(c)P_{m+r}^m(y),$$
(3.38)
$$S_{mn}^{(2)}(c,y)=\underset{r=\mathrm{}}{\overset{\mathrm{}^{}}{}}d_r^{mn}(c)Q_{m+r}^m(y),$$
(3.39)
where $`P_n^m(y)`$ and $`Q_n^m(y)`$ are the associated Legendre polynomials of first and second kind, respectively.
Expressions for radial and angular spheroidal functions, and corresponding eigenvalues $`\lambda `$, for particular values of $`m`$ and $`n`$, are presented in Ref. .
Coefficients $`d_k^{mn}(c)`$ are calculated with the help of the following recurrence relation:
$$\alpha _kd_{k+2}+(\beta _k\lambda _{mn})d_k+\gamma _kd_{k2}=0,$$
(3.40)
where
$$\alpha _k=\frac{(2m+k+2)(2m+k+1)c^2}{(2m+2k+3)(2m+2k+5)},$$
(3.41)
$$\beta _k=(m+k)(m+k+1)+\frac{2(m+k)(m+k+1)2m^21}{(2m+2k1)(2m+2k+3)}c^2,$$
(3.42)
$$\gamma _k=\frac{k(k1)c^2}{(2m+2k3)(2m+2k1)}.$$
(3.43)
The calculation is made by the following procedure. First, one calculates $`N_r^m`$,
$$N_{r+2}^m=\gamma _r^m\lambda _{mn}\frac{\beta _r^m}{N_r^m}(r2),$$
(3.44)
$$N_2^m=\gamma _0^m\lambda _{mn};N_3^m=\gamma _1^m\lambda _{mn},$$
(3.45)
$$\gamma _r^m=(m+r)(m+r+1)+\frac{1}{2}c^2\left[1\frac{4m^21}{(2m+2r1)(2m+2r+3)}\right](r0).$$
(3.46)
Second, one calculates the fractions $`d_0/d_{2r}`$ and $`d_1/d_{2p+1}`$ with the use of
$$\frac{d_0}{d_{2r}}=\frac{d_0}{d_2}\frac{d_2}{d_4}\mathrm{}\frac{d_{2r2}}{d_{2r}},$$
(3.47)
$$\frac{d_1}{d_{2p+1}}=\frac{d_1}{d_3}\frac{d_3}{d_5}\mathrm{}\frac{d_{2p1}}{d_{2p+1}},$$
(3.48)
and
$$N_r^m=\frac{(2m+r)(2m+r1)c^2}{(2m+2r1)(2m+2r+1)}\frac{d_r}{d_{r2}}$$
(3.49)
The coefficients $`d_0`$, for even $`r`$, and $`d_1`$, for odd $`r`$, are determined via the normalization of the solution.
### 3.5 The general case
In this Section, we consider the general solution of our basic equations (3.9) and (3.10). To have more general set up, we consider the case of different charges of nuclei, $`Z_1`$ and $`Z_2`$. This leads to appearance of additional linear in $`y`$ term in Eq.(3.10), so that both the ordinary differential equations become of similar structure. Also, we restrict consideration by analyzing discrete spectrum, i.e. we assume that the energy $`E<0`$.
Let us denote
$$p=\frac{R}{2}\sqrt{2E},a=R(Z_2+Z_1),b=R(Z_2Z_1).$$
(3.50)
Then, Eqs.(3.9) and (3.10), for the general case of different charges of nuclei, can be written as
$$\frac{d}{dx}\left((x^21)\frac{d}{dx}f_{mk}(p,a;x)\right)$$
$$+\left(\lambda _{mk}^{(x)}p^2(x^21)+ax\frac{m^2}{x^21}\right)f_{mk}(p,a;x)=0,$$
(3.51)
$$\frac{d}{dy}\left((1y^2)\frac{d}{dy}g_{mq}(p,b;y)\right)$$
$$+\left(\lambda _{mq}^{(y)}p^2(1y^2)+by\frac{m^2}{1y^2}\right)g_{mq}(p,b;y)=0,$$
(3.52)
where we assume that the solutions obey
$$|f_{mk}(p,a;1)|<\mathrm{},\underset{x\mathrm{}}{lim}f_{mk}(p,a;x)=0,|g_{mq}(p,b;\pm 1)|<\mathrm{}.$$
(3.53)
The eigenvalues $`\lambda `$ in Eqs.(3.51) and (3.52) should be equal to each other,
$$\lambda _{mk}^{(x)}(p,a)=\lambda _{mq}^{(y)}(p,b).$$
(3.54)
The general solution $`\psi (x,y,\phi )`$ of Eq.(3.6) is represented in the following factorized form:
$$\psi _{kqm}(x,y,\phi ;R)=N_{kqm}(p,a,b)f_{mk}(p,a;x)g_{mq}(p,b;y)\frac{\mathrm{exp}(\pm im\phi )}{\sqrt{2\pi }}.$$
(3.55)
The normalization coefficients $`N_{kqm}(p,a,b)`$ in Eq.(3.55) are represented with the help of derivatives of the eigenvalues, $`\lambda _{mk}^{(x)}(p,a)`$ and $`\lambda _{mq}^{(y)}(p,b)`$, namely,
$$N_{kqm}^2(p,a,b)=\frac{16p}{R^3}\left[\frac{\lambda _{mq}^{(y)}(p,b)}{p}\frac{\lambda _{mk}^{(x)}(p,a)}{p}\right]^1.$$
(3.56)
For a given indices $`k`$, $`q`$, $`m`$, and fixed values of $`Z_1`$, $`Z_2`$, and $`R`$, the discrete energy spectrum $`E`$ can be determined from Eq.(3.54). This equation has unique solution, $`p=p_{kqm}(a,b)`$. Then, by solving the relation stemming from (3.50)
$$p_{kqm}(R(Z_2+Z_1),R(Z_2Z_1))=\frac{R}{2}\sqrt{2E}$$
(3.57)
in respect to $`E`$, we can find the discrete spectrum of energy,
$$E_j(R)=E_{kqm}(R,Z_1,Z_2).$$
(3.58)
Number of zeroes, $`k`$, $`q`$, and $`m`$, of the functions $`g(y)`$, $`f(x)`$, and $`\mathrm{exp}\pm im\phi `$ are the angular, radial and azimuthal quantum numbers, respectively. However, instead of $`k`$, $`q`$, and $`m`$ one can use their linear combinations, namely, $`N=k+q+m+1`$ is main quantum number and $`l=q+m`$ is orbital quantum number.
To construct the general solution $`u(z)`$, which is called Coloumb spheroidal function (csf), in terms of angular csf $`g(y)`$ and radial csf $`f(x)`$, let us, again, use the form which accounts for singularities at the points $`z=\pm 1`$ and $`z=\mathrm{}`$,
$$u(z)=(1z^2)^{m/2}\mathrm{exp}[p(1\pm z)]v(z).$$
(3.59)
Then, we represent $`v(z)`$ as an expansion,
$$v(z)=\underset{s=0}{\overset{\mathrm{}}{}}a_s(p,b,\lambda )w_s(z),$$
(3.60)
in some set of basis functions $`w_s(z)`$.
Now, the complexity of the recurrence relations depends on the basis. In the preceding sections, where the particular cases, $`R=0`$ and $`q=0`$, have been considered, we used a power series representation. One can try other forms of the representation as well. For a good choice of the basis functions $`w_s(z)`$, we can obtain three-term recurrence relation of the form
$$\alpha _sa_{s+1}\beta _sa_s+\gamma _sa_{s1}=0,$$
(3.61)
where $`\alpha _s`$, $`\beta _s`$, and $`\gamma _s`$ are some polynomials in $`p`$, $`b`$, and $`\lambda `$. Then, using the tridiagonal matrix $`\widehat{A}`$ consisting of the coefficients $`\alpha _s`$, $`\beta _s`$, and $`\gamma _s`$ entering Eq.(3.61), we can write down the equation to find out eigenvalues $`\lambda _{mk}^{(x)}(p,a)`$ and $`\lambda _{mq}^{(y)}(p,b)`$. Namely,
$$\mathrm{det}\widehat{A}=F(p,b,\lambda )=0.$$
(3.62)
The matrix $`\widehat{A}`$ has a tridiagonal form. This leads directly to one-to-one correspondence between $`\mathrm{det}\widehat{A}`$ and the infinite chain fraction,
$$F(p,b,\lambda )=\beta _0\frac{\alpha _0\gamma _1}{\beta _1{\displaystyle \frac{\alpha _1\gamma _2}{\beta _2\mathrm{}{\displaystyle \frac{\alpha _N\gamma _{N+1}}{\beta _{N+1}\mathrm{}}}}}}=\beta _0\frac{\alpha _0\gamma _1}{\beta _1}\frac{\alpha _1\gamma _2}{\beta _2}\mathrm{}\frac{Q_N}{P_N}.$$
(3.63)
In numerical computations, this relation allows one to find out eigenvalues $`\lambda `$ in an easier way due to simpler algorithm provided by the chain fraction. Consequently, one can compute the energy and coefficients $`a_s`$ of the expansion of eigenfunctions $`g(y)`$ and $`f(x)`$ by using the chain fraction.
The result of this approach in constructing of the solutions depends on the convergence of the chain fraction. Analysis of the convergence can be made from a general point of view. Sufficient conditions of the convergence of the chain (3.63), and of the expansion (3.60), are the following two relations:
$$\left|\frac{\alpha _{s1}\gamma _s}{\beta _{s1}\beta _s}\right|<\frac{1}{4},\frac{a_{s+1}}{a_s}_{|s\mathrm{}}\frac{\beta _s}{2\alpha _s}\left[1\left(14\frac{\alpha _s\gamma _s}{\beta _s^2}\right)^{1/2}\right].$$
(3.64)
Further analysis of the convergence depends on specific choice of the basis functions $`u_s(z)`$.
(i) Series expansion, $`v_s(z)=z^s`$. In this case, the radius, $`Z_v`$, of convergence is
$$Z_v=\underset{s\mathrm{}}{lim}\left|\frac{a_s}{a_{s+1}}\right|.$$
(3.65)
Particularly, when $`a_{s+1}/a_s0`$ at $`s\mathrm{}`$ the series (3.60) converges at any $`z`$.
(ii) For the choice of basis function $`v_s(z)`$ in the form of orthogonal polynomials, the sufficient condition for convergence of Fourier series (3.60) is
$$\left|\frac{a_s}{a_{s+1}}\right|_s\mathrm{}1\frac{1}{s}.$$
(3.66)
Below, we consider separately angular and radial csf entering the general solution.
#### 3.5.1 The angular Coloumb spheroidal function
For the angular Coloumb spheroidal function (acsf), it is natural to choose the basis functions $`v_s(y)`$ in the form of associated Legendre polynomials, $`P_{s+m}^m(y)`$. Indeed, they form complete system in the region $`y[1,1]`$, and reproduce acsf at $`p=b=0`$ (see Sec. 3.3). Inserting of the expansion
$$g_{mq}(p,b;y)=\underset{s=0}{\overset{\mathrm{}}{}}c_sP_{s+m}^m(y)$$
(3.67)
into Eq.(3.52) entails five-term recurrence relation. However, this relation, which is sometimes used, is not so suitable as the three-term relation. This is because the determinant of the corresponding pentadiagonal matrix can not be represented as a chain fraction. Nevertheless, in the case $`b=0`$, i.e. for $`Z_1=Z_2`$, this five-terms recurrence relation is reduced to two three-terms recurrence relations, separately for even ($`c_2=0`$, $`c_0=1`$) and odd ($`c_1=0`$, $`c_1=1`$) solutions of Eq.(3.52) presented in previous Section.
For the general case $`b0`$, the expansions of $`g(p,b;y)`$, handling singularities at the points $`y=\pm 1`$ and $`y=\mathrm{}`$, respectively, as considered by Baber and Hasse , are
$$g_{mq}(p,b;y)=\mathrm{exp}[p(1+y)]\underset{s=0}{\overset{\mathrm{}}{}}c_sP_{s+m}^m(y),$$
(3.68)
$$g_{mq}(p,b;y)=\mathrm{exp}[p(1y)]\underset{s=0}{\overset{\mathrm{}}{}}c_s^{}P_{s+m}^m(y),$$
(3.69)
These expansions yield three-terms recurrence relation,
$$\rho _sc_{s+1}\kappa _sc_s+\delta _sc_{s1}=0,c_1=0,$$
(3.70)
where the coefficients for the case of expansion (3.68) have the following form:
$$\rho _s=\frac{(s+2m+1)[b2p(s+m+1)]}{2(s+m)+3},$$
$$\kappa _s=(s+m)(s+m+1)\lambda ,$$
(3.71)
$$\delta _s=\frac{s[b+2p(s+m)]}{2(s+m)1}.$$
To estimate convergence of these expansions, one can use the above made estimation of the convergence, with the following replacements: $`\alpha _s\rho _s`$, $`\beta _s\kappa _s`$, $`\gamma _s\delta _s`$, and $`a_sc_s`$. For the expansion (3.79) we have
$$\left|\frac{\rho _{s1}\delta _s}{\kappa _{s1}\kappa _s}\right|_s\mathrm{}\left(\frac{p}{s}\right)^2,$$
(3.72)
i.e., at $`p>1`$, convergence takes place only at $`s>2p`$. We should take into account this condition when choosing minimal number of terms in the chain fraction (3.73) which is sufficient to calculate $`\lambda `$, to a required accuracy.
The recurrence relation for the coefficients $`c_s^{}`$ of the expansion (3.69) differs from that of Eq.(3.70) by the replacement $`pp`$ in formulas (3.71). Clearly, this replacement does not change the form of the chain fraction,
$$F^{(y)}(p,b,\lambda )=\kappa _0\frac{\rho _0\delta _1}{\kappa _1}\frac{\rho _1\delta _2}{\kappa _2}\mathrm{}$$
(3.73)
So, in both the cases, (3.68) and (3.69), the eigenvalues $`\lambda `$ can be found from one and the same equation,
$$F^{(y)}(p,b,\lambda )=0.$$
(3.74)
In practical calculations with the help of this algorithm, the infinite chain fraction (3.63) is, of course, replaced by the finite one, $`F_{N+1}^{(y)}(p,b,\lambda )`$, in which one retains a sufficiently big number $`N`$ of terms. Typically, $`N>10`$ provides very good accuracy. So, the eigenvalues are computed as the roots of the polynomial $`Q_{N+1}(p,b,\lambda )`$ of degree $`N+1`$, namely,
$$F_{N+1}^{(y)}(p,b,\lambda )=\frac{Q_{N+1}(p,b,\lambda )}{P_{N+1}(p,b,\lambda )}.$$
(3.75)
Such a representation allows one to exclude singularities, associated to zeroes of the polynomial $`P_{N+1}(p,b,\lambda )`$, from Eq.(3.74). Further, from the definitions (3.63) and (3.75) we obtain the following recurrence relation for the polynomial $`Q_k(p,b,\lambda )`$:
$$Q_{k+1}=Q_k\overline{\kappa }_{Nk}Q_{k1}\overline{\rho }_{Nk}\overline{\delta }_{Nk+1},Q_1=0,Q_0=1,$$
(3.76)
with the use of which one can find $`Q_{N+1}`$. Here, the coefficients $`\overline{\kappa }_s`$, $`\overline{\rho }_s`$, and $`\overline{\delta }_s`$ differ from that of Eq.(3.71) by the factor $`(1+\kappa _s^2)^{1/2}`$. This factor does not change the recurrence relation (3.70). However, it makes possible to avoid accumulating of big numbers at intermediate computations. Indeed, from Eq.(3.71) for $`\kappa _s`$ it follows that the leading coefficients of the polynomials $`Q_k`$ would behave as $`k^{4k}`$, for example, for $`k=4`$ we would have $`4^{16}`$, if we would not made the above mentioned renormalization of the coefficients $`\rho _s`$, $`\kappa _s`$, and $`\delta _s`$. The eigenvalue is found as an appropriate root of the polynomial $`Q_{N+1}(p,b,\lambda ^{(y)})`$. Clearly, for big $`N`$, there is no way to represent in general the roots of $`Q_{N+1}`$ analytically so one is forced to use numerical computations.
In the numerical computations, to pick up the appropriate eigenvalue $`\lambda _{mq}^{(y)}(p,b)`$ among $`N+1`$ roots of the polynomial $`Q_{N+1}`$ it is necessary to choose some starting value of $`\lambda `$. For example, one can put the starting value at the point $`p=b=0`$, where $`\lambda _{mq}^{(y)}(0,0)=(q+m)(q+m+1)`$. The first step is to increase discretely $`pp+\mathrm{\Delta }p`$ and $`bb+\mathrm{\Delta }b`$ beginning from the starting point $`p=b=0`$, at fixed values of $`m`$ and $`q`$, and the second step is to find $`\lambda _{mq}^{(y)}(p+\mathrm{\Delta }p,b+\mathrm{\Delta }b)`$ with the help of Eq.(3.74). Repeating these steps one can find $`\lambda _{mq}^{(y)}`$ numerically as a function of $`p`$ and $`b`$ in some interval of interest.
Also, asymptotics of $`\lambda `$ which will be studied in Sec. 3.6 are of much help here to choose the appropriate root. For example, for $`b=0`$ and $`N=5`$ we obtain numerically from the determinant of the tridiagonal matrix consisting of the coefficients defined by Eq.(3.71), with $`\kappa _s`$ on the main digonal, and $`\rho _s`$ and $`\delta _s`$ on the upper and lower adjacent diagonals respectively, the polynomial,
$$\text{det}\widehat{A}=0.003\lambda ^60.2\lambda ^5+(0.2p^2+5.5)\lambda ^4(6.3p^2+56)\lambda ^3+(1.5p^4+66p^2+231)\lambda ^2$$
(3.77)
$$(19p^4+226p^2+277)\lambda +p^6+44p^4+186p^2.$$
Only one of its six roots has asymptotics,
$$\lambda _{|p0}=0.667p^20.0148p^4+O(p^5),$$
(3.78)
which reproduces, to a good accuracy, the asymptotics (3.124). So, this is the desired root to be used in subsequent calculations. Also, observe the decrease of the numerical coefficients at higher degrees of $`\lambda `$ which control the convergence.
Note that, at $`p1`$, the acsf is concentrated around the points $`y=\pm 1`$ so that expansion (3.68) converges slowly. In this case one uses another, more appropriate, expansions,
$$g_{mq}(p,b;y)=(1y^2)^{m/2}\mathrm{exp}[p(1+y)]\underset{s=0}{\overset{\mathrm{}}{}}c_s(1+y)^s,$$
(3.79)
$$g_{mq}(p,b;y)=(1y^2)^{m/2}\mathrm{exp}[p(1y)]\underset{s=0}{\overset{\mathrm{}}{}}c_s^{}(1y)^s.$$
(3.80)
Evidently, expansion (3.79) converges faster in the region $`[1,0]`$ while the expansion (3.80) converges faster in the region $`[0,1]`$. Here, the coefficients $`c_s`$ of the expansion (3.79) obey the three-term recurrence sequence (3.70), with
$$\rho _s=2(s+1)(s+m+1),$$
$$\kappa _s=s(s+1)+(2s+m+1)(2p+m)+b\lambda ,$$
(3.81)
$$\delta _s=b+2p(s+m).$$
It is remarkable to note that expansions (3.79) and (3.80) converge at any $`y`$, and the corresponding chain fractions (3.63) converge at any $`p`$ since
$$\frac{c_{s+1}}{c_s}_{|s\mathrm{}}\frac{2p}{s},\left|\frac{\rho _{s1}\delta _s}{\kappa _{s1}}\kappa _s\right|_s\mathrm{}\frac{4p}{s}.$$
(3.82)
Similarly, the coefficients $`c_s^{}`$ obey the same relation, with the replacement $`bb`$ in Eq.(3.81).
In practical calculations, one can use a combination of expansions (3.68) and (3.79). Namely, the procedure is: from expansion (3.68) one finds eigenvalues while the eigenfunctions are calculated from to Eq.(3.79). Of course, both solutions (3.79) and (3.80) should be sewed, for example, at the point $`y=0`$, because the recurrence relations do not determine, in this case, a general normalization of the coefficients $`c_s`$ and $`c_s^{}`$. Particularly, the sewing condition, which defines the normalization of $`c_s`$ and $`c_s^{}`$, has the form
$$\underset{s=0}{}c_s=\underset{s=0}{}c_s^{}.$$
(3.83)
To derive the asymptotics of acsf and its eigenvalues we can use an expansion in Laguerre polynomials,
$$g_{mq}(p,b;y)=(1y^2)^{m/2}\mathrm{exp}[p(1\pm y)]\underset{s=0}{\overset{\mathrm{}}{}}c_sL_{s+m}^m(2p(1\pm y)),$$
(3.84)
$$L_n^m(z)=\frac{e^zz^m}{n!}\frac{d^n}{dz^n}(e^zz^{n+m}).$$
(3.85)
The insertion of this expansion into Eq.(3.52) and the use of the differential equation for Laguerre polynomials,
$$z\frac{d^2}{dz^2}L_n^m(z)+(1z+m)\frac{d}{dz}L_n^m(z)+nL_n^m(z)=0$$
(3.86)
yield recurrence relation (3.70). For the case of positive sign in Eq.(3.84), we should put
$$\rho _s=(s+m+1)\left(s+1+\frac{b}{2p}\right),$$
$$\kappa _s=(2s+m+1)\left(s+m+1+\frac{b}{2p}2p\right)+(s+m)(m+1)+b\lambda ,$$
(3.87)
$$\delta _s=s\left(s+m+\frac{b}{2p}\right).$$
#### 3.5.2 The radial Coloumb spheroidal function
The radial Coloumb spheroidal function (rcsf) obviously should be written in a form suitable to handle singularities at the points $`x=1`$ and $`x=\mathrm{}`$, namely,
$$f_{mk}(p,a;x)=(1x^2)^{m/2}\mathrm{exp}[p(x1)]f(x).$$
(3.88)
So, the equation for $`f(x)`$ takes the form
$$(x^21)f^{\prime \prime }(x)+[2p(x^21)+2(m+1)x]f^{}(x)+[\lambda +m(m+1)+2p\sigma x]f(x)=0,$$
(3.89)
where we have denoted $`\sigma =\frac{a}{2p}(m+1)`$. In the case when the expansion
$$f(x)=\underset{s=0}{}a_su_s(x)$$
(3.90)
implies a three-terms recurrence relation, the eigenvalues $`\lambda _{mk}^{(x)}(p,a)`$ can be found from the chain fraction equation,
$$F^{(x)}(p,a;\lambda )=0.$$
(3.91)
Also, the expansion which is of practical use has been considered by Jaffe . In this case, the expansion series (3.90) becomes
$$f(x)=(x+1)^\sigma \underset{s=0}{}a_s\chi ^s,$$
(3.92)
where $`\chi =(x1)/(x+1)`$ is Jaffe’s variable. By inserting (3.92) into the equation for the function $`f(x)`$, we get recurrence relation (3.61), where the coefficients are
$$\alpha _s=(s+1)(s+m+1),$$
$$\beta _s=2s^2+(2s+m+1)(2p\sigma )am(m+1)+\lambda =$$
(3.93)
$$=2s(s+2p\sigma )(m+\sigma )(m+1)2p\sigma +\lambda ,$$
$$\gamma _s=(s1\sigma )(sm1\sigma ).$$
Also, for Jaffe series expansion, we have
$$\left|\frac{\alpha _{s1}\gamma _s}{\beta _{s1}\beta _s}\right|_s\mathrm{}=\frac{1}{4}\left(1\frac{4p}{s}\right)+O\left(\frac{p^2}{s^2}\right),$$
(3.94)
i.e., the chain fraction converges at $`p>0`$. One can see also that the Jaffe expansion converges at any $`x`$.
In addition, the function $`f(x)`$ can be expanded in associated Laguerre polynomials,
$$f(x)=(x+1)^\sigma \underset{s=0}{\overset{\mathrm{}}{}}a_sL_{s+m}^m(\overline{x}),\overline{x}=2p(x1).$$
(3.95)
In this case,the recurrence relation is of three-terms form, and the coefficients are
$$\alpha _s=(s+m+1)\left[\frac{a}{2p}(s+1)\right]=(s+m+1)(sm\sigma ),$$
$$\beta _s=(2s+m+1)\left[\frac{a}{2p}(s+m+1)\right]+$$
(3.96)
$$+2p(2s+m+1)(s+m)(m+1)a+\lambda ,$$
$$\gamma _s=s\left[\frac{a}{2p}(s+m)\right]=s(s1\sigma ).$$
As to numerical computation of the eigenvalues $`\lambda ^{(x)}`$, expansions (3.92) and (3.95) are equivalent because the chain fraction depends, in fact, only on $`\beta _s`$ and $`\alpha _s\gamma _{s+1}`$. Indeed, by comparing Eq.(3.93) and Eq.(3.96), one can easily see that in both cases $`\beta _s`$ and $`\alpha _s\gamma _{s+1}`$ are the same. Evidently, it then follows that the associated chain fractions are equivalent to each other.
However, we should note that Jaffe’s recurrence sequence, in general, is more stable, while Laguerre expansion (3.95) is more suitable to find out the asymptotics of $`f_{mk}(p,a;x)`$.
Also, we note that the associated ”radial” polynomials $`Q_{N+1}`$, the root $`\lambda ^{(x)}(p,a)`$ of which should be found, contain a much bigger number of terms, in comparison to the ”angular” case. So, practically finding of radial eigenvalues is much harder than that of angular eigenvalues.
At equal charges of nuclei, $`Z_1=Z_2`$, the equation for $`g`$, and the recurrence relation for $`g_k`$, are the same as they are in the particular case $`q=0`$ considered in Sec. 3.4. A general solution for $`g`$ is then given by acsf (3.38) and (3.39), with coefficients $`d_r^{mn}`$ given by recurrence relation (3.40).
In the reminder of this Section we would like to note that, in general, solving the recurrence relations can be made equivalent to solving associated ordinary differential equations by making the $`z`$ transform. In many cases the $`z`$ transform helps to solve recurrence relations. Namely, one defines the function
$$Z(z)=\underset{s=0}{\overset{\mathrm{}}{}}\frac{a_s}{z^n}$$
(3.97)
associated to the coefficients $`a_s`$ entering Eq.(3.61) viewed as a function of discrete variable $`s`$. For $`\alpha _s`$, $`\beta _s`$, and $`\gamma _s`$ given by Eq.(3.93) we obtain from Eq.(3.61)
$$z(z1)^2Z^{\prime \prime }+\left[(1m)z^2+2(2p\sigma 1)z+2\sigma +m+1\right]Z^{}+$$
(3.98)
$$+\left[(\sigma +m)(m+1)+2p\sigma \lambda +\frac{(m+\sigma )\sigma }{z}\right]Z=0.$$
For the coefficients $`c_s`$, we define
$$Y(z)=\underset{s=0}{\overset{\mathrm{}}{}}\frac{c_s}{z^n}.$$
(3.99)
and for $`\rho _s`$, $`\kappa _s`$, and $`\delta _s`$ given by Eq.(3.81) we obtain from Eq.(3.70)
$$z(z1)^2Y^{\prime \prime }+\left[(\tau 2)z^2+2(2p\tau +1)z+\tau \right]Y^{}+$$
(3.100)
$$+\left[(m(m+1)m\tau )z2p(m+1)+(\tau +m)(m+1)+b\lambda \frac{\tau }{z}\right]Z$$
$$m(m\tau +1)c_0z=0,$$
where we have denoted
$$\tau =\frac{b}{2p}+m+1.$$
(3.101)
If one has solved these differential equations for $`Z(z)`$ and $`Y(z)`$, then, by making the inverse $`z`$ transform, one can find the expansion coefficients $`a_s`$ and $`c_s`$ (and thus the general solution of the problem).
### 3.6 Asymptotics of csf and their eigenvalues
To analyze the exact solution, which is of rather complicated nonclosed form (infinite series) given in the previous Sections, it is much instructive to derive its asymptotics, which can be represented in a closed form. In this Section, we present the asymptotics at large ($`R\mathrm{}`$) and small ($`R0`$) distances between the nuclei, with a particular attention paid to the ground state.
#### 3.6.1 Asymptotics at $`R\mathrm{}`$
For increasing distances $`R`$ between the nuclei, at fixed quantum numbers $`k`$, $`q`$, and $`m`$, we have increasing values of the parameters $`p`$, $`a`$, and $`b`$,
$$p=(2E)^{1/2}R/2\mathrm{},a=(Z_2+Z_1)R\mathrm{},b=(Z_2Z_1)R\pm \mathrm{}.$$
(3.102)
Let us introduce the notation
$$\alpha =\frac{a}{2p}=\frac{Z_2+Z_1}{\sqrt{2E}},\beta =\frac{b}{2p}=\frac{Z_2Z_1}{\sqrt{2E}}$$
(3.103)
and assume that $`\alpha 1`$ and $`\beta 1`$.
acsf at $`R\mathrm{}`$.
Let us consider the asymptotic expansion of acsf. In this case, the equation for the Whittaker function, $`M_{\kappa ,\mu }(y)`$, builds ansatz around the poles $`y=\pm 1`$. Here, the solution is constructed in two overlapping intervals, $`𝒟_{}=[1,y_1]`$ and $`𝒟_+=[y_2,1]`$, with $`y_2<y_1`$. Then, the asymptotics of acsf $`g_{mq}(p,2p\beta ;y)`$ in the interval $`𝒟_{}`$ have the form
$$g_{mq}(p,2p\beta ;y)=\frac{d_{}}{\mathrm{\Gamma }(m+1)}\left[\frac{2\mathrm{\Gamma }\left(\kappa +\frac{1+m}{2}\right)}{\mathrm{\Gamma }\left(\kappa +\frac{1m}{2}\right)}\right]^{1/2}\times $$
$$\times \frac{M_{\kappa ,m/2}\left(2p(1+y)+2(\kappa +\beta )\mathrm{ln}\frac{1y}{2}\right)}{\sqrt{1y^2}}[1+O(p^1)],y𝒟_{}.$$
(3.104)
while in the interval $`𝒟_+`$ it is
$$g_{mq}(p,2p\beta ;y)=\frac{d_+}{\mathrm{\Gamma }(m+1)}\left[\frac{2\mathrm{\Gamma }\left(\kappa ^{}+\frac{1+m}{2}\right)}{\mathrm{\Gamma }\left(\kappa ^{}+\frac{1m}{2}\right)}\right]^{1/2}\times $$
$$\times \frac{M_{\kappa ^{},m/2}\left(2p(1y)+2(\kappa ^{}+\beta )\mathrm{ln}\frac{1+y}{2}\right)}{\sqrt{1y^2}}[1+O(p^1)],y𝒟_+.$$
(3.105)
Here, the coefficients $`d_{}`$ and $`d_+`$ ($`d_{}^2+d_+^2=1`$) are defined by the relations
$$d_{}=|\frac{\mathrm{sin}\pi (2\kappa ^{}m1)}{\mathrm{sin}\pi (2\kappa m1)+\mathrm{sin}\pi (2\kappa ^{}m1)}|^{1/2}\times $$
(3.106)
$$\times \mathrm{sgn}\left[\frac{\mathrm{cos}\pi (\kappa (m+1)/2)}{\mathrm{sin}\pi (\kappa ^{}(m+1)/2)}\right]$$
$$d_+=\left|\frac{\mathrm{sin}\pi (2\kappa m1)}{\mathrm{sin}\pi (2\kappa m1)+\mathrm{sin}\pi (2\kappa ^{}m1)}\right|^{1/2}.$$
(3.107)
rcsf at $`R\mathrm{}`$.
Now, let us consider the asymptotic expansion of rcsf. The replacements $`xy,`$ $`pp,`$ $`\alpha \beta `$ convert the radial equation around the point $`x=1`$ to the angular equation around the point $`y=1`$. Thus, the corresponding asymptotics of rcsf are directly related to the above found asymptotics of acsf.
The rcsf, normalized to the first order in $`p`$, has the form
$$f_{mk}(p,2p\alpha ;x)=\frac{1}{m!}\left[\frac{2(k++m)!}{k!(x^21)}\right]^{1/2}\times $$
$$\times M_{\kappa ,m/2}\left(2p(x1)\right)[1+O(p^1)].$$
(3.108)
Since the first index of Whittaker function in Eq.(3.108) is $`\kappa =k+(m+1)/2`$, the function can be expressed in terms of Laguerre polynomials.
Energy at $`R\mathrm{}`$.
In the limit $`R\mathrm{}`$, the Coloumb two-center problem is evidently reduced to two separate problems of Coloumb centers, with the charges $`Z_1`$ and $`Z_2`$. Each of the atoms, $`eZ_1`$ and $`eZ_2`$, is characterized by a set of parabolic quantum numbers, $`[n,n_1,n_2,m]`$ and $`[n^{},n_1^{},n_2^{},m]`$, which are related to each other by the relations
$$n=n_1+n_2+m+1,n^{}=n_1^{}+n_2^{}+m+1.$$
(3.109)
The number $`k`$ of zeroes of rcsf coincides with the number $`n_1`$, for the angular functions of the left center, $`eZ_1`$, and with the number $`n_1^{}`$, for the angular functions of the right center, $`eZ_2`$.
A series expansion in inverse power of $`R`$ can be obtained in the form
$$E_{[nn_1n_2m]}(Z_1,Z_2,R)=\frac{Z_1^2}{2n^2}\frac{Z_2}{R}+\frac{3Z_2n\mathrm{\Delta }}{2R^2Z_1}\frac{Z_2n^2}{2R^3Z_1^2}(6\mathrm{\Delta }^2n^2+1)+$$
(3.110)
$$+\frac{Z_2n^3}{16R^4Z_1^4}[Z_1\mathrm{\Delta }(109\mathrm{\Delta }^239n^29m^2+59)Z_2n(17n^23\mathrm{\Delta }^29m^2+19)]+$$
$$+\frac{\epsilon _5}{R^5}+\frac{\epsilon _6}{R^6}+O\left(\frac{1}{R^7}\right),$$
where $`\mathrm{\Delta }=n_1n_2`$, and $`\epsilon _{5,6}`$ are defined via the expressions,
$$\epsilon _5=\frac{n^3}{64Z_1^3}[n_Z(1065\mathrm{\Delta }^4594n^2\mathrm{\Delta }^2+1230\mathrm{\Delta }^2234m^2\mathrm{\Delta }^2+9m^4+$$
(3.111)
$$+33n^418n^2m^218m^2+105138n^2)+4n_Z^2\mathrm{\Delta }(21\mathrm{\Delta }^2111n^2+63m^2189)].$$
$$\epsilon _6=\frac{n^4}{64Z_1^4}[n_Z\mathrm{\Delta }(2727\mathrm{\Delta }^4+2076n^2\mathrm{\Delta }^25544\mathrm{\Delta }^2+1056m^2\mathrm{\Delta }^293m^4$$
(3.112)
$$273n^4+78n^2m^2+450m^21533+1470n^2)+2n_Z^2(207\mathrm{\Delta }^4+1044n^2\mathrm{\Delta }^2+$$
$$+2436\mathrm{\Delta }^2576\mathrm{\Delta }^2m^242n^2+371162m^2+42m^2n^289n^4+15m^4)+$$
$$+2n_Z^3\mathrm{\Delta }(3\mathrm{\Delta }^269n^211733m^2)],$$
where $`n_Z=nZ_2/Z_1`$.
Eq.(3.110) gives the multipole expansion in the electrostatic energy of the interaction between the atom $`eZ_1`$ and the far-distant charge $`Z_2`$ (so called $`eZ_1`$-terms).
Note that expansion (3.110) can be obtained by ordinary perturbation techniques as well. Indeed, the degrees of $`Z_1`$ display the orders of the multipole moment of the atom $`eZ_1`$.
The series of terms corresponding to the other atom, $`eZ_2`$, is obtained from Eq.(3.110) with the use of self-evident replacements, $`Z_1Z_2`$, $`nn^{}`$, $`\mathrm{\Delta }\mathrm{\Delta }^{}`$, and $`n_2n_2^{}`$.
Finally, the energy of the ground state $`1s\sigma _g`$ of the molecular ion, for which $`Z_1=Z_2=1`$ (equal charges of nuclei), can be written, to a high accuracy, as
$$E_{1000}(1,1,R)=\frac{1}{2}\frac{9}{4R^4}\frac{15}{2R^6}\frac{213}{4R^7}\frac{7755}{64R^8}\frac{1733}{2R^9}\frac{86049}{16R^{10}}O\left(\frac{1}{R^{11}}\right).$$
(3.113)
#### 3.6.2 Asymptotics at $`R0`$
Energy at $`R0`$.
In the case of positive total charge, $`Z=Z_1+Z_2>0`$, and at $`R0`$, we can use perturbative approach to $`Z_1eZ_2`$ problem, without using a separation of variables. Namely, the Hamiltonian of the system $`Z_1eZ_2`$ is represented as the sum
$$\widehat{H}=\widehat{H}^{UA}+\widehat{W}=\frac{\widehat{P}^2}{2m}\frac{Z_1}{r_1}\frac{Z_2}{r_2}.$$
(3.114)
The operator $`\widehat{H}^{UA}`$ is usually chosen as the Hamilton operator of the so called united atom,
$$\widehat{H}^{UA}=\frac{\widehat{P}^2}{2m}\frac{Z}{r_c},$$
(3.115)
which is placed on the $`z`$-axis at the point $`z=z_0`$,
$$z_0=\left(\frac{1}{2}+\frac{Z_2}{Z}\right)R=\left(\frac{1}{2}+\frac{Z_1}{Z}\right)R.$$
(3.116)
The point $`(0,0,z_0)`$ is called center of charges due to the fact that it lies at the distances
$$R_1=\frac{Z_2}{Z}R\text{ and }R_2=\frac{Z_1}{Z}R,$$
(3.117)
from the left and right atoms, respectively.
We choose a spherical coordinate system, $`(r_c,\vartheta _c,\phi )`$, with the origin at point $`(0,0,z_0)`$, and the angle $`\vartheta _c`$ measured from $`z`$-axis. Then, the eigenstates $`\psi _{Nlm}^{UA}`$ of the operator $`\widehat{H}^{UA}`$ are
$$\psi _{Nlm}^{UA}(\stackrel{}{r}_c)=R_{Nl}(r_c)Y_l^m(\vartheta _c,\phi ),$$
(3.118)
while the eigenvalues are given by
$$E_{Nlm}^{UA}=\frac{Z^2}{2N^2}.$$
(3.119)
The matrix $`W_{Nlm}^{Nl^{}m^{}}`$ of the perturbation operator $`\widehat{W}`$ is diagonal on the functions $`\psi _{Nlm}^{UA}(\stackrel{}{r}_c)`$ of the atom if $`z_0`$ is defined by Eq.(3.116). Below, the first two terms of the expansion of energy in powers of $`R`$ are given,
$$E_{Nlm}(Z_1,Z_2,R)=$$
$$\frac{Z^2}{2N^2}\frac{2Z_1Z_2[l(l+1)3m^2]}{N^3l(l+1)(2l1)(2l+1)(2l+3)}(ZR)^2+O((ZR)^3).$$
(3.120)
For the ground state of the $`Z_1eZ_1`$ system, with equal values of the charges, one can find the following expression for the energy, up to the second order of perturbation:
$$E_{000}^{(2)}(Z_1,Z_1,R)=$$
$$Z^2\left[\frac{1}{2}+\frac{1}{6}(ZR)^2\frac{1}{6}(ZR)^3+\frac{43}{2160}(ZR)^4\frac{1}{36}(ZR)^5\mathrm{ln}ZR+\mathrm{}\right].$$
(3.121)
csf at $`R0`$.
For fixed quantum numbers we have, at $`R0`$,
$$p=(2E)^{1/2}R/20,a=(Z_2+Z_1)R0,b=(Z_2Z_1)R0.$$
(3.122)
Let us denote
$$\alpha =\frac{a}{2p}=\frac{Z_2+Z_1}{\sqrt{2E}}=\sigma +m+1,\beta =\frac{b}{2p}=\frac{Z_2Z_1}{\sqrt{2E}}.$$
(3.123)
In this notation, the energy is $`E=Z^2/(2\alpha )^2`$. Let us consider asymptotics of csf of the ground state of the molecular ion.
acsf at $`R0`$.
The power series expansion of acsf $`g_{00}(p,2p\beta ;y)`$ in small parameter $`p`$ can be obtained by expanding it in the Legendre polynomials. For the eigenvalue $`\lambda _{00}^{(y)}(p,2p\beta )`$, we then get
$$\lambda _{00}^{(y)}(p,2p\beta )=(1\beta ^2)\left[\frac{2}{3}p^2\frac{2}{135}p^4(1+11\beta ^2)+O(p^6)\right].$$
(3.124)
rcsf at $`R0`$.
To expand rcsf $`f_{00}(p,2p(1+\sigma );x)`$, at $`p0`$, $`\sigma =O(p^2)`$, we use Jaffe’s expansion,
$$f_{00}(p,2p(1+\sigma );x)=\mathrm{exp}(px)(1+x)^\sigma \underset{s=0}{\overset{\mathrm{}}{}}a_s\chi ^s,\chi =\frac{x1}{x+1},$$
(3.125)
where $`a_s`$’s obey three-term recurrence relation with the coefficients (3.93).
For the eigenvalue $`\lambda _{00}^{(x)}(p,2p(1+\sigma ))`$, we get
$$\lambda _{00}^{(x)}(p,2p(1+\sigma ))=\sigma (1+2p)+\sigma ^2(1+4p\mathrm{ln}4p\gamma )+o(p^5).$$
(3.126)
rcsf of the ground state of the molecular ion, $`Z_1=Z_2=1`$, can be presented as
$$f_{00}(p,2p(1+\sigma );x)=\mathrm{exp}(px)(1+x)^\sigma [1+\sigma ^2\mathrm{Li}_2(\chi )+o(p^4)],$$
(3.127)
where $`\chi =(x1)/(x+1)`$ is Jaffe’s variable and $`\mathrm{Li}_2(\chi )`$ is dilogarithm function,
$$\mathrm{Li}_2(\chi )=\underset{n=1}{\overset{\mathrm{}}{}}\frac{\chi ^n}{n^2}=\underset{0}{\overset{\chi }{}}\mathrm{ln}(1\xi )\frac{d\xi }{\xi }.$$
(3.128)
The ground state energy is defined as a function of three parameters, $`Z_1`$, $`Z_2`$, and $`R`$,
$$\lambda _{00}^{(y)}(p,2p\beta )=\lambda _{00}^{(x)}(p,2p(1+\sigma ))$$
(3.129)
Combining Eqs.(3.124), (3.126) and (3.129), we get series expansion for the ground state energy of the $`Z_1eZ_2`$ system in the form
$$E_{000}(Z_1,Z_2,R)=$$
$$\frac{1}{2}Z^2+\frac{2}{3}Z_1Z_2(ZR)^2\frac{2}{3}Z_1Z_2(ZR)^3+\frac{2}{5}Z_1Z_2\left(1\frac{64Z_1Z_2}{27Z^2}\right)(ZR)^4$$
$$\frac{8}{45}Z_1Z_2\left[\frac{5Z_1Z_2}{Z^2}\mathrm{ln}(2RZ\gamma )+1\frac{199Z_1Z_2}{12Z^2}\right](ZR)^5+o((ZR)^5).$$
(3.130)
Comparing Eqs.(3.121) and (3.130) we see that the terms proportional to $`(ZR)^2`$ and $`(ZR)^3`$ coincide. The next order corrections makes a difference; Eq.(3.121) obtained by the second-order perturbation is of less accuracy. A practically achieved accuracy of the first-order perturbation (3.120) and of Eq.(3.130) is the same; at $`ZR<0.1`$, the discrepancy is not bigger than 1%, and becomes sharply smaller with the increase of the parameter $`ZR`$.
#### 3.6.3 Quasiclassical asymptotics
At $`R\mathrm{}`$, for $`eZ_1`$ solutions we have
$$\lambda _{mk}^{(x)}(p,2p\alpha )=2p(2\kappa \alpha )\kappa (2\kappa \alpha m)+$$
(3.131)
$$+\frac{\kappa }{2p}(2\kappa ^23\kappa \alpha +\alpha ^2m^2)+o(p^2),$$
$$\lambda _{mq}^{(y)}(p,2p\beta )=2p(2\chi +\beta )\chi (2\chi +2\beta m)$$
(3.132)
$$\frac{\chi }{2p}(2\chi ^2+3\chi \beta +\beta ^2m^2)+o(p^2).$$
From the equality $`\lambda _{mk}^{(x)}=\lambda _{mq}^{(y)}`$, we get the expansion for $`E_j(R)`$ which coincides with the asymptotics (3.110), up to terms of the order of $`R^2`$.
In the limit $`R0`$, the following expansions are justified,
$$\lambda _{mk}^{(x)}=\left[\frac{a}{2p}(k+1/2)\right]^2+O(p^2),$$
(3.133)
$$\lambda _{mq}^{(y)}=(l+1/2)^2+\frac{p^2}{2}\frac{b^2}{8(l+1/2)^2}+O(p^4).$$
(3.134)
and we get, by using the equation $`\lambda ^{(x)}=\lambda ^{(y)}`$,
$$E_{Nlm}(R)=\frac{1}{2}\left(\frac{Z_1+Z_2}{N}\right)^2R^2\frac{Z_1Z_2(Z_1+Z_2)^2}{4N^3(l+1/2)^5}[(l+1/2)^23m^2].$$
(3.135)
Expression (3.135) for the energy coincides with the asymptotics (3.120), up to $`O(l^2)`$. Clearly, an accuracy of the quasiclassical Eqs.(3.131)-(3.135) becomes higher for a greater number of zeroes of the solutions, $`k`$ and $`q`$. However, even for the ground state, $`k=q=m=0`$, these equations give a good approximation for the energy in both limiting cases, $`R0`$ and $`R\mathrm{}`$. Also, we note that corresponding numerical calculations showed that for the intermediate values of $`R`$ the terms $`E_j(R)`$ can be determined within the quasiclassical approach with accuracy of about 5%, or more .
## 4 Scaling method and binding energy of three-body Santilli-Shillady isochemical model $`\widehat{H}_2`$
To find the ground state energies of $`H_2^+`$ and $`\widehat{H}_2`$, we use computations of the $`1s\sigma `$ terms of $`H_2^+`$ ion and of $`\widehat{H}_2`$ based on the above presented exact csf solution by solving the corresponding equations $`\lambda ^{(x)}=\lambda ^{(y)}`$.
The angular and radial eigenvalues $`\lambda `$ are found as solutions of the equations containing infinite chain fractions presented in previous Sections which should be interrupted and then solved numerically, to a required accuracy.
Our primary interest is the study of Santilli-Shillady model $`\widehat{H}_2`$ system. However, we present the results for $`H_2^+`$ ion as well to check our calculations and to use them in the scaling method described below.
Also, the reader should keep in mind that we are primarily interested in ascertaining whether there exist a non-zero value of $`R`$ for which a fully stable and point-like isoelectronium permits an exact representation of the binding energy of $`H_2`$ molecule.
It should be noted that the ground state electronic energy is obtained as a function of the parameter $`R`$ due to Eq.(3.58). By adding to it the internuclear potential energy $`1/R`$, we obtain the total ground state energy of the system, so that at some value $`R=R_{min}`$, the total energy $`E`$ necessarily has a minimum, if the system is stable. This is the way to determine uniquely the internuclear distance under an exact representation of the total energy $`E`$.
To have an independent check of the result for the total ground state energy of $`\widehat{H}_2`$ (with the stable and point-like isoelectronium) obtaining from the exact solution, we develop a scaling method based on the original Schrödinger equation for $`H_2^+`$ ion like system. Namely, it appears that one is able to calculate the ground state energy as a function of $`R`$ for any $`H_2^+`$ ion like system with equal charges of nuclei, $`Z_1=Z_2`$, provided that one knows the ground state energy as a function of $`R`$ for the $`H_2^+`$ ion itself. It should be pointed out that the scaling method does not depend on the obtained solution because it reflects, in fact, the scaling properties of the Schrödinger equation itself.
In addition, we use below Ritz’s variational approach to $`H_2^+`$ like systems to find out the approximate value of the ground state energy of $`\widehat{H}_2`$, as well as to check the result provided by the exact solution, and to demonstrate the accuracy of the variational approximation.
Our general remark is that in both approaches, the exact solution and Ritz’s variational solution, we use Born-Oppenheimer approximation (fixed nuclei). Clearly, taking into account the first order correction, i.e., zero harmonic oscillations of the nuclei in $`H_2`$ around their equilibrium positions, we achieve greater accuracy.
But we still have a significant inaccuracy in the value of dissociation energy due to the fact that $`H_2`$ system has the lightest possible nuclei (two single protons).
To estimate this inaccuracy, one can invoke Morse’s potential customarily used for diatomic molecules. In particular, the analysis for $`H_2`$ molecule shows that the ground state energy of harmonic oscillations of the nuclei receives 1.4% correction due to the first anharmonic term.
### 4.1 Exact representation of binding energies of $`H_2^+`$ ion and three-body $`\widehat{H}_2`$ system
The csf based computations for $`H_2^+`$ ion were presented, for example, by Teller , Bates et al. , and Wind , and we do not repeat this study here for brevity, while we shall just describe the procedure and present our final numerical results in Appendix, Table 2.
In particular, Teller presented a plot of the resulting function $`E_{1s\sigma }(R)`$ which provides a good accuracy, and Wind used the exact solution to present a table of energy values in seven decimal places for distance values of $`R`$ up to 20 a.u. in steps of 0.05 a.u. However, these results cannot be used directly in our case since the repulsive potential between the nuclei, $`1/R`$, has not been accounted for, and, as the main reason, we have the isoelectronium instead of one single electron.
In the Appendix, we present the results obtained from the csf based recurrence relations by numerical calculations for ordinary $`H_2^+`$ ion and for $`\widehat{H}_2`$ system, at $`M=2m_e`$. These results are presented in Tables 2 and 4. Tables 3 and 5 have been derived from Tables 2 and 4, respectively, by simple adding internuclear potential energy $`1/R`$, to obtain the total energy of the system.
In Table 2, we present the $`1s\sigma _g`$ electronic term of $`H_2^+`$. In Table 3, we present the total energy of $`H_2^+`$. In Table 4, we present the $`1s\sigma _g`$ term of $`\widehat{H}_2`$, at the mass $`M=2m_e`$. In Table 5, we present the total energy of $`\widehat{H}_2`$, at the mass $`M=2m_e`$. Also, in Table 6, we present the total minimal energies of $`\widehat{H}_2`$ and optimal distances $`R`$, for various values of the isoelectronium mass parameter, $`M=\eta m_e`$. All the data of these Tables are purely theoretical and, additionally, we plot them in Figures 1–8, for the reader convenience. Figures 6 and 8 give more detailed view on the interval $`0.26<\eta <0.34`$.
The analysis of the data in Tables 25 is simple. Namely, one should identify the minimal value of the energy in each Table. One can use Figures 1–4 for visual identification of the minima, and then turn to the corresponding Tables 25, to reach a higher numerical accuracy.
Let us consider, as an example, Table 2. One can see that the energy minimum for $`H_2^+`$ is
$$E_{1s\sigma }=2.0\text{ a.u.}\text{ at }R=0\text{ a.u.}$$
(4.1)
Note that $`E_{1s\sigma }`$ is the electronic energy, i.e. the internuclear repulsion has not been taken into account here. Our remark is that this energy value corresponds to $`He^+`$ ion due to the fact that two $`H_2^+`$ nuclei are superimposed and form $`He`$ nucleous at $`R=0`$.
Let us consider now Table 3. In this Table, one can find the line $`|2.0|0.602634|`$ corresponding to visual minimal value of the energy. To identify a more precise value of the minimal energy, one should use the interpolation of the data. This gives us the minimum of the total energy of $`H_2^+`$,
$$E=0.6026346\text{ a.u.}\text{ at }R=1.9971579\text{ a.u.}$$
(4.2)
This theoretical value represents rather accurately the known experimental value $`E_{exper}[H_2^+]=0.6017\text{ a.u.}`$ for $`H_2^+`$ ion, thus establishing the validity of our csf based calculations. For completeness, we note that the experimental dissociation energy of $`H_2^+`$ ion is $`D_{exper}[H_2^+]0.0974\text{ a.u.}=2.65\text{ eV}`$, and the internuclear distance $`R_{exper}[H_2^+]2.00\text{ a.u.}=1.0584\AA `$.
Let us now consider Table 4. One can see that the energy minimum for $`\widehat{H}_2`$, at $`M=2m_e`$, is
$$E_{1s\sigma }=16.0\text{ a.u.}\text{ at }R=0\text{ a.u.}$$
(4.3)
Note that $`E_{1s\sigma }`$ only yields the isoelectronium’s energy, i.e. the internuclear repulsion has not been taken into account here. Our remark is that this energy value corresponds to the $`He`$ atom, where the two electrons form a stable point-like isoelectronium of mass $`M=2m_e`$.
Let us consider now Table 5 which is of striking interest for our study. In this Table, one can find the line $`|0.250|7.61428940411169996|`$ corresponding to the visual minimal value of the energy. To identify a more precise value of the minimal energy, one should use the interpolation of the data given in this Table. This gives us the minimum of the total energy of $`\widehat{H}_2`$,
$$E=7.617041\text{ a.u.}\text{ at }R=0.258399\text{ a.u.}$$
(4.4)
This theoretical value is in quite good agreement with the preceding theoretical result by Santilli and Shillady obtained via structurally different variational numerical method, $`E_{var}=7.61509174`$ at $`R_{var}=0.2592`$ (see third column of Table 1 in Ref.). It is quite naturally to observe that this variational energy is a bit higher (by $`0.002\text{ a.u.}`$) than the above one obtained from the exact solution (as it is expected to be for any variational solution).
However, this exact theoretical value (4.4) does not meet the experimental value $`E_{exper}[H_2]=1.17\text{ a.u.}`$ known for $`H_2`$ molecule. Indeed, adopted approximation that the isoelectronium is point-like, stable, and has mass $`M=2m_e`$ leads us to the theoretical value (4.4) while the known experimental value, $`E_{exper}[H_2]=1.17\text{ a.u.}`$, differs much from it.
Essentially the same conclusion is due to numerical program SASLOBE by Santilli and Shillady , where Gaussian screened Coloumb potential interaction between the electrons, rather than the stable point-like isoelectronium approximation, has been used to achieve final precise fit of $`E=1.174474\text{ a.u.}`$, with the obtained bond length $`R=1.4011\text{ a.u.}`$, at the isoelectronium correlation length $`r_c=0.01125\text{ a.u.}`$ (see Table 1 in Ref.). We discuss on this issue in Sec. 5.
Our remark is that, due to Table 5, the experimental value $`E=1.17\text{ a.u.}`$ is fitted by the distance $`R=0.072370\text{ a.u.}`$ However, this energy value is not minimal and thus can not be ascribed reasonable physical treatment in Table 5.
Our conclusion from the above analysis is that we have two main possibilities to overcome this sharp discrepancy between our theoretical and the experimental binding energy values which has place at $`M=2m_e`$:
1. Consider unstable isoelectronium, i.e. the four-body Santilli-Shillady model of $`H_2`$ molecule;
2. Treat the mass $`M`$ of isoelectronium as a free parameter, instead of fixing it to $`M=2m_e`$, assuming thus some defect of mass discussed in Introduction,
in order to fit the experimental data on $`H_2`$ molecule.
The first possibility will be considered in a subsequent paper because it needs in application of different technique, while the second possibility can be studied within the three-body Santilli-Shillady model under consideration to which we turn below.
In the next Section, we develop simple formalism allowing one to deal with the mass and charge of isoelectronium viewed as free parameters, and arrive at the conclusion (see Table 6) that the restricted three-body Santilli-Shillady model of $`H_2`$ molecule is capable to fit the experimental binding energy, with the total mass of isoelectronium equal to $`M=0.308381m_e`$, although with the internuclear distance about 19.6% bigger than the experimental value.
### 4.2 The scaling method
In order to relate the characteristics of $`H_2^+`$ ion like system to that of thoroughly studied $`H_2^+`$ ion, we develop scaling method based on the Schrödinger equation. The neutral $`\widehat{H}_2`$ system with stable point-like isoelectronium is an example of $`H_2^+`$ ion like system in which we are particularly interested here. Below, we develop scaling method for the case of arbitrary mass and charge of the particle.
Let us write the Schrödinger equation for a particle of the rescaled charge
$$e\zeta e,$$
(4.5)
(we turn here from $`e=1`$ to $`e=1`$ representation), and the rescaled mass
$$m_e\eta m_e,$$
(4.6)
with equal charges of nuclei, $`+eZ_1=+eZ_2=+eZ`$,
$$\left[\frac{\mathrm{}^2}{2\eta m_e}_r^2\frac{\zeta Ze^2}{r_a}\frac{\zeta Ze^2}{r_b}+\frac{Z^2e^2}{R_{ab}}\right]\psi =E\psi .$$
(4.7)
where $`\eta `$ and $`\zeta `$ are scaling parameters, and $`R_{ab}`$ is distance between the nuclei. The condition $`Z_1=Z_2`$ is an essential point to stress here because owing to which we can successfully develop the scaling method. We introduce the unit of length,
$$r_0=\frac{1}{\eta \zeta Z}r_B\frac{1}{\eta \zeta Z}\frac{\mathrm{}^2}{m_ee^2},$$
(4.8)
where $`r_B`$ is Bohr’s radius. Dividing Eq.(4.7) by $`\zeta Ze^2`$, and multiplying it by $`r_0`$, we get
$$\left[\frac{\mathrm{}^2}{\eta \zeta Zm_ee^2}r_0\frac{1}{2}_r^2r_0\frac{1}{r_a}r_0\frac{1}{r_b}+r_0\left(\frac{Z}{\zeta }\right)\frac{1}{R_{ab}}\right]\psi =\frac{r_0E}{\zeta Ze^2}\psi .$$
(4.9)
We introduce dimensionless entities $`\rho =r/r_0`$, $`\rho _a=r_a/r_0`$, $`\rho _b=r_b/r_0`$, and $`=R_{ab}/r_0`$. Then, Laplacian in Eq.(4.9) becomes $`r_0^2_r^2=_\rho ^2`$. Further, introducing unit of energy,
$$E_0=\frac{\eta m_e\zeta ^2Z^2e^4}{\mathrm{}^2}\eta \zeta ^2Z^2\frac{m_ee^4}{\mathrm{}^2},$$
(4.10)
we have dimensionless energy $`\epsilon =E/E_0`$ so that Eq.(4.9) can be rewritten as
$$\left[\frac{1}{2}_\rho ^2\frac{1}{\rho _a}\frac{1}{\rho _b}+\frac{1}{(\frac{\zeta }{Z})}\right]\psi =\epsilon \psi .$$
(4.11)
Note that, at $`\eta =1`$, $`\zeta =1`$, and $`Z=1`$, the constants $`r_0(\eta ,\zeta ,Z)`$ and $`E_0(\eta ,\zeta ,Z)`$ reproduce ordinary atomic units,
$$r_0(1,1,1)=r_B=\frac{\mathrm{}^2}{m_ee^2},E_0(1,1,1)=2E_B=\frac{m_ee^4}{\mathrm{}^2},$$
(4.12)
and we recover the case of $`H_2^+`$ ion. On the other hand, in terms of dimensionless entities the original Schrödinger equation for $`H_2^+`$ ion is
$$\left[\frac{1}{2}_\rho ^2\frac{1}{\rho _a}\frac{1}{\rho _b}+\frac{1}{R}\right]\psi _0=\epsilon (R)\psi _0,$$
(4.13)
where $`R=R_{ab}/r_B`$. Comparison of Eq.(4.11) and Eq.(4.13) shows that by putting $`R=(\zeta /Z)`$ in Eq.(4.11), we obtain the equation,
$$\left[\frac{1}{2}_\rho ^2\frac{1}{\rho _a}\frac{1}{\rho _b}+\frac{1}{R}\right]\psi =\epsilon (R)\psi ,$$
(4.14)
which identically coincides with the original Eq.(4.13). The difference is that Eq.(4.14) is treated in terms of the rescaled units, $`r_0(\eta ,\zeta ,Z)`$ and $`E_0(\eta ,\zeta ,Z)`$, instead of the ordinary Bohr’s units, $`r_B`$ and $`E_B`$. As the result, we have one and the same form of Schrödinger equation for any $`H_2^+`$ like system characterized by equal charges of nuclei. This makes a general ground to calculate some characteristic entity of any $`H_2^+`$ like system when one knows its value for $`H_2^+`$ ion.
Particularly, one can easily derive $`R_{ab}`$ and $`E`$ for the system with arbitrary parameters $`\eta `$, $`\zeta `$, and $`Z`$ from their values, $`R_{ab}[H_2^+]`$ and $`E[H_2^+]=2E_B\epsilon (R)`$, obtained for $`H_2^+`$ ion (for which $`\eta =1`$, $`\zeta =1`$, and $`Z=1`$). Indeed, since for arbitrary $`\eta `$, $`\zeta `$, and $`Z`$
$$R=\frac{\zeta }{Z}=\frac{\zeta }{Z}\frac{R_{ab}}{r_0}=\frac{\zeta }{Z}\frac{R_{ab}}{r_B}\eta \zeta Z,$$
(4.15)
we can establish the following relationship between the distances corresponding to arbitrary $`Z\zeta Z`$ system and $`H_2^+`$ ion,
$$R_{ab}=\frac{R[H_2^+]}{\eta \zeta ^2}.$$
(4.16)
It is remarkable to note that the dependence on $`Z`$ disappeared in Eq.(4.16). In the case of isoelectronium of mass $`M=2m_e`$ and charge $`2e`$, we have $`\eta =2`$ and $`\zeta =2`$, so that
$$R_{ab}=\frac{R[H_2^+]}{8}.$$
(4.17)
Also, the energy $`E(R)`$ of $`Z\zeta Z`$ system and energy $`\epsilon (R)`$ of $`H_2^+`$ ion are related to each other according to the equation,
$$E(R)=\eta \zeta ^2Z^2\left(\frac{m_ee^4}{\mathrm{}^2}\right)\epsilon (R).$$
(4.18)
#### 4.2.1 The case $`M=2m_e`$
So, in the case of isoelectronium of mass $`M=2m_e`$ and charge $`2e`$, we get
$$E(R_{ab})=8\epsilon (R).$$
(4.19)
Note however that the factor $`\zeta /Z=2`$ arised due to $`R=(\zeta /Z)`$ is hidden here so that in order to calculate the values of $`E(R_{ab})`$ and $`R_{ab}`$ from $`\epsilon (R)`$ and $`R`$ respectively one should multiply $`\epsilon `$ by 8 and $`R`$ by $`1/4`$.
As the result, in accordance with the scaling method the points can be calculated due to the following rule:
$$(R,E)(R,E+1/R)(R/4,8E)(R/4,8E+4/R),$$
(4.20)
for Tables 2 $``$ 3 $``$ 4 $``$ 5. One can easily check numerically that these properties indeed hold true for the presented Tables. Thus, the scaling method can be used instead of the independent numerical calculations for $`\widehat{H}_2`$ system if one has the data for $`H_2^+`$ ion.
It is highly important to note here that the energy minimum in Table 3 is not rescaled to the energy minimum in Table 5 due to the absence of energy scaling between these Tables; see Eq.(4.20), from which one can observe that $`(8E+4/R)`$ can not be expressed as $`n(E+1/R)`$, where $`n`$ is a number. So one needs to identify minimum in Table 5 independently (after calculating all the points), rather than direct rescale the minimum from Table 3 to try to get minimum for Table 5.
#### 4.2.2 The case $`M=\eta m_e`$
For a more general case of isoelectronium mass, $`M=\eta m_e`$, and charge $`2e`$, we should keep the following sequence of calculations:
$$(R,E)(R,E+1/R)(\frac{R}{2\eta },4\eta E)(\frac{R}{2\eta },4\eta E+\frac{2\eta }{R}).$$
(4.21)
starting from Table 2 to obtain, at the last step, the table of values (similar to Table 5) from which we should extract a minimal value of the energy and the corresponding optimal distance, at each given value of mass $`\eta `$. The result of the analysis of a big number of such tables is collected in Table 6, where the interval $`0.26<\eta <0.34`$ appears to be of interest; $`M=\eta `$, in atomic units. Plots of the data of Table 6 are presented in Figures 5 and 7 (Figures 6 and 8 give more detailed view on the interval of interest) show that
$$E_{min}(M)3.808M,R_{opt}(M)\frac{0.517}{M},$$
(4.22)
to a good accuracy. Note that $`E_{min}(M)`$ unboundedly decreases with the increase of $`M`$ (there is no local minimum), so we can use a fit, instead of the minimization in respect with $`M`$. From this Table, we obtain the following final fit of the binding energy for the restricted three-body Santilli-Shillady model of $`H_2`$ molecule:
$$M=0.308381m_e,E=1.174475\text{ a.u.},R=1.675828\text{ a.u.},$$
(4.23)
where the mass parameter $`M`$ of the isoelectronium has been varied in order to meet the experimental energy $`E_{exper}[H_2]=1.174474\text{ a.u.}=31.9598\text{ eV}`$. Using this value of mass, $`M=0.308381m_e`$, we computed the total energy as a function of the internuclear distance $`R`$, and depicted it in Figure 9 to illustrate that $`R=1.675828\text{ a.u.}`$ indeed corresponds to a minimal value of the energy. Note that the predicted optimal distance $`R=1.675828\text{ a.u.}=0.886810\AA `$ appears to be about 19.6% bigger than the conventional experimental value $`R_{exper}[H_2]=1.4011\text{ a.u.}=0.742\AA `$.
This rather big (19.6%) discrepancy can not be ascribed to the Born-Oppenheimer approximation used in this paper since it gives relatively small uncertainty in the energy value, even in the case of $`H_2`$ molecule. We stress here that in the Born-Oppenheimer approximation, the three-body problem (the Schrödinger equation) can be given exact solution owing to separation of the electronic and nuclear degrees of freedom while the full three-body problem (accounting for the wave functions of the nuclei, etc.) can not be solved exactly.
In a strict consideration, we should calculate the dissociation energy of $`H_2`$ molecule, $`D=2E_0EE^{nucl}`$, to make comparison to the experimental value, $`D_{exper}[H_2]0.164\text{ a.u.}=4.45\text{ eV}`$ . Here, $`E_0=0.5\text{ a.u.}=13.606\text{ eV}`$ is the ground state energy of separate $`H`$-atom and $`E^{nucl}`$ is the energy of zero mode harmonic oscillations of the nuclei, with the experimental value $`E_{exper}^{nucl}[H_2]`$ $`0.01\text{ a.u.}`$ = $`0.27\text{ eV}`$ . One can see that the zero mode energy $`E^{nucl}`$ (which is taken to be $`E^{nucl}=0`$, in the Born-Oppenheimer approximation) is estimated to be less than 1% of the predicted $`E`$. The leading anharmonic correction to the harmonic oscillation energy is estimated to be 1.4% of $`E^{nucl}`$, i.e. it is of the order of $`0.00014\text{ a.u.}=0.004\text{ eV}`$, in the case of $`H_2`$ molecule. So, in total the Born-Oppenheimer approximation makes only up to 1% uncertainty, which is obviously insufficient to treat the predicted $`R=1.675828\text{ a.u.}`$ as an acceptable value, from the experimental point of view.
Note that, at the given $`M`$, we can not ”fix” $`R`$ to be equal to the desired experimental value $`R_{exper}=1.4011\text{ a.u.}`$ unless we shift $`E`$ to some nonminimal value, which is, as such, meaningless. Conversely, if we would fit experimental $`R_{exper}`$ by varying $`M`$, we were obtain $`E_{min}`$ between $`1.52\text{ a.u.}`$ and $`1.33\text{ a.u.}`$ (see Table 6), which is much deviated from the experimental $`E_{exper}`$. In other words, the relation between $`E`$ and $`R`$, governed by the Schrödinger equation, is such that at some value of $`R`$ there is a minimum of $`E`$ so that $`R`$ is not some kind of free parameter here since the system tends to minimize its own energy. In accordance to the exact solution of the model, our single free parameter, $`M`$, can not provide us with the exact fit of both the experimental values, $`E_{exper}`$ and $`R_{exper}`$.
Thus, we arrive at the conclusion that the three-body Santilli-Shillady model of $`H_2`$ molecule yields the result (4.23), which indicates that the assumption of stable point-like isoelectronium builds a crude approximation to the general (four-body) Santilli-Shillady model. This means that we are forced to possess that the isoelectronium is not stable point-like quasi-particle, to meet the experimental data on $`H_2`$ molecule.
### 4.3 Variational solution
In studying $`H_2^+`$ ion like systems, one can use Ritz variational approach to obtain the value of the ground state energy as well. This approach assumes analytical calculations, which are easier than that used in finding the above exact solution but they give approximate value of the energy. It is helpful in making simplified analysis of the system. This can be made for the general case of isoelectronium total mass, which eventually undergoes some ”defect” while its ”bare” total mass is assumed to be $`M=2m_e`$. Ritz variational solution of the $`H_2^+`$ like problem yields, of course, similar result for the energy of $`\widehat{H}_2`$. Below, we present shortly results of our calculations. However, we stress that the variational solution is given here just to make some support to the exact solution, and to see the typical order of the variational approximation.
Using hydrogen ground state wave function and one-parameter Ritz variation, we obtain the following expression for the energy of $`H_2^+`$ like system:
$$E(\rho )=\frac{1}{2}\frac{e^2}{a_0}+\frac{e^2}{a_0}\frac{1}{\rho }\frac{(1+\rho )e^{2\rho }+(1\frac{2}{3}\rho ^2)e^\rho }{1+(1+\rho +\frac{1}{3}\rho ^2)e^\rho },$$
(4.24)
where $`\rho =R/a_0`$ is variational parameter. For the general case of mass $`m`$ and charge $`q=\zeta e`$, Eq.(4.24) can be rewritten in the following form:
$$E(\rho ,\zeta )=\frac{Me^4\zeta ^2}{\mathrm{}^2}\left(\frac{1}{2}+F(\rho )\right),$$
(4.25)
where
$$F(\rho )=\frac{1}{\rho }\frac{(1+\rho )e^{2\rho }+(1\frac{2}{3}\rho ^2)e^\rho }{1+(1+\rho +\frac{1}{3}\rho ^2)e^\rho }.$$
(4.26)
Numerically, the function $`F(\rho )`$ reaches minimum at the value $`\rho =2.5`$, which should be used in the above expression for $`E(\rho ,\zeta )`$. So, putting $`\zeta =1`$ we obtain the variational value of $`H_2^+`$ ion energy, $`E(\rho )=0.565`$. Note, to make a comparison, that we have the value $`E_{exact}=0.6026`$ due to the exact solution (4.2), and the value $`E_{exper}[H_2^+]=0.6017`$ as the experimental value of the energy of $`H_2^+`$ ion. Thus, the optimal distance between the protons in $`H_2^+`$ ion is $`R_m=a_0\rho =2.5\text{ a.u.}`$, and the obtained variational energy $`E`$ is slightly higher than both the values $`E_{exact}`$ and $`E_{exper}`$, as it is normally expected to be in the variational approach. Now, we should replace electron by isoelectronium to describe the associated $`\widehat{H}_2`$ model. Substituting $`M=2m_e`$ and $`\zeta =2`$, we see that the r.h.s. of Eq.(4.25) contains overall factor 8, in comparison to the $`H_2^+`$ ion case ($`M=m_e`$ and $`\zeta =1`$),
$$\stackrel{~}{E}(\rho )=8|2E_B|\left(\frac{1}{2}+F(\rho )\right),$$
(4.27)
The function $`F(\rho )`$ remains the same, and its minimum is reached again at $`\rho =2.5`$. Then, energy of $`H_2`$ molecule due to Eq.(4.27) is $`\stackrel{~}{E}(\rho )=8|2E_B|0.565=4.520\text{ a.u.}`$ This value should be compared with the one given by Eq.(4.4).
Below, we collect the above mentioned data and results of this Section in Table 1.
## 5 Concluding remarks
In this paper we have shown that the restricted three-body Santilli-Shillady isochemical model of the hydrogen molecule admits an exact analytic solution capable of representing the molecular binding energy in a way accurate to the sixth digit, $`E=1.174475\text{ a.u.}`$, and the internuclear distance $`R=1.675828\text{ a.u.}`$, which is about 19.6% bigger than the conventional experimental value, $`R_{exper}[H_2]=1.4011\text{ a.u.}`$
We should emphasize that the presented exact analytical solution includes infinite chain fractions. They still need numerical computation to reach the characteristic values of $`H_2^+`$ ion like systems, such as the ground state energy, with the understanding that these values can be reached with any needed accuracy. For example, at the lengths of the chain fractions $`N=100`$ and $`N=50`$ for the angular and radial eigenvalues, one achieves accuracy of the ground state energy of about $`10^{12}`$.
The general (four-body) Santilli-Shillady isochemical model of $`H_2`$ cannot be, apparently, solved exactly, even in Born-Oppenheimer approximation, so that Ritz variational approach can be applied here to get the approximate values of the ground state energy and corresponding internuclear distance.
Ritz variational approach is a good instrument to analyze few-body problems, and restricted $`H_2`$ molecule is such a system. It is wellknown that the variational solution of the ordinary model of $`H_2`$ molecule includes rather complicated analytical calculations of the molecular integrals, with the hardest part of work being related to the exchange integral. Particularly, evaluated exchange integral for $`H_2`$ molecule is expressed in terms of a special function (Sugiura’s result, 1927). It is quite natural to expect that even more complications will arise when dealing with the Hulten potential.
The reason to consider the general four-body Santilli-Shillady model of $`H_2`$ molecule, after the made analysis of $`H_2^+`$ like system approximate approach to it, is that the stable point-like isoelectronium $`\widehat{H}_2`$ model-based theoretical prediction does not meet the experimental data on $`H_2`$ molecule, for the ”bare” isoelectronium mass $`M=2m_e`$, although we achieved essentially exact representation of the binding energy taking $`M=0.308381m_e`$. Also, this stable point-like isoelectronium (three-body) model does not account for essential effect existing in the general (four-body) model. This effect is related to the potential barrier between the region associated to the attractive Hulten potential, $`r_{12}<r_0`$, and the region associated to repulsive Coloumb potential, $`r_{12}>r_0`$, where $`r_0`$ is the distance between the electrons at which Hulten potential is equal to Coloumb potential, $`V(r_{12})=0`$; see Eq.(2.5). Characteristics of the barrier can be extracted from the function $`V(r_{12})`$. The barrier is finite for the used values of the parameters $`V_0`$ and $`r_c`$ so that the electrons penetrate it. The two $`1s`$ electrons are thus simultaneously in two regimes, the first is strongly correlated regime due to short-range attractive Hulten potential (isoelectronium) and the second is weakly correlated regime due to the ordinary Coloumb repulsion. Also, there exist a transient regime corresponding to the region about the equilibrium point, $`r_{12}r_0`$, i.e. inside the barrier. Schematically, one could thought of that the electrons are, for instance, 10% in the isoelectronium regime, 1% in the transient regime, and 89% in the Coloumb regime. We stress that in the three-body approach to $`H_2`$ molecule considered in this paper we have 100% for the isoelectronium regime.
Numerical computation by Santilli and Shillady based on Gaussian transform techniques and SASLOBE computer program has shown excellent agreement of the general four-body model with experimental data on $`H_2`$ molecule. They used Gaussian screened Coloumb potential as an approximation to the Hulten potential. It would be instructive to use Ritz variational approach, which deals with analytical calculations, in studying the four-body Santilli-Shillady model of $`H_2`$ molecule. One can try it first for the Gaussian screened Coloumb potential, or exponential screened Coloumb potential in which case there is a hope to achieve exact analytical evaluation of the Coloumb and exchange integrals. Being a different approach, this would give a strong support to the numerical results on the ground state energy obtained by Santilli and Shillady. Also, having analytical set up one can make qualitative analysis of the four-body Santilli-Shillady model of $`H_2`$ molecule. However, we should to note that these potentials, being approximations to the Hulten potential, will yield some approximate models, with corresponding approximate character of the results.
## Appendix
We use $`N=16`$ power degree approximation, the polynomials $`Q_N^{(x)}`$ and $`Q_N^{(y)}`$, to find both the radial, $`\lambda ^{(x)}(p,a)`$, and angular, $`\lambda ^{(y)}(p,b)`$, eigenvalues of the csf. $`Q_N`$’s are obtained by the use of recurrence relations (3.76) and definitions of the coefficients $`\alpha _s`$, $`\beta _s`$, $`\gamma _s`$, $`\rho _s`$, $`\kappa _s`$, and $`\delta _s`$, where we put $`b=0`$, i.e. $`Z_1=Z_2=1`$, and quantum number $`m=0`$. Each of the two polynomials has 16 roots for $`\lambda `$ from which we select one root which is appropriate due to its asymptotic behavior at $`R0`$. Numerical solution of the equation $`\lambda ^{(x)}(p,a)=\lambda ^{(y)}(p,b)`$ gives us the list of values of the electronic ground state energy $`E(R)=E_{1s\sigma }(R)`$, which corresponds to $`1s\sigma _g`$ term of the $`H_2^+`$ ion, as a function of the distance $`R`$ between the nuclei. Table 2 presents the result, where no interpolation has been used. Numerical computation of each point in Table 2 took about 88 sec on ordinary Pentium desktop computer.
Below, we present some useful numerical values enabling one to convert atomic units, at which $`m_e=e=\mathrm{}=1`$, to the other units. Note also that for the energy $`1\text{ a.u.}1\text{ hartree}`$, and for the length $`1\text{ a.u.}1\text{ bohr}`$.
Atomic units in terms of the other units
| 1 a.u. of mass, $`m_e`$ | 9.10953$`10^{28}`$ | gramms |
| --- | --- | --- |
| 1 a.u. of charge, $`e`$ | 1.60219$`10^{19}`$ | Coloumbs |
| 1 a.u. of action, $`\mathrm{}`$ | 1.05459$`10^{27}`$ | $`\mathrm{erg}\mathrm{sec}`$ |
| 1 a.u. of length, $`\frac{\mathrm{}^2}{m_ee^2}`$ | 0.529177$`10^8`$ | cm |
| 1 a.u. of energy, $`\frac{m_ee^4}{\mathrm{}^2}`$ | 27.2116 | eV |
| 1 a.u. of time, $`\frac{\mathrm{}^3}{m_ee^4}`$ | 2.41888$`10^{17}`$ | sec |
| 1 a.u. of velocity, $`e^2/\mathrm{}`$ | 2.18769$`10^8`$ | cm/sec |
| $`\alpha =\frac{e^2}{\mathrm{}c}`$ | 1/137.0388 | |
Conversion of the energy units
| | a.u. | eV | Kcal$``$mole | cm<sup>-1</sup> |
| --- | --- | --- | --- | --- |
| a.u. | 1 | 27.212 | 6.2651$`10^2`$ | 2.1947$`10^5`$ |
| eV | 3.6749$`10^2`$ | 1 | 23.061 | 8065.48 |
| Kcal$``$mole | 1.5936$`10^3`$ | 4.3364$`\dot{1}0^2`$ | 1 | 3.4999$`10^2`$ |
| cm<sup>-1</sup> | 4.5563 $`10^6`$ | 1.2398$`\dot{1}0^4`$ | 2.8573$`10^3`$ | 1 |
Minimum of the energy $`E_{1s\sigma }(R)`$ is $`E_{1s\sigma }=1.9999999976\text{ a.u.}`$ at $`R=0`$, which reproduces the known value $`E_{1s\sigma }=2\text{ a.u.}`$ to a very high accuracy. Moreover, one can compare Table 2 and the table of Ref. to see that each energy value in Table 2 does reproduce Wind’s result up to seven decimal places. This means that our numerical calculations are correct.
We remark that Wind used $`N=50`$ approximation and presented seven decimal places while we use $`N=16`$ approximation and present seventeen decimal places. Alas, there is no need to keep such a high accuracy, and also Wind mentioned that even $`N=10`$ approximation gives the same result, up to seven digits.
By making 16th-order interpolation of the points in Table 2 and adding to it the potential of interaction between the nuclei, $`1/R`$, we obtain the list of values of the total energy presented in Table 3. It reveals the only minimum of the total energy, $`E(R)+R^1=E_{min}=0.6026346\text{ a.u.}`$ at the distance $`R=R_{opt}=1.9971579\text{ a.u.}`$
The results collected in Table 4 have been obtained directly by numerical calculations with the use of replacements $`p^22p^2`$ and $`a4a`$, where $`p`$ and $`a`$ are defined by Eq.(3.50), in the coefficients $`\alpha _s`$, $`\beta _s`$, $`\gamma _s`$, $`\rho _s`$, $`\kappa _s`$, and $`\delta _s`$ of the recurrence relations. These replacements have been made due to Eq.(3.15), with the mass parameter $`M=2`$ and the charge parameter $`q=2`$, corresponding to the stable point-like isoelectronium of mass $`M=2m_e`$ and charge $`2e`$. In addition, it turns out that Table 4 can be derived directly from Table 2 by the use of rescalements $`RR/4`$ and $`E8E`$. This remarkable property is confirmed by the scaling method developed in Sec. 4.2, and proves that the scaling method is correct. By adding $`1/R`$ to the isoelectronic energy values of Table 4 we obtain Table 5 showing the total energy of the $`\widehat{H}_2`$ system, at the mass $`M=2m_e`$. The minimum of the total energy is found $`E(R)+R^1=E_{min}=7.617041\text{ a.u.}`$ at $`R=R_{opt}=0.258399\text{ a.u.}`$
Table 6 presents result of calculations of the minimal total energies and corresponding optimal distances, at various values of the isoelectronium mass parameter $`M=\eta m_e`$ ($`M=\eta `$, in atomic units). We have derived some 27 tables (such as Table 5) from Table 2 by the scaling method according to Eq.(4.21), and find minimum of the total energy in each table, together with the corresponding optimal distance. Then we collected all the obtained energy minima and optimal distances in Table 6. With the fourth order interpolation/extrapolation, the graphical representations of Table 6 show (see Figures 5–8) that the minimal total energy behaves as $`E_{min}(M)3.808M`$, and the optimal distance behaves as $`R_{opt}(M)0.517/M`$, to a good accuracy. One can see that at $`M=2m_e`$ we have $`E_{min}(M)=7.617040\text{ a.u.}`$ and $`R_{opt}(M)=0.258396\text{ a.u.}`$, which recover the earlier obtained values $`E_{min}=7.617041\text{ a.u.}`$ and $`R_{opt}=0.258399\text{ a.u.}`$ of Table 5, to a high accuracy, thus showing once again correctness of the used scaling method. In fact, the values of $`E`$ and $`R`$ for $`M=1.50m_e`$, $`M=1.75m_e`$, and $`M=2.00m_e`$ in Table 6 have been obtained by extrapolation so they are not as much accurate as they are in Table 5. However, this is not of much importance here because we use them only to check the results of the scaling method.
The main conclusion following from Table 6 is that the mass parameter value $`M=0.308381m_e`$ fits the energy value $`E_{min}(M)=1.174475\text{ a.u.}`$, with the corresponding $`R_{opt}(M)=1.675828\text{ a.u.}`$, which appears to be about 19.6% bigger than the experimental value $`R_{exper}[H_2]=1.4011\text{ a.u.}`$ The total energy as a function of internuclear distance, for this value of mass, $`M=0.308381m_e`$, is shown in Figure 9 to illustrate that the obtained optimal distance $`R_{opt}=1.675828\text{ a.u.}`$ corresponds to a minimal value of the total energy. |
warning/0001/astro-ph0001040.html | ar5iv | text | # Updated Information on the Local Group
## 1 Introduction
The study of individual nearby galaxies that belong to the Local Group is presently one of the most active areas of extragalactic research. As a result the recent reviews by van den Bergh (1999a, 2000) are already slightly out of date. The purpose of the present paper is to provide an updated report on the status of the most recent research on individual Local Group members, and on the structure of the Local Group itself.
## 2 Information on Individual Galaxies
The individual galaxies for which new information has become available are listed below in approximate order of decreasing luminosity.
### 2.1 The Andromeda galaxy (M31 = NGC 224)
Kormendy & Bender (1999) have used spectroscopy with high angular resolution to study the binary nucleus of M31. Their observations support Tremaine’s model in which P<sub>1</sub> is the brightest part of a single eccentric disk, where stars linger while at the apocenters of their orbits around P<sub>2</sub>. The latter object is found to contain a $`3.3\times 10^7\mathrm{M}_{\mathrm{}}`$ black hole. Barmby et al. (1999) have published a catalog of 435 probable globular clusters in M31, of which 330 have UBV photometry, and 158 have been observed spectroscopically. These authors find that the metal-poor clusters have larger projected distances from the galaxy center, and show lower rotation, than do the metal-rich clusters. They estimate that M31 contains $`450\pm 100`$ globulars, from which the specific globular cluster frequency $`\mathrm{S}=0.9\pm 0.2`$. Hamilton & Fesen (1999) have used the Hubble Space Telescope to image absorption by the remnant of the supernova S Andromedae (SN 1885A), which is observed to have a diameter of $`0^{\prime \prime }.55\pm 0^{\prime \prime }.15`$, in the light of Fe II. From their observations these authors conclude that this remnant contains between 0.1 and 1.0 $`\mathrm{M}_{\mathrm{}}`$ of iron.
### 2.2 The Milky Way system
Many years ago Eggen & Sandage (1959) suggested that RR Lyrae and Groombridge 1830 (HR 4550) belong to a small physical group of five stars that are traversing the solar neighborhood at high velocity. The view that such clumps of high-velocity stars are real is supported by observations with the HIPPARCOS satellite (Helmi et al. 1999). From space motions of a nearly complete sample of nearby high-velocity halo stars these authors conclude that seven objects are members of a single debris stream. From their data Helmi et al. conclude that $``$8% of all metal-poor stars outside the solar radius represent the remnants of a single disintegrated dwarf galaxy. However, this estimate may turn out to be too high if a single narrow debris stream presently happens to be passing through the solar neighborhood.
Figer et al. (1999) have determined the slope of the mass spectrum of in the Arches and Quintuplet clusters, which are located near the Galactic center. They find that the mass spectra of these very young clusters have a slope $`\mathrm{\Gamma }=0.65`$, which is less steep than that for young clusters elsewhere in the Galaxy which typically have $`\mathrm{\Gamma }1.4`$. Each cluster has a mass of $`1\times 10^4\mathrm{M}_{\mathrm{}}`$, which are among the highest known in the Galaxy. Taken at face value these results suggest that the region near the Galactic center is particularly prone to forming very massive open clusters. From near-infrared echelle spectra Carr, Sellgren & Balachandran (1999) find that \[Fe/H\] = -0.02 $`\pm `$ 0.13, i.e. nearly solar, for the Galactic center supergiant IRS 7.
### 2.3 The Triangulum galaxy (M33 = NGC 598)
Corbelli & Salucci (1999) have measured the rotation curve of M33 out to a distance of 16 kpc (13 disk scale-lengths) and find that the rotation curve rises out to the last measured point. This result implies a dark halo mass of $`5\times 10^{10}\mathrm{M}_{\mathrm{}}`$. Beyond 3 kpc the gravitational potential is dominated by a dark halo with a density that decreases radially as R<sup>-1.3</sup>. From an unbiased sample of 60 clusters Chandar et al. (1999a,b) find that cluster formation in M 33 has been continuous over the last 10 Gyr, i.e. unlike the LMC, M 33 did not have a gap in its cluster formation history. Young clusters in M 33 have masses in the range $`6\times 10^2\mathrm{M}_{\mathrm{}}`$ to $`2\times 10^4\mathrm{M}_{\mathrm{}}`$, which is smaller than those of the old clusters which typically have masses of a few $`\times 10^5\mathrm{M}_{\mathrm{}}`$. Gordon et al. (1999) have doubled the sample of supernova remnants that are known in the Triangulum galaxy. Among their sample of 53 SNRs they find no evidence for a strong correlation between surface brightness and diameter. Many of these remnants are found to be associated with, or embedded in, H II regions. This (not unexpectedly) suggests that the majority of these remnants were produced by supernovae of Type II.
### 2.4 The Large Magellanic Cloud (LMC)
Detached eclipsing variables are powerful tools for the determination of extragalactic distances. Recently Nelson et al. (1999) have redetermined the reddening of the eclipsing variable HV 2274 and find E(B-V) $`=0.083\pm 0.006`$, which is significantly lower than previously published values. From their new reddening Nelson et al. derive a distance modulus (m-M)$`{}_{0}{}^{}=18.40\pm 0.07`$. Gibson (1999) has reviewed recent distance determinations for the LMC, which range from (m-M)$`{}_{0}{}^{}=18.20`$ to 18.75. This large spread shows that significant unappreciated sources of systematic error still exist in modern determinations of the distance to the Clouds of Magellan.
Sakai, Zaritsky & Kennicutt (1999) have recently used the magnitude level of the tip of the LMC red giant branch to derive a distance modulus (m-M)$`{}_{0}{}^{}=18.59\pm 0.09`$ (random) $`\pm 0.16`$ (systematic), which agrees well with previous determinations via Cepheid variables.
Dolphin (1999b) has studied the star formation history in two fields in the LMC. His results, which are summarized in Table 1, appear to show (1) a steady rise of the metallicity index \[Fe/H\] with time, and (2) that the rate of star formation between 2.5 and 7 Gyr ago was an order of magnitude lower than it has been during the most recent 2–3 Gyr period. Since the two fields studied by Dolphin are separated by $`3^{}.0`$ (2.6 kpc) his data refer to global star formation rates. Taken at face value Dolphin’s results appear to weaken the previous conclusion that the rate of cluster formation in the Large Cloud increased more rapidly $``$3 Gyr ago than did the rate of star formation. However, Holtzman et al. (1999) reach a very different conclusion from Hubble Space Telescope observations of two fields in the outer disk of the LMC. They conclude that there was no gap in the age distribution of these stars. Clearly we are still far from understanding the evolutionary history of the disk component of the LMC. From new HST color- magnitude diagrams Johnson et al. (1999) find that the age difference between the Large Cloud clusters NGC 1466, NGC 2257 and Hodge 11, on the one hand, and the Galactic clusters M 92 and M 3 on the other, is less than 1.5 Gyr. It is presently not understood why globular cluster formation occurred simultaneously in the LMC, the main body of the Galactic halo, and in the outer Galactic halo (NGC 2419).
From the relatively large number of red clump stars in the LMC Bar Holtzman et al. conclude that the stellar population in the Bar is older than that in the outer fields of the Large Cloud. This conflicts with some previous work that had suggested that the LMC Bar might be a relatively recent feature. From their integrated spectra Dutra et al. (1999) conclude that NGC 1928 (\[Fe/H\] $`=1.2`$) and NGC 1939 (\[Fe/H\] $`=2.0`$) may be globular clusters that had not previously been recognized as such. If confirmed this would increase the number of LMC globular clusters from 13 to 15. A catalog of X-ray sources in the LMC has been published by Haberl & Pietsch (1999). These authors give likely identifications for 144 sources. Of these objects 46 (32%) appear to be associated with supernova remnants, 17 (12%) are X-ray binaries, and nine (6%) are “supersoft” sources. The majority of unidentified sources are probably associated with background galaxies or foreground stars. Demers & Battinelli (1999) have surveyed the periphery of the LMC for young blue stars that might be associated with the Bridge linking the Large and the Small Clouds. Few such stars are found suggesting that the Bridge does not extend deep into the LMC. Liebert (1999) has pointed out that the hot blue star found in the LMC cluster NGC 1818 has the wrong luminosity and radius to be a luminous white dwarf.
Murali (1999) finds that the motion of the Magellanic Stream through ambient gas can strongly heat the Stream clouds, driving mass loss and causing evaporation. Survival of the stream for 500 Myr sets an upper limit $`<10^5\mathrm{cm}^3`$ for the Galactic halo gas through which the Stream in orbiting.
### 2.5 The Small Magellanic Cloud
Most previous investigators have concluded that cluster formation in the Small Cloud has proceeded at a more-or-less uniform rate. However new cluster age determinations by Rich et al. (1999) now suggest that cluster formation in the SMC may have been enhanced during bursts that occurred 2 Gyr and 8 Gyr ago. More, and more accurate, age determinations for Small Cloud clusters will be required to strengthen this conclusion. An unbiased survey of 93 star clusters in a 2.4 square degree area of the SMC (Pieterzyński & Udalski 1999) shows an age distribution that is very strongly biased towards young clusters, with only 3% of the clusters having ages $`>1`$ Gyr. The fact that 60% of all SMC clusters are younger than 100 Myr should probably be interpreted as evidence for short cluster life-times, rather than as evidence for a recent burst of cluster formation.
Bica & Dutra (1999) have published an updated census of SMC clusters, and of clusters in the Bridge between the LMC and SMC, that is based on the OGLE survey. Their paper contains a map of the distribution of SMC clusters, which shows a strong concentration in the Small Cloud “Bar” and a lesser concentration of clusters in the Wing of the SMC.
Cold atomic hydrogen has been detected in the Bridge between the Magellanic Clouds by Kobulnicky & Dickey (1999). The early-type stars observed in the Bridge could have formed in these cold clouds. These objects therefore need not have migrated from the main body of the SMC.
Rolleston et al. (1999) have compared the abundances of three early-type stars in the Bridge between the LMC and the SMC with those observed in “normal” B-type stars near the Sun. They found an average metal deficiency \[m/H\] $`=1.05\pm 0.12`$ for these objects. Surprisingly this value falls well below the present (Luck et al. 1998) SMC metallicity \[Fe/H\] $`=0.74`$. In fact it lies between that of the globular cluster NGC 121 (\[Fe/H\] $`=1.19`$), which has an age of $``$12 Gyr, and that of the old open cluster L 1 (\[Fe/H\] $`=1.01`$) that has an age of $``$10 Gyr. Taken at face value this result suggests that the material in the Bridge was tidally detached from the Small Cloud $``$10 Gyr ago. However, Putman (1999) has argued that the Bridge was formed only 0.2 Gyr ago, when the LMC and SMC had a very close encounter with a minimum separation of only $``$7 kpc. Alternatively it might be hypothesized that the metallicity of gas in the Bridge was lowered by swept-up pristine inter-galactic gas. The metallicity of NGC 330, which is the most luminous young SMC cluster, remains controversial. Hill (1999) finds $`[\mathrm{Fe}/\mathrm{H}]=0.82\pm 0.11`$ for six cool cluster supergiants, compared to $`[\mathrm{Fe}/\mathrm{H}]=0.69\pm 0.10`$ for six cool SMC field stars. Finally Gonzalez & Wallerstein (1999) obtain $`[\mathrm{Fe}/\mathrm{H}]=0.94\pm 0.02`$ for seven stars in NGC 330.
The ninth nova to be discovered in the SMC was found near the optical center of this galaxy by Glicenstein (1999) at $`\alpha =0^\mathrm{h}59^\mathrm{m}23.^\mathrm{s}0,\delta =73^{}07^{}56^{\prime \prime }`$ (equinox 2000).
### 2.6 The spheroidal<sup>1</sup><sup>1</sup>1The galaxies NGC 147, NGC 185, and NGC 205 belong to the same morphological family as do the dwarf spheroidals, such as Sculptor and Fornax. However, their relatively high luminosity makes the term dwarf spheroidal seem inappropriate. I therefore call such objects “spheroidals”. It seems likely that such spheroidals are lower luminosity examples of the form family that de Vaucouleurs (1959) refers to as “lenticulars”. galaxy NGC 205
A nova has been discovered in NGC 205 by Johnson & Modjaz (1999).
### 2.7 The starburst galaxy IC 10
From observations in the far infrared Bolatto et al. (1999) find that dust in the mild starburst galaxy IC 10 appears to deficient in small grains. It seems likely that such grains were destroyed by intense UV radiation in the neighborhood of the hot luminous stars that were formed during the recent burst of star formation in this galaxy.
### 2.8 The spheroidal galaxy NGC 185
The history of star formation in NGC 185 has been studied by Martínez-Delgado, Aparicio & Gallart (1999), who find that the bulk of the star formation in this galaxy took place at early times. Stars only formed near the center of this galaxy during the last $``$1 Gyr. Most of the young blue objects discovered by Baade (1951) turn out to be star clusters, rather than individual stars. Martínez-Delgado et al. have also discovered a supernova remnant near the center of NGC 185.
### 2.9 The dwarf irregular IC 1613
Hubble Space Telescope observations by Cole et al. (1999) show that the dominant stellar population in IC 1613 has an age of $``$7 Gyr. From its Hess diagram these authors find that the evolutionary history of IC 1613 may have been similar to that of the Pegasus dwarf (DDO 216). Both of these objects are classified Ir V on the DDO system. Antonello et al. (1999) have found five Cepheids of Population II in IC 1613, thus providing prima facie evidence for the existence of a very old stellar population component in this galaxy. King, Modjaz & Li (1999) have observed what is believed to be the first nova ever observed in IC 1613.
### 2.10 The Wolf-Lundmark-Melotte system (DDO 221)
A recent color-magnitude diagram obtained with the Hubble Space Telescope (Dolphin 1999a) shows that about half of all star formation in the WLM galaxy occurred during a burst that began $``$13 Gyr ago. During the course of this burst the metallicity increased from \[Fe/H\] $``$$`2.2`$ to \[Fe/H\] $`=1.3`$. From the apparent absence of a horizontal branch population Dolphin places an upper limit of $``$$`20\mathrm{M}_{\mathrm{}}`$ per Myr on the rate of star formation between 12 Gyr and 15 Gyr ago. Between 2.5 and 9 Gyr ago the average rate of star formation was $`100200\mathrm{M}_{\mathrm{}}`$ per Myr.
The WLM system contains a single globular cluster, for which Hodge et al. (1999) find M$`{}_{\mathrm{V}}{}^{}=8.8`$, \[Fe/H\] $`=1.5`$ and an age of $``$15 Gyr. It is of interest to note that the lone globular in this dwarf galaxy has a luminosity that falls slightly above that of the mean for all globular clusters. The apparent absence of faint low-mass globulars is of particular interest because destruction of such objects probably can not be attributed to the weak tidal forces of the WLM dwarf galaxy.
### 2.11 The disintegrating Sagittarius galaxy
From observations of tidal debris of the Sagittarius system Johnston et al. (1999) conclude that this object has orbited the Galaxy for at least 1 Gyr, and that the mass of this galaxy has decreased by a factor of 2–3 during this period. The orbits of Sgr are found to have Galactocentric distances that oscillate between $``$13 kpc and $``$41 kpc and periods in the range 550–750 Myr. Its most recent perigalacticon occurred 50 Myr ago. \[Jiang & Binney (1999) find that the Sagittarius dwarf might have started its infall from a distance greater than 200 kpc if its initial mass was as great as $``$10$`{}_{}{}^{11}\mathrm{M}_{\mathrm{}}^{}`$.\] Both the orbit of Sgr, and the Galactic potential field, could be constrained by improved proper motion observations of the stellar debris associated with this object. Burton & Lockman (1999) have found no neutral hydrogen gas associated with the Sagittarius dwarf.
### 2.12 The Fornax dwarf spheroidal
Buonanno et al. (1999) have recently used the Hubble Space Telescope to obtain a color-magnitude diagram for the globular cluster Fornax No. 4. Whereas the clusters Fornax 1, 2, 3 and 5 have horizontal branches that extend over a wide range of colors (and include RR Lyrae variables), Fornax 4 is seen to have a red horizontal branch. Fornax 4 is $``$3 Gyr younger than the other Fornax globulars. Buonanno et al. draw attention to the fact that the color-magnitude diagram of Fornax 4 exhibits a strong resemblance to that of the “young” Galactic globular cluster Ruprecht 106 (Fusi Pecci et al. 1995, and references therein). This observation raises two questions: (1) What is it about dwarf spheroidals that allows them to form globular clusters such as Fornax 4 and Terzan 7 (which is associated with the Sagittarius dwarf) long after the formation of globulars ceased in the main body of the Galactic halo, and (2) could it be that “young” outer halo globulars such as Ruprecht 106 were originally formed in now defunct dwarf spheroidals?
### 2.13 The Sagittarius dwarf irregular galaxy (SagDIG)
On the basis of a very uncertain distance of 1.4 Mpc SagDIG (UKS 1927-177 $`=`$ UGA 438) had previously been regarded as a possible Local Group member. New photometry by Lee & Kim (1999) yields a distance (based on the tip of the red giant branch) of only $`1.18\pm 0.10`$ Mpc, and an improved integrated luminosity M$`{}_{\mathrm{V}}{}^{}=11.97`$, which makes this object slightly more luminous that Leo A (M$`{}_{\mathrm{V}}{}^{}=11.5`$). A similar distance of $`1.06\pm 0.10`$ Mpc has recently been found by Karachentsev, Aparicio & Makarova (1999). If, following Courteau & van den Bergh (1999), we assume that the barycenter of the Local Group is situated at a distance of 454 kpc towards $`\mathrm{}=121.^{}7,\mathrm{b}=21.^{}3`$, then SagDIG is located at a distance of $`1.29\pm 0.09`$ Mpc from the LG barycenter. This makes SagDIG the most remote object suspected of Local Group membership. Its distance is marginally greater than the radius of the zero-velocity surface of the Local Group, for which Courteau & van den Bergh (1999) find a value of R$`{}_{0}{}^{}=1.18\pm 0.15`$ Mpc.
### 2.14 The Leo I (Regulus system)
From HST observations Gallart et al. (1999) conclude that 70% – 80% of the star forming activity in Leo I took place between 7 Gyr and 1 Gyr ago. There is little or no evidence for the presence of stars with ages $`>10`$ Gyr. About 1 Gyr ago the rate of star formation appears to have dropped abruptly to a near-negligible level. However, some very low-level star formation may have continued until $``$300 Myr ago.
### 2.15 The M31 companion And II
From spectra of seven stars in And II, that were obtained with the Keck telescope, Côté et al. (1999) have found a mean velocity V$`{}_{\mathrm{r}}{}^{}=188\pm 3\mathrm{km}\mathrm{s}^1`$, and a velocity dispersion of $`9.3\pm 2.6\mathrm{km}\mathrm{s}^1`$. From these data they obtain a mass-to-light ratio M/L$`{}_{\mathrm{V}}{}^{}=21_{10}^{+14}`$ in solar units, i.e. this dwarf spheroidal appears to contain a significant amount of dark matter. Côté, Oke & Cohen (1999) have also obtained Keck spectra of 42 red giants in And II, from which they find a mean metallicity $`[\mathrm{Fe}/\mathrm{H}]=1.47\pm 0.19`$, with a dispersion of $`0.35\pm 0.10`$ dex. Da Costa et al. (1999) have studied the color-magnitude diagram of And II and find that the majority of stars have ages in the range 6 Gyr to 9 Gyr, although the presence of RR Lyrae variables and blue horizontal branch stars attests to the existence of a population component with an age $`>10`$ Gyr. And II differs from And I in that it does not exhibit a radial gradient in horizontal branch morphology. Furthermore the dispersion in abundance is considerably larger in And II than it was found to be in And I. These results show that these two dwarf spheroidal companions to M 31 must have had quite different evolutionary histories. It would be interesting to know if there is a correlation between a radial horizontal branch gradient and the metallicity dispersion among dwarf spheroidal galaxies. On the basis of its 680 kpc distance Da Costa et al. conclude that And II is physically associated with M 31, rather than with M 33.
### 2.16 The dwarf spheroidals And V, And VI and And VII
Caldwell (1999) has derived accurate luminosities and surface brightness profiles for Andromeda V, VI and VII. And V turns out to be fainter than previously believed, whereas And VI and And VII were found to be more luminous than previously thought. And V has a metallicity that lies above the average metallicity versus luminosity relation for Local Group dwarf galaxies.
### 2.17 The Aquarius dwarf (AqrDIG = DDO 210)
The membership of this dwarf irregular in the Local Group has now been firmly established by Lee et al. (1999), who derive a distance of $`950\pm 50`$ kpc from the magnitude of the giants at the tip of the red giant branch. The corresponding distance of this object from the barycenter of the Local Group is also 950 kpc. This implies that DDO 210 is rather isolated in space.
### 2.18 The recently discovered Cetus system
Whiting, Hau & Irwin (1999) have searched the region with $`\delta <+3^{}`$ for faint dwarfs that might be members of the Local Group. Two objects were discovered during this program. One of these was the Antlia system, which lies just beyond the zero-velocity surface of the Local Group (van den Bergh 1999b). The other was a faint dwarf spheroidal galaxy in Cetus. From the position of the tip of the giant branch of the Cetus system Whiting et al. derive a Galactocentric distance of $`775\pm 50`$ kpc, and a distance of 615 kpc from the adopted center of the Local Group. This places Cetus comfortably inside the 1.18 Mpc radius (Courteau & van den Bergh 1999) of the Local Group zero-velocity surface. It would be of great importance to obtain radial velocities for individual red giants in the Cetus dwarf. This would enable one to determine the amount of dark matter in this galaxy. Furthermore knowledge of the systemic velocity of Cetus would add weight to the Local Group mass determination. No neutral hydrogen gas has been found to be associated with the Cetus dwarf.
### 2.19 The Sculptor dwarf spheroidal
Majewski et al. (1999) find that the metallicity distribution in Sculptor appears to be bimodal with components having \[Fe/H\] $`=2.3`$ and \[Fe/H\] $`=1.5`$. As is the case in many other Local Group galaxies the older metal-poor component appears more extended than the younger metal-rich component. However, Hurley-Keller, Mateo & Grebel (1999), although confirming the central concentration of red horizontal branch stars, do not find a radial age (or metallicity) gradient.
### 2.20 The Phoenix dwarf galaxy
Martínez-Delgado, Gallart & Aparicio (1999) find that this dIr/dSph galaxy has an inner component that contains young stars which is stretched in the east-west direction, and an outer component that is extended north-south, which is mainly populated by old stars. The rate of star formation in the central region of this galaxy appears to have remained approximately constant over time. St.-Germain et al. (1999) have found a cloud of $``$10$`{}_{}{}^{5}\mathrm{M}_{\mathrm{}}^{}`$ of H I located just west of Phoenix, that has a radial velocity of $`23\mathrm{km}\mathrm{s}^1`$, which may be physically associated with this object. Optical radial velocities of stars in Phoenix will be required to confirm this suspicion.
### 2.21 The Ursa Minor dwarf
New Hubble Space Telescope observations by Mighell & Burke (1999) confirm that the UMi system had a simple evolutionary history with a single $``$2 Gyr long burst of star formation occurring $``$14 Gyr ago.
## 3 DISCUSSION
New data by Armandroff et al. (1999), Caldwell (1999), Da Costa et al. (1999), Dolphin (1999a), Lee & Kim (1990), Lee et al. (1999), and Whiting et al. (1999) are incorporated in Table 2, which lists classification types on the DDO system, luminosities, Galactocentric distances, and distances from the adopted barycenter, for all 36 presently known Local Group members. Figure 1 shows a plot of the integral frequency of galaxy distances from the Local Group barycenter adopted by Courteau & van den Bergh (1999). To minimize the influence of observational selection effects galaxies with M$`{}_{\mathrm{V}}{}^{}>10.0`$ have not been plotted. The figure shows that only one galaxy (SagDIG at R$`{}_{\mathrm{LG}}{}^{}=1.29\pm 0.09`$ Mpc) lies marginally beyond the dynamically determined radius of the Local Group zero-velocity surface, for which Courteau & van den Bergh (1999) determined R$`{}_{0}{}^{}=1.10\pm 0.15`$ Mpc. From Eqn. (10) of Courteau & van den Bergh (1999) it is seen that the age of the Local Group can be determined from its mass M, and the radius of its zero-velocity surface. Adopting a Local Group mass M $`=(2.3\pm 0.6)\times 10^{12}\mathrm{M}_{\mathrm{}}`$ and assuming that SagDIG, which is located at a distance of $`1.29\pm 0.09`$ Mpc from the barycenter of the Local Group, lies on the Local Group zero-velocity surface one obtains a value of $`17.9\pm 2.7`$ Gyr for the dynamical age of the Local Group. This value compares favorably with the 12–16 Gyr age derived from evolutionary models of the most metal-poor stars, the $`16\pm 6`$ Gyr age that Sneden et al. (1996) derive from the Th/Eu abundance ratio in CS 22893-052, and the 10–20 Gyr ages that Cowan, Thielemann & Truran (1991) have derived from cosmochronology. If the Local Group zero-velocity surface, in fact, lies slightly beyond SagDIG, then the LG age calculated above becomes a lower limit. The figure shows that 71% of all Local Group galaxies<sup>2</sup><sup>2</sup>2In van den Bergh (2000) all galaxies within $``$1.5 Mpc were initially regarded as candidate Local Group members. Subsequently it was found that the Local Group has a zero-velocity surface with a radius of $`1.10\pm 0.15`$ Mpc, as measured from its barycenter. All galaxies within 1.25 Mpc of the barycenter have therefore been regarded as possible, or probable, Local Group members. are situated within 0.5 Mpc of the Local Group barycenter. On the other hand only one (4%) of the known Local Group galaxies brighter than M$`{}_{\mathrm{V}}{}^{}=10`$ are located beyond R$`{}_{\mathrm{LG}}{}^{}=1.0`$ Mpc. The total luminosity of the Local Group is found to be M$`{}_{\mathrm{V}}{}^{}=22.0`$, of which only 0.5% originates beyond R$`{}_{\mathrm{LG}}{}^{}=0.5`$ Mpc. This strongly supports Hubble’s (1936) claim that the Local Group “is isolated in the general field”. The only Local Group member, in addition to SagDIG, that is known to have R$`{}_{\mathrm{LG}}{}^{}>1.0`$ Mpc is the faint Tucanae system (M$`{}_{\mathrm{V}}{}^{}=9.6`$) at R$`{}_{\mathrm{LG}}{}^{}=1.1`$ Mpc.
Among the most important Local Group problems that remain to be resolved are the following:
* Why does the luminosity distribution of Local Group galaxies contain at least an order of magnitude fewer faint dwarfs than theory (Klypin et al. 1999) predicts?
Is theory wrong, or have the majority of low-mass dwarfs formed no stars? A deep search for stars associated with compact high-velocity clouds might answer this question.
* Why does the Large Magellanic Cloud contain very old globular clusters that have disk kinematics, while M 33 is embedded in an apparently younger globular cluster system that exhibits halo kinematics?
* Why did the LMC form so few clusters with ages between 4 Gyr and 10 Gyr? Was there a similar hiatus in the rate of star formation, or did the fraction of all Large Cloud stars that ended up in clusters suddenly increase $``$4 Gyr ago?
* Did any of the Local Group galaxies change their morphological type during the last $``$10 Gyr? It would be particularly important to establish if the Bar of the LMC is a relatively young feature.
* Why did different Local Group dwarf galaxies have such different star forming histories, and how do such differences depend on environment?
* It would be important to strengthen and confirm Freeman’s (1999) conclusion that the metal-rich r<sup>1/4</sup> halo of M 31 resulted from violent relaxation after the merger of two massive ancestral galaxies, whereas the metal-poor outer halo of the Milky Way system was mainly built up via capture of numerous low-mass metal-poor objects.
The new generation of powerful optical, radio and space telescopes, in conjunction with improved wide-field and infrared-sensitive detectors, should enable us to answer many of these questions in the first years of the next millennium.
I am indebted to Alan Whiting for making his data on the Cetus system available before publication. |
warning/0001/astro-ph0001221.html | ar5iv | text | # Electromagnetic showers in a strong magnetic field
## 1 Introduction
Electromagnetic showers are a universal phenomenon. Besides occurring in matter or radiation field cascade, multiplication of electrons and photons can arise in a strong magnetic field. Such super-strong fields ($`10^{12}\mathrm{G}`$) probably exist in the vicinity of some astrophysical objects such as pulsars, for example. In this case rotating neutron stars induce strong electric fields above the polar cap. Accelerated by these fields, high-energy particles (with energy up to $`10\mathrm{TeV}`$) move along curved magnetic field lines and emit curvature photons. The energy of these photons is enough to produce electron-positron pairs in magnetic and electric fields. The subsequent quantized synchrotron radiation by pairs will convert to a second generation of pairs and then an electromagnetic cascade develops in the pulsar magnetosphere. The shower development determines, to a considerable extent, the properties of the observed radiation from these objects. This was proposed for the first time in . Since electromagnetic cascades in strong magnetic fields were considered in many works mainly in connection with specific models of radio pulsars, gamma-ray bursts, blazars (see, for example, -).
It is well known that the essentially non-zero probabilities for magnetic bremsstrahlung and pair production require both strong field and high energies . The relevant parameter determining the criteria for this is:
$$\chi =\frac{\epsilon }{mc^2}\frac{H}{H_{cr}}$$
where $`\epsilon `$ is the particle energy, $`H`$ is the magnetic field strength, $`m`$ is the electron mass and $`H_{cr}=4.41\times 10^{13}\mathrm{G}`$.
The total probabilities (cross sections) for radiation and pair production for a given value of the magnetic field strength depend only on $`\chi `$ and are shown in figure 1. One can see that magnetic pair production has significant probability for $`\chi 0.1`$ (photon energy must be $`2mc^2`$). For effective shower development one needs even higher values of $`\chi `$ ($`\chi 1`$) because with increasing $`\chi `$ the radiated photon spectrum becomes harder. For $`\chi 1`$ (quantum region) the energy of the radiated photon is of the order of the electron energy. It is interesting to note that for a photon with energy $`7,5\times 10^{19}\mathrm{eV}`$ even Earth’s magnetic field ($`0.3\mathrm{G}`$) is strong enough ($`\chi 1`$) to be a good environment for creating an electromagnetic shower. If such extremely high-energy photons are presented in the primary cosmic ray flux they will undergo cascading in the geomagnetic field before entering the Earth’s atmosphere. This problem has been intensively discussed recently in connection with the detection of the highest cosmic ray events .
As mentioned above, most of the treatments of cascades in a magnetic field were connected with specific models of astrophysical objects and are numerical in nature. The most general treatment of cascade properties emphasizing an analytical approach is made in , where the steady-state kinetic equations for the electron-positron and photon distributions are solved in a strong magnetic field.
A different approach to the shower study in a magnetic field is applied in . It is similar to those for showers in matter and was motivated by the study of the primary gamma rays with extremely high energies ($`10^{20}\mathrm{eV}`$) propagating through the geomagnetic field and Earth’s atmosphere. The average shower characteristics obtained by numerically solving the system of cascade equations, show some of the main features of the cascade. While the shower is similar to those in matter for $`\chi 1`$, its nature changes sharply for $`\chi 1`$ which is connected with sharp increase of the photon free path.
Recently, a kinetic theory of electromagnetic showers in a strong magnetic field has been developed in a similar to the cascade theory in matter, in approximation A . Electromagnetic shower theories have been developed since 1937 following the works of Bhabha and Heitler and Carlson and Oppenheimer . Landau and Rumer developed a complete theory in approximation A. Their work contains the formalism which is widely used in later shower theories.
In the shower theory the mathematical description of the cascade process is based on the Boltzmann kinetic equation for particle flux density. The system of integro-differential equations of the one-dimensional Landau-Rumer theory is universal because it describes the shower development in any substance. Only the expressions for pair creation and bremsstrahlung probabilities per unit length are different. Using asymptotic forms of both processes for very high energies in a strong magnetic field, analytic formulae similar to those of standard cascade theory in approximation A for one-dimensional shower characteristics are obtained in . But, as in matter, the kinetic equations were solved within certain approximations.
In this work we give the results from the Monte Carlo simulation of the longitudinal development of electromagnetic showers in a strong magnetic field for $`\chi 1`$. We present shower profiles for different ratios of the primary and threshold energies, $`E_0/E`$, and energy spectra of shower particles at different depths. We analyse the behaviour of the shower maximum with respect to $`E_0/E`$. We compare our modelled results with theoretical ones in order to estimate theoretical approximations and the range of their validity.
## 2 The probability functions
The main elementary processes leading to particle multiplication in a magnetic field are magnetic bremsstrahlung and magnetic pair production. The corresponding probabilities per unit length are :
$`\pi (\epsilon ,\omega )\omega ={\displaystyle \frac{\alpha m^2}{\pi \sqrt{3}}}{\displaystyle \frac{\omega }{\epsilon ^2}}\left[\left({\displaystyle \frac{\epsilon \omega }{\epsilon }}+{\displaystyle \frac{\epsilon }{\epsilon \omega }}\right)K_{\frac{2}{3}}\left({\displaystyle \frac{2u}{3\chi }}\right){\displaystyle _{\frac{2u}{3\chi }}^{\mathrm{}}}K_{\frac{1}{3}}\left(y\right)y\right]`$
(1)
$`\gamma (\omega ,\epsilon )\epsilon ={\displaystyle \frac{\alpha m^2}{\pi \sqrt{3}}}{\displaystyle \frac{\epsilon }{\omega ^2}}\left[\left({\displaystyle \frac{\omega \epsilon }{\epsilon }}+{\displaystyle \frac{\epsilon }{\omega \epsilon }}\right)K_{\frac{2}{3}}\left({\displaystyle \frac{2u_1}{3\chi }}\right)+{\displaystyle _{\frac{2u_1}{3\chi }}^{\mathrm{}}}K_{\frac{1}{3}}\left(y\right)y\right]`$
where $`\epsilon `$ and $`\omega `$ are the electron and photon energy and $`u=\frac{\omega }{\epsilon \omega }`$$`u_1=\frac{\omega ^2}{\epsilon (\omega \epsilon )}`$. Parameter $`\chi `$ was defined above. Here $`\mathrm{}=c=1`$. $`K_\nu \left(z\right)=_0^{\mathrm{}}\mathrm{}^{z\mathrm{ch}(t)}\mathrm{ch}\left(\nu t\right)t`$ is a modified Bessel function known as MacDonald’s function. Parameter $`\nu `$ can have any real or complex values, here $`\nu =\frac{2}{3}`$ and $`\frac{1}{3}.`$ For simplicity we assume that the electron (positron) is moving perpendicular to the magnetic field $`H`$. As already mentioned, the probabilities of both processes are essentially different from zero under the condition that the parameter $`\chi 1`$ , which means that the particle energy $`\epsilon >>\epsilon _c`$, where $`\epsilon _c=mc^2\frac{H_{cr}}{H}`$. It is mentioned in that this condition corresponds to the conditions $`u/\chi <<1`$ and $`u_1/\chi <<1`$, i.e. one can use in (1) the asymptotic form of $`K_\nu \left(z\right)`$ for $`z<<1`$, $`K_\nu (z1)=\frac{\mathrm{\Gamma }(\nu )}{2^{1\nu }}(\frac{1}{z})^\nu +\mathrm{}`$. Then expressions (1) may be simplified:
$`\pi (\epsilon ,\omega )`$ $`=`$ $`q\left[{\displaystyle \frac{\left(\epsilon \omega \right)^{\frac{5}{3}}}{\epsilon ^{\frac{7}{3}}\omega ^{\frac{2}{3}}}}+{\displaystyle \frac{1}{\left(\epsilon \omega \right)^{\frac{1}{3}}\epsilon ^{\frac{1}{3}}\omega ^{\frac{2}{3}}}}\right],`$
$`\gamma (\omega ,\epsilon )`$ $`=`$ $`q\left[{\displaystyle \frac{\left(\omega \epsilon \right)^{\frac{5}{3}}}{\omega ^{\frac{8}{3}}\epsilon ^{\frac{1}{3}}}}+{\displaystyle \frac{\epsilon ^{\frac{5}{3}}}{\omega ^{\frac{8}{3}}\left(\omega \epsilon \right)^{\frac{1}{3}}}}\right],`$
where $`q=3.9\times 10^6\left[\frac{H}{H_c}\right]^{\frac{2}{3}}\frac{\mathrm{GeV}^{\frac{1}{3}}}{\mathrm{cm}}`$. Using $`u=\omega /\epsilon `$ (correspondingly $`u=\epsilon /\omega `$) we can rewrite (LABEL:eq2) in the form:
$`\pi (\epsilon ,\omega )d\omega `$ $`=`$ $`{\displaystyle \frac{q}{\epsilon ^{\frac{1}{3}}}}\left[{\displaystyle \frac{\left(1u\right)^{\frac{5}{3}}}{u^{\frac{2}{3}}}}+{\displaystyle \frac{1}{\left(1u\right)^{\frac{1}{3}}u^{\frac{2}{3}}}}\right]u,`$
$`\gamma (\omega ,\epsilon )d\epsilon `$ $`=`$ $`{\displaystyle \frac{q}{\omega ^{\frac{1}{3}}}}\left[{\displaystyle \frac{\left(1u\right)^{\frac{5}{3}}}{u^{\frac{1}{3}}}}+{\displaystyle \frac{u^{\frac{5}{3}}}{\left(1u\right)^{\frac{1}{3}}}}\right]u.`$
The total probabilities per unit length for bremsstrahlung and pair production are given by
$`W_r\left(\epsilon \right)`$ $`=`$ $`{\displaystyle _0^\epsilon }\pi (\epsilon ,\omega )\omega ={\displaystyle \frac{q}{\epsilon ^{\frac{1}{3}}}}{\displaystyle _0^1}\left[{\displaystyle \frac{\left(1u\right)^{\frac{5}{3}}}{u^{\frac{2}{3}}}}+{\displaystyle \frac{1}{\left(1u\right)^{\frac{1}{3}}u^{\frac{2}{3}}}}\right]u`$ (4)
$`=`$ $`5.642{\displaystyle \frac{q}{\epsilon ^{\frac{1}{3}}}},`$
$`W_p\left(\omega \right)`$ $`=`$ $`{\displaystyle _0^\omega }\gamma (\omega ,\epsilon )\epsilon =1.467{\displaystyle \frac{q}{\omega ^{\frac{1}{3}}}}.`$ (5)
Unlike the Bethe-Heitler probability,$`\pi (\epsilon ,\omega )`$ does not contain an infrared divergence and because of this $`W_r\left(\epsilon \right)`$ is finite.
As mentioned earlier, the probabilities of both processes were obtained using the asymptotic form of the function $`K_\nu \left(z\right)`$ for $`z<<1`$. This explains the behaviour of the total cross sections as a function of the particle energy of power $`(\frac{1}{3})`$ (at fixed $`H`$) and in figure 1 this is the region of $`\chi 1`$. Expression (5) coincides with the expression of the photon attenuation coefficient given in Erber’s review when the asymptotic form of auxiliary function $`T(\chi )`$ for $`\chi >>1`$, $`T(\chi )0.60\chi ^{\frac{1}{3}}`$, is used. Here $`\chi =\frac{1}{2}\frac{h\nu }{mc^2}\frac{H}{H_{cr}}`$, $`h\nu `$ is the photon energy.
## 3 Simulation
To investigate the shower characteristics in a strong magnetic field we developed our own Monte Carlo code. The probabilities (LABEL:eq3) were used to sample energies of the secondary particles - photon in bremsstrahlung and electron in pair creation. We constructed tables with cumulative distributions from probabilities (LABEL:eq3) through a small step by $`u`$. The mean interaction length is the inverse value of the total probability, (4) and (5).
It was shown in that the average ranges of both electron and photon, with energy $`E_0`$ in a magnetic field $`H`$, are of the same order of magnitude, $`\frac{E_0^{\frac{1}{3}}}{q}`$. The quantity $`L=\frac{E_0^{\frac{1}{3}}}{q}`$ which includes the matter parameters (magnetic field strength $`H`$) could play the role of a radiation length. In this case $`E_0`$ is the energy of the particle initiating the shower. It is important to note, that unlike matter, here $`L`$ depends on the energy $`E_0`$ of the primary particle. $`L`$ is a relatively small quantity. For example, in magnetic field $`H=10^5\mathrm{G}`$ for $`E_0=10^6\mathrm{GeV}`$$`L=14.86\mathrm{cm}`$, for $`E_0=10^9\mathrm{GeV}`$$`L=148.6\mathrm{cm}`$.
We considered the problem as one dimensional, i.e. we assumed that the shower is only developing in the direction of the primary particle entering the magnetic field at $`t=0`$. The distance is measured in units $`L`$. All shower particles were followed down to some threshold energy $`E`$.
## 4 Results and discussion
We performed simulation for various sets of primary and threshold energies ($`E_0`$ and $`E`$) and different magnetic field strengths $`H.`$
Our results confirmed the theoretical prediction that the above-mentioned quantity, $`L`$, plays the role of a radiation length. Similar to the standard shower theory under approximation A, the longitudinal cascade development is independent of an absorber (magnetic field) when distances are measured in radiation lengths and the average behaviour of a shower is expressed by a function of $`E_0/E.`$
Shower profiles for electron-initiated showers and different $`E_0/E`$ are shown in figure 2. The corresponding curves for photons are very close to those of electrons and because of this they are not shown. When the primary particle is a photon (figure 3), the shower maximum is shifted by $`\frac{1}{2}\mathrm{r}.\mathrm{l}.`$ (where r.l. denotes radiation length) deeper, which is close to the difference between the mean interaction lengths of primary electron and photon.
The differential energy spectra of shower particles at different stages of the shower development are shown in figure 4. The depth of $`2.25\mathrm{r}.\mathrm{l}`$. is near the shower maximum. The spectra can be described by power-law of energy, $`E\frac{n}{E}E^\delta `$, with $`\delta `$ increasing with the depth and approaching $`1`$ for $`t\mathrm{}`$.
As can be seen, the typical distance over which the shower develops is a few radiation lengths. The depth of the shower maximum increases very slowly with the increase of $`E_0/E`$ approaching a limit. In very strong fields, e.g. in pulsars, $`L`$ becomes very small ($`10^4\mathrm{cm}`$ for $`E_0=10^6\mathrm{GeV}`$) which means that the shower has the longitudinal spread of the same order. This confirms the theoretical prediction in that the strong magnetic fields are effective screens for very high energy electrons, positrons and photons.
It should be noted here that these properties of the shower are valid for such values of $`E_0`$, $`E`$ and $`H`$ that satisfy the condition $`\chi 1`$ and the asymptotic expressions (2) can also be used. In addition, to consider the shower as one-dimensional, the electron energies must obey the condition
$$\left(\frac{\epsilon }{mc^2}\right)^{\frac{2}{3}}\left(\frac{H_{cr}}{H}\right)^{\frac{1}{3}}\frac{\alpha }{2\pi }1.$$
This criteria comes from the requirement that the electron gyroradius $`R=\epsilon /eH`$ must be much greater than the typical length of the shower, i.e. $`RL`$.
However, detailed comparison of a numerical results shows that the Monte Carlo cascade curves differ significantly from the theoretical ones. This is illustrated in figure 5 where shower profiles for the electron-induced showers and $`E_0/E=1000`$ are shown.
This disagreement could probably be explained with the approximations used in to get a solution of cascade equations. In the case of a magnetic field, the shower theory is more complicated than the conventional theory in approximation A because the probabilities $`\pi (\epsilon ,\omega )\omega `$ and $`\gamma (\omega ,\epsilon )\epsilon `$ are not scaling functions, i.e. they do not depend only on the ratio of secondary to primary energies (functions of $`\omega /\epsilon `$or $`\epsilon /\omega `$ in our notation). Mellin transforms lead to differential-difference equations for the distribution functions and its solutions are found by authors of in the so-called adiabatic and modified adiabatic approximations. As pointed out in , the adiabatic approximation (AA) is limited to the region after the maximum. In the modified adiabatic approximation it is assumed that the explicit dependence of $`\pi (\epsilon ,\omega )`$ and $`\gamma (\omega ,\epsilon )`$ on primary energy can be eliminated by its substitution with some mean energy per interaction in the shower. This leads to the same solutions with the modified radiation length but, obviously, the shower development is distorted. If we use the probabilities $`\pi (\epsilon ,\omega )`$ and $`\gamma (\omega ,\epsilon )`$ in this modified way in our Monte Carlo code we obtain results very close to the theoretical ones. Figure 6 shows the number of electrons in the shower maximum $`N_{\mathrm{max}}`$ as function of $`E_0/E`$. One can see that the rise of $`N_{\mathrm{max}}`$ with $`E_0/E`$ is slower than that of AA of the theory. $`N_{\mathrm{max}}`$ is exactly proportional to $`\left(E_0/E\right)^\xi `$ . AA gives $`\xi `$ close to one as it is in the standard theory with Bethe-Heitler cross sections (curve labelled BH in the figure) while our simulation (curve labelled MC) gives $`\xi \frac{2}{3}`$. Figure 7 shows the depth of the maximum $`t_{\mathrm{max}}`$ as a function of $`E_0/E`$ for both theoretical and simulated showers. As can be seen, the behaviour of the modelled $`t_{\mathrm{max}}`$ is too different from those of logarithmically increasing theoretical $`t_{\mathrm{max}}`$. After some $`E_0/E`$ the modelled $`t_{\mathrm{max}}`$ almost ceases to increase.
It is easy to demonstrate this feature of showers in a magnetic field using Heitler’s elementary cascade model . In this model, electromagnetic particles subdivide into two particles with half the initial energy. In matter this takes place on every radiation length. In a magnetic field, however, the situation is substantially different. After each interaction the interaction length, which is a function of energy power of $`\frac{1}{3}`$, decreases due to particle energy splitting. If the primary particle has an energy $`E_0`$ and the interaction length $`t_0\mathrm{r}.\mathrm{l}.`$, then after the first interaction the number of particles will be two, each of energy $`E_0/2`$. The next interaction length will be $`t_0/2^{\frac{1}{3}}`$ and the number of particles four, each of energy $`E_0/4`$. After $`N`$ interactions the particle number is already $`2^N`$ and their energy $`E=E_0.2^N`$. The next interaction length will be $`\frac{t_0}{\left[2^{\frac{1}{3}}\right]^N}\mathrm{r}.\mathrm{l}.`$ The distance at which this takes place is
$`T=t_0+{\displaystyle \frac{t_0}{2^{\frac{1}{3}}}}+{\displaystyle \frac{t_0}{\left(2^{\frac{1}{3}}\right)^2}}+\mathrm{}+{\displaystyle \frac{t_0}{\left(2^{\frac{1}{3}}\right)^{N1}}}=t_0\left[1+{\displaystyle \frac{1}{2^{\frac{1}{3}}}}+{\displaystyle \frac{1}{\left(2^{\frac{1}{3}}\right)^2}}+\mathrm{}+{\displaystyle \frac{1}{\left(2^{\frac{1}{3}}\right)^{N1}}}\right].`$
The expression in the square brackets is a sum of a geometrical progression with $`q=2^{1/3}`$ and thus
$`T`$ $`=`$ $`4.847t_0\left[1{\displaystyle \frac{1}{\left(2^{\frac{1}{3}}\right)^{N1}}}\right].`$ (6)
This expression shows that if $`N`$ increases, i.e. $`E_0/E`$ increases, $`t_{\mathrm{max}}`$ approaches a limit. This means that after some $`N`$ (or $`E_0/E`$) $`t_{\mathrm{max}}`$ practically ceases to increase. In our rather simplified model the maximum $`t_{\mathrm{max}}`$is $`4.847t_0\mathrm{r}.\mathrm{l}.`$ The first interaction length $`t_0`$ is $`0.117\mathrm{r}.\mathrm{l}`$. for the electron and $`0.682\mathrm{r}.\mathrm{l}.`$ for the photon. If we take the greater value this will lead to an estimation of the maximum $`t_{\mathrm{max}}`$ of $`3.31\mathrm{r}.\mathrm{l}.`$ whose value does not contradict with the curve modelled in figure 7.
Another important characteristic of electromagnetic showers in a strong magnetic field which can be easily obtained from the simulation are fluctuations in the shower development. There is no theoretical consideration of this problem. The problem is complicated even for the standard cascade theory (see, e.g. ).
In figure 8 the fluctuations of the number of shower electrons $`\left(\sigma /N\right)`$ as a function of the depth for different $`E_0/E`$ are shown. The primary particle is an electron. The behaviour of the curves is very similar to those of BH showers but the fluctuations in a magnetic field are significantly larger. This is illustrated in figure 9 where fluctuations in the shower maximum as a function of $`E_0/E`$ are shown, compared with BH showers. Calculations for BH showers are performed by direct MC simulation in air for very high energies $`E_0`$ and $`E`$, i.e. at the conditions where the standard cascade theory in approximation A is valid.
The larger fluctuations in a magnetic field compared to BH showers are not an unexpected result. The main sources for the fluctuations of the number of particles in the cascade are fluctuations of interaction lengths and the random energy distribution of created secondary particles. Unlike the BH cross sections the mean interaction lengths for magnetic bremsstrahlung and pair-production processes are functions of the particle energy of power $`\frac{1}{3}`$ (see (4,5)). The differential cross sections (LABEL:eq3) depend on the particle energy as well. Their features lead to a significant probability that the secondary electron or photon will take a large fraction of the primary energy.
## 5 Conclusions
A direct MC simulation of the longitudinal development of electromagnetic showers in a strong magnetic field has been performed. The processes of the magnetic bremsstrahlung and of the magnetic pair production were included in the simulation with asymptotic expressions for both probabilities valid for very high energies. Simulated results were compared with the theory of showers in a strong magnetic field developed in . The main predictions of this theory - the dependence of the radiation length on magnetic field strength and the energy of the primary particle, the very small shower longitudinal spread and the closeness of electron and photon profiles, are confirmed by our simulation.
However, the AA of the theory used in was inadequate for a precise qualitative estimate of shower characteristics.
## Acknowledgments
The authors thanks A. Rekalo and T.Stanev for helpful discussions. This work was partially supported by a grant F-460 of the Bulgarian NFSR and by the Bulgarian Science and Culture Foundation.
## References |
warning/0001/astro-ph0001212.html | ar5iv | text | # The triple degenerate star WD 1704+481
## 1 Introduction
WD 1704+481 was identified as a pair of white dwarfs of similar brightness separated by about 6 arcscec by Sanduleak and Pesch (1982) from objective prism plates. Spectrophotometry of the pair by Greenstein, Dolez & Vauclair (1983) showed that, although their visual magnitudes are almost identical, the SE component (WD 1704+481.1 = Sanduleak A = GR 576, V=14.48, (B$``$V)=-0.09) is bluer and, therefore, hotter than the NW component (WD 1704+481.2 = Sanduleak B = GR 577, V=14.45, (B$``$V)=0.14). This immediately suggests that the cooler component must have a larger radius which also implies a lower mass for a cool white dwarf such as Sanduleak B. Greenstein et al. estimated a mass of 0.32 – 0.43$`\mathrm{M}_{}`$ for Sanduleak B which is well below the typical mass of white dwarfs (0.55$`\mathrm{M}_{}`$, Bergeron, Saffer & Liebert 1992).
Low mass white dwarfs such as Sanduleak B are thought to be the result of binary star evolution, in which the evolution of a star during the red giant phase is interrupted by interactions with a nearby star. The physics of this interaction is complex but it is thought to lead to the stripping of the outer hydrogen layers from the red giant in a “common-envelope” phase, halting the formation of the degenerate helium core and leading to the formation of an anomalously low mass white dwarf (Iben & Livio 1993). The hypothesis that binary star evolution forms low mass white dwarfs was confirmed by the discovery by Marsh, Dhillon & Duck (1995) of at least 5 short period binary white dwarfs in a sample of 7 low mass white dwarfs.
In this paper we present spectra of the H$`\alpha `$ line of Sanduleak B which clearly show that it is a close binary with two white dwarf components – a double degenerate star (DD). We derive the spectroscopic orbits of both components and show that it is likely to be the first known example of a binary white dwarf with one white dwarf composed of helium and one composed of carbon and oxygen.
## 2 Observations and reductions.
The data for this study come from observations obtained with the 2.5m Isaac Newton Telescope (INT) in September 1998 and the 4.2m William Herschel Telescope (WHT) in October 1998. The INT spectra were obtained with the intermediate dispersion spectrograph using the 500mm camera, a 1200 line/mm grating, a 0.9arcsec slit and a TEK charge coupled device (CCD) as a detector at a dispersion of 0.39Å per pixel. The WHT spectra were obtained using the blue arm of the ISIS spectrograph, a 1200 line/mm grating, a 0.83arcsec slit and an EEV CCD at a dispersion of 0.22Å per pixel. Integration times varied between 600s and 1200s. Each observation of the stars was bracketed by observations of a copper-neon arc. We set the angle of the slit so that spectra of both stars were obtained simultaneously. The position angle estimated from this slit angle was 295 for the INT and 290 for the WHT, which agrees well with the PA of 289$`{}_{}{}^{}\pm 5^{}`$ given by Greenstein et al. The separation of the stars measured from our INT images is 6.0 arcsec with an uncertainty of a few tenths of an arcsecond. The seeing was good for all the observations ($``$ 1arcsec) so the spectra of the stars are clearly separated in all our images. Of the 49 spectra, 41 were obtained with the INT.
The bias level in all the images determined from the clipped-mean in the overscan region was subtracted from all the images before further processing. For the INT spectra, several images of a tungsten lamp were combined to form a normalized master flat-field image for each night’s data. Similar images for the WHT spectra show mild fringing. We therefore combined flat-field images taken immediately before and after each observation of the stars in order to form a normalized flat-field image for each image of the star. Extraction of the spectra from the images was performed automatically using optimal extraction to maximize the signal-to-noise of the resulting spectra (Horne 1986). Uncertainties due to photon statistics are rigorously followed through the data reduction process so that reliable uncertainties are known for every point in the final spectra. The arcs associated with each stellar spectrum were extracted using the same weighting determined for the stellar image to avoid possible systematic errors due to tilted spectra. The wavelength scale was determined from a fourth-order polynomial fit to measured arc-line positions. We calibrated the wavelength response of the WHT spectra using observations of the standard star BD +33 2642 (Bohlin 1996). To calibrate the wavelength response of the INT spectra we used a least-squares fit of a smooth function to the ratio of the average WHT spectrum of WD 1740+481.2 and the average INT spectrum of the same star. The core of the H$`\alpha `$ line and regions affected by telluric absorption were excluded from the fit. We then normalized all the spectra using a linear fit to the continuum either side of the H$`\alpha `$ line.
We measured the resolution of the spectra by fitting a Gaussian profile by least-squares to a neon arc line at 6532Å in a series of spectra. The full-width at half-maximum (FWHM) of the model profile was typically 0.94Å for the INT spectra and 0.45Å for the WHT spectra. The INT spectra for one night are affected by poor spectrograph focus and have a resolution of 1.2Å.
## 3 Measurement of the radial velocities.
Our spectra of the cooler component of the visual binary clearly show two sharp cores to the H$`\alpha `$ line (Fig. 1) which vary in position periodically in a sinusoidal manner, indicating the presence of two white dwarfs in a close binary orbit. An initial estimate of the period was obtained by measuring the radial velocity of the deeper core using a model profile. We first used a simultaneous least-squares fit to all the spectra of four Gaussian profiles with a common position but independent depths and widths to determine the model profile. The shape of the profile was then fixed in a least-squares fit to the individual spectra in which only the position of the model profile was allowed to vary. These initial velocities are affected by the presence of the weaker core, but were good enough to determine the orbital period. A periodogram of these radial velocities showed a clear peak at 6.91 cycles per day, which confirmed our initial impression that the period is around 0.145d.
To measure the radial velocities of the two components more precisely we used a simultaneous fit to all the spectra of two model profiles, one for each star, in which the position of each model profile is predicted from a circular orbit of the form $`\gamma +K\mathrm{sin}(2\pi (TT_0)/P)`$. In this way we are able to determine the shape of the two profiles and the parameters of the two circular orbits directly. The (variable) resolution of each of the spectra and the effects of smearing due to orbital motion are included in the fitting process. There are many free-parameters in this fitting process so we used a series of least-squares fits in which first the profile shapes were fixed while the parameters of the orbit were varied and then vice versa, until we had established values for all the parameters which were nearly optimal. Only data within 5000 km s<sup>-1</sup> of H$`\alpha `$ and unaffected by telluric absorption are included in the fitting process. We used four Gaussian profiles to model the broad wings of the H$`\alpha `$ line and the core of star with the deeper core and two Gaussian profiles for the other star. A polynomial is also included in the fitting process to allow for smooth, asymmetric features in the profile. For the final least-squares fit the parameters of the profile shapes and the orbit were all varied independently. The parameters of this final fit are given in Table 1. We designate the star with the deeper H$`\alpha `$ core component B and the star with the weaker H$`\alpha `$ core is then component C as we will refer to Sanduleak A as component A throughout to avoid confusion.
We also measured the radial velocities of Sanduleak A using the same method employed to measure the initial radial velocities of Sanduleak B. There is no significant variability in these radial velocities, which have an average value of $`\gamma _\mathrm{A}=0.6\pm 0.3\mathrm{km}\mathrm{s}^1`$.
## 4 Discussion
The first obvious result of our analysis is the mass ratio $`q=\frac{M_\mathrm{B}}{M_\mathrm{C}}=0.70\pm 0.03`$. WD 1704+481.2 is one of several double degenerates (DDs) which have been identified among low mass white dwarfs ($`M0.49`$$`\mathrm{M}_{}`$) where the fainter companion is sufficiently young, i.e, hot, for its H$`\alpha `$ core to be detected and a mass ratio derived (Moran, Marsh & Maxted, 1999). These observed mass ratios for DDs are shown in Fig. 2. Component C has the weaker core in our H$`\alpha `$ spectra. The depth of the core does not vary strongly with temperature over the temperature range expected for components B and C (Koester et al. 1998), so we can confidently state that component C is the fainter component. For comparison, we also plot the results of Han (1998). We applied a set of selection criteria to the simulated population of white dwarf binaries to approximate these selection effects, namely, that the younger white dwarf is less massive than 0.49$`\mathrm{M}_{}`$and that the older white dwarf is younger than 1 Gyr. The mass ratio distribution for white dwarf binaries predicted by the models of Iben, Tutukov & Yungelson (1997) using similar selection criteria but their own “standard model” is quite similar, i.e. it shows a single peak near $`q=0.7`$. WD 1704+481.2 lies near the peak of the theoretical distribution, which is a success for the theoretical model. What is less clear is why the majority of DDs with measured mass ratios do not have mass ratios which accord with the models.
We can use the observed difference $`\gamma _\mathrm{B}\gamma _\mathrm{C}=11.5\pm 2.3\mathrm{km}\mathrm{s}^1`$ to determine the masses of the stars as follows. The gravitational redshift of the light from the stars is given by $`v_{\mathrm{gr}}=0.635(M/\mathrm{M}_{})/(R/\mathrm{R}_{})\mathrm{km}\mathrm{s}^1`$. Since one star is more massive than the other and more massive white dwarfs have smaller radii (for a given temperature), the more massive component will appear to have a larger (more positive) systemic velocity, as is observed. To determine the radius of component B we use the models of Althaus & Benvenuto (1997) for helium white dwarfs. We assume a temperature of 9500K for component B (Greenstein, Dolez & Vauclair 1983) and use parabolic interpolation to find the radius. Note that a 1000K difference in temperature changes the radius derived by less than one percent so the temperature assumed has a negligible effect on our results. We assume that component C is slightly cooler (8500K) and use the mixed-composition models of Wood (1990) to determine its radius. In Fig. 3 we show the predicted gravitational redshift difference, $`\mathrm{\Delta }v_{\mathrm{gr}}`$ as a function of the mass of each component. The observed value of $`\mathrm{\Delta }v_{\mathrm{gr}}`$ is also indicated and we see that this observed value implies $`M_\mathrm{B}=0.39\pm 0.05\mathrm{M}_{}`$ and $`M_\mathrm{C}=0.56\pm 0.07\mathrm{M}_{}`$. White dwarfs more massive than about 0.5$`\mathrm{M}_{}`$will have passed through a helium-burning stage and have a carbon-oxygen composition, so these masses suggest that component B is a white dwarf composed of helium and component C is composed of carbon and oxygen. In Fig. 4 we see that the masses of the two components are roughly what would be expected from the theoretical models of Han (1998).
The mass of component C is typical for single white dwarfs (Bergeron, Saffer & Liebert 1992) and component A has the same gravitational redshift. The change in radius with temperature in this regime is very small, so we expect component A has a normal mass. Component A is $``$200 AU distant from components B and C (Greenstein et al. 1983) and the initial separation of components B and C must have been smaller than $`1`$ AU for a common-envelope phase to have occurred. The ratio of the orbital periods is sufficiently large that we need not worry about dynamical instability of the orbits (Kiseleva, Eggleton & Anosova 1994) and so we can be confident that the evolution of the inner binary will have been unaffected by the presence of component A and vice versa.
From the parameters of the orbit we find $`M_\mathrm{B}\mathrm{sin}^3i=0.075\pm 0.003\mathrm{M}_{}`$ and $`M_\mathrm{C}\mathrm{sin}^3i=0.108\pm 0.006\mathrm{M}_{}`$ which, combined with the masses derived above, yields an inclination of $`61^{}`$. The separation of components B and C is only 0.74$`\mathrm{R}_{}`$ but the small radii of white dwarfs ($`0.01\mathrm{R}_{}`$) rules out the possibility of observable eclipses in this binary.
## 5 Conclusion
We have shown that WD 1704+481 is a hierarchical triple star in which all three components are white dwarfs, i.e., a triple degenerate star. The outermost star, component A appears to be a typical white dwarf. Components B and C are a close binary with an orbital period of only 0.145d. The mass ratio ($`q=\frac{M_{\mathrm{bright}}}{M_{\mathrm{faint}}}=\frac{M_\mathrm{B}}{M_\mathrm{C}}=0.70\pm 0.03`$) for the inner binary is the first measured mass ratio for a double degenerate which is close to the peak of the mass ratio distribution predicted by theoretical models. Similarly, the masses of components B and C derived from the difference between their gravitational redshifts ($`M_\mathrm{B}=0.39\pm 0.05\mathrm{M}_{}`$ and $`M_\mathrm{C}=0.56\pm 0.07\mathrm{M}_{}`$) appear near the peaks of the theoretical mass distributions. These suggest that the component C is a typical white dwarf composed of carbon and oxygen and that component B is composed of helium.
## Acknowledgements
PFLM was supported by a PPARC post-doctoral grant. CM was supported by a PPARC post-graduate studentship. The William Herschel Telescope and the Isaac Newton Telescope are operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias. |
warning/0001/cond-mat0001415.html | ar5iv | text | # Two-channel Kondo model as a generalized one-dimensional inverse square long-range Haldane-Shastry spin model
\[
## Abstract
Majorana fermion representations of the algebra associated with spin, charge, and flavor currents have been used to transform the two-channel Kondo Hamiltonian. Using a path integral formulation, we derive a reduced effective action with long-range impurity spin-spin interactions at different imaginary times. In the semiclassical limit, it is equivalent to a one-dimensional Heisenberg spin chain with two-spin, three-spin, etc. long-range interactions, as a generalization of the inverse-square long-range Haldane-Shastry model. In this representation the elementary excitations are ”semions”, and the non-Fermi-liquid low-energy properties of the two-channel Kondo model are recovered.
PACS numbers: 72.15.Qm, 71.10.+x, 71.27.+a, 75.20.Hr
\]
The two-channel Kondo model is known to have a non-Fermi liquid (non-FL) low-energy fixed point in the overscreening case, and has been put forward as a model to explain non-FL behavior observed in several quite different physical systems at low temperatures, such as certain heavy fermion alloys and two-level systems . There are exact solutions for the ground state and thermodynamics of this model derived from the Bethe Ansatz , but there are still continuing efforts to find an intuitive understanding of the nature of excitations in the neighborhood of the low-energy fixed point. Numerical renormalization group , conformal field theory (CFT) , and bosonization approach made a number of predictions for the many-body excitations at the fixed point, but they do not provide an intuitive interpretation of these excitations as it is possible at the FL fixed point of the single-channel Kondo model .
It was shown a long time ago that the single-channel Kondo model can be reduced to an inverse-square one-dimensional Ising model , which is a prototype classical model in statistical physics . Moreover, such a reduced effective model had helped Anderson and coworkers to establish the correct FL behavior of the low-energy fixed point for the one-channel model in the early 70s , namely, the Ising spin-spin correlation function should behave as $`\frac{1}{\tau ^2}`$ with $`\tau `$ as the imaginary time. On the other hand, the non-FL behavior of the two-channel Kondo model is characterized by the dynamic correlation function of the impurity spin, $`<T_\tau \stackrel{}{S_d}(\tau )\stackrel{}{S_d}(0)>\frac{1}{\tau },`$ but its physical meaning, unfortunately, has not been fully understood yet. To our knowledge, the possible connection of the overscreened two-channel Kondo model with quantum spin models has not been explored so far.
In this Letter, the algebra of total spin, charge, and flavor currents of the two-channel Kondo model have been represented in terms of Majorana fermions, and in the path integral formulation, we derive a reduced equivalent quantum Heisenberg spin model, in which the impurity spins at different imaginary times are strongly correlated, including two-spin, three-spin, etc. long-range exchange interactions. In particular, the two-spin interaction has exact the same form as the integrable inverse-square one-dimensional Heisenberg spin chain — the so-called Haldane-Shastry (HS) model , while the three-spin and four-spin interaction parts, etc. are natural generalizations of the HS model. The non-FL fixed point action of the two-channel Kondo model is identified with the two-spin interaction part (the HS model), and the elementary excitations of this low-energy non-FL fixed point are spinons ($`S=1/2`$ objects) obeying semion (half-fractional) statistics intermediate between bosons and fermions . Actually, the two-spin long-range interaction part forms an ideal semion gas, while the three-spin long-range term is a dangerous irrelevant interaction, leading to important corrections to the thermodynamic properties. The fourth and higher order terms are irrelevant variables. The non-FL properties of the two-channel Kondo model are thus recovered.
We start with the Hamiltonian of the two-channel Kondo model in the form,
$`H=H_0+H_I`$ (1)
$`H_0={\displaystyle \frac{v_f}{2\pi }}{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \underset{\sigma =,}{}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\psi _{j,\sigma }^{}(x)(i_x)\psi _{j,\sigma }(x)`$ (2)
$`H_I={\displaystyle \underset{a=x,y,z}{}}J_aS_d^aJ_s^a(0),`$ (3)
where we have retained only the s-wave scattering, linearized the fermion spectrum, and replaced the incoming and outgoing waves with two left-moving electron fields $`\psi _{j,\sigma }(x)`$. $`J_s^a(x)`$ are the conduction electron spin current operators: $`J_s^a(x)=_{j,\sigma ,\sigma ^{}}:\psi _{j,\sigma }^{}(x)s_{\sigma ,\sigma ^{}}^a\psi _{j,\sigma ^{}}(x):`$, where $`s^a`$ being spin-1/2 matrices, :: means normal ordering. We introduce charge and flavor currents: $`J_c(x)=_{j,\sigma }:\psi _{j,\sigma }^{}(x)\psi _{j,\sigma }(x):`$ and $`J_f^a(x)=_{j,j^{},\sigma }:\psi _{j,\sigma }^{}(x)t_{j,j^{}}^a\psi _{j^{},\sigma }(x):`$, where $`t_{j,j^{}}^a`$ are generators of an SU(2) symmetry group. Following Affleck and Ludwig , the free part of the Hamiltonian can be rewritten as a sum of three commuting terms by the usual point-splitting procedure (Sugawara construction):
$`H_0={\displaystyle \frac{v_f}{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x`$ $`[{\displaystyle \frac{1}{8}}:J_c(x)J_c(x):+{\displaystyle \frac{1}{4}}:\stackrel{}{J}_f(x)\stackrel{}{J}_f(x):`$ (5)
$`+{\displaystyle \frac{1}{4}}:\stackrel{}{J}_s(x)\stackrel{}{J}_s(x):],`$
while the interaction term is expressed only in terms of the electron spin currents and the impurity spin. The information about the number of channels is contained in the commutation relations obeyed by the spin currents
$`[J_s^a(x),J_s^b(x^{})]=iϵ^{abc}J_s^a(x)\delta (xx^{})+{\displaystyle \frac{ki}{4\pi }}\delta _{a,b}\delta ^{}(xx^{})`$indicating that $`J_s^a(x)`$ form an SU(2) level $`k=2`$ Kac-Moody algebra. Meanwhile, the charge and flavor currents satisfy
$`[J_c(x),J_c(x^{})]`$ $`=`$ $`2ki\delta ^{}(xx^{}),`$
$`[J_f^a(x),J_f^b(x^{})]`$ $`=`$ $`iϵ^{abc}J_f^a(x)\delta (xx^{})+{\displaystyle \frac{ki}{4\pi }}\delta _{a,b}\delta ^{}(xx^{}).`$
They form a U(1) Kac-Moody and another SU($`k=2`$) level-2 Kac-Moody algebra, respectively.
It is now quite natural to introduce a Majorana representation of the spin current operators,
$`J_s^x(x)=i\chi _2(x)\chi _3(x),`$ (6)
$`J_s^y(x)=i\chi _3(x)\chi _1(x),`$ (7)
$`J_s^z(x)=i\chi _1(x)\chi _2(x),`$ (8)
where $`\chi _1(x),\chi _2(x),`$ and $`\chi _3(x)`$ are left-moving free Majorana fermion fields, and it can be shown to reproduce the SU(2) level-2 Kac-Moody commutation relations. It is important to note that this representation is only appropriate for the two-channel model as it leads to a level-2 algebra. It would be inappropriate for the single-channel Kondo model where the corresponding spin current generates a level-1 algebra.
In a similar way, we can also introduce Majorana representations for the flavor currents
$`J_f^x(x)=i\chi _2^{}(x)\chi _3^{}(x),`$ (9)
$`J_f^y(x)=i\chi _3^{}(x)\chi _1^{}(x),`$ (10)
$`J_f^z(x)=i\chi _1^{}(x)\chi _2^{}(x),`$ (11)
which reproduces the commutation relations satisfied by the flavor currents, and $`J_c(x)=2i\chi _4^{}(x)\chi _5^{}(x)`$ representing the charge current operator. Note that $`\chi _\alpha ^{}`$ with $`\alpha =1,2,3,4,5`$ are also left-moving free Majorana fermion fields. It is well-known that the dynamics of charge, flavor, and spin are completely determined by the commutation relations of these current operators. Though the spin currents of the two-channel Kondo model can be represented in terms of three Majorana fermion fields $`\chi _\alpha (x)`$ ($`\alpha =1,2,3`$), they can not be given any simple physical interpretation in terms of the original conduction electrons $`\psi _{j,\sigma }(x)`$.
Now using these current operators, the Hamiltonian (2) is presented as a quartic form in the Majorana fields. This form is convenient if one pursues a purely algebraic approach as in the CFT . However, for our purpose it is more convenient to perform an inverse Sugawara construction using again the point-splitting procedure, and rewrite the terms quartic in the Majorana fermions as kinetic energy terms which are quadratic :
$`:`$ $`J_c(x)J_c(x):=4{\displaystyle \underset{\alpha =4}{\overset{5}{}}}\chi _\alpha ^{}(i_x)\chi _\alpha ^{}(x);`$ (12)
$`:`$ $`\stackrel{}{J}_f(x)\stackrel{}{J}_f(x):=2{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\chi _\alpha ^{}(i_x)\chi _\alpha ^{}(x);`$ (13)
$`:`$ $`\stackrel{}{J}_s(x)\stackrel{}{J}_s(x):=2{\displaystyle \underset{\alpha =1}{\overset{3}{}}}\chi _\alpha (i_x)\chi _\alpha (x).`$ (14)
The model Hamiltonian is thus transformed into the following two parts ,
$`H_c+H_f={\displaystyle \frac{v_f}{4\pi }}{\displaystyle \underset{\alpha =1}{\overset{5}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\chi _\alpha ^{}(x)(i_x)\chi _\alpha ^{}(x),`$ (15)
$`H_s={\displaystyle \frac{v_f}{4\pi }}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\chi _\alpha (x)(i_x)\chi _\alpha (x)`$ (16)
$`\text{ }{\displaystyle \frac{iJ}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\delta (x)\text{ }\stackrel{}{S}_d\stackrel{}{\chi }(x)\times \stackrel{}{\chi }(x).`$ (17)
$`H_c+H_f`$ describes the non-interacting charge and flavor degrees of freedom. It is invariant under the symmetry group $`U(1)SU(2)_2`$ described by the five free Majorana fermion fields $`\chi _\alpha ^{}(x)`$ ($`\alpha =1,2,3,4,5`$). $`H_s`$ is the main part of the model and it describes the spin degrees of freedom with three left-moving Majorana fermion fields $`\chi _\alpha `$ ($`\alpha =1,2,3`$) interacting with the impurity spin. It has the $`SU(2)_2`$ symmetry so that the full Hamiltonian is described by eight different Majorana fermion fields.
In the two-channel model Hamiltonian, $`H_s`$ will give rise to the essential low-energy physics of the model because it is the only part which includes the interaction. The reduced partition function of the model Hamiltonian with $`H_s`$ can be written in the form of a functional integral:
$`Z`$ $`=`$ $`{\displaystyle }D\widehat{\mathrm{\Omega }}{\displaystyle }{\displaystyle \underset{\alpha =1}{\overset{3}{}}}D\chi _\alpha \mathrm{exp}\{iS_d\omega \left(\widehat{\mathrm{\Omega }}\right){\displaystyle _0^\beta }d\tau `$ (19)
$`({\displaystyle _{\mathrm{}}^+\mathrm{}}dx{\displaystyle \underset{\alpha =1}{\overset{3}{}}}{\displaystyle \frac{1}{2}}\chi _\alpha _\tau \chi _\alpha +H_s[\widehat{\mathrm{\Omega }},\chi _\alpha ])\},`$
where the impurity spin part has been expressed in terms of a spin coherent state path integral , $`\widehat{\mathrm{\Omega }}=(\theta ,\varphi )`$ is a unit vector describing the family of spin states $`|S_d,m_d`$, the eigenstates of $`S_d^2`$ and $`S_d^z`$ with eigenvalues $`S_d(S_d+1)`$ and $`m`$, respectively, and the periodic boundary condition is assumed for the spin vector variable.
$$iS_d\omega \left(\widehat{\mathrm{\Omega }}\right)=iS_d_0^\beta 𝑑\tau (1\mathrm{cos}\theta )\stackrel{}{\varphi }$$
(20)
is known as the Berry phase of the spin history, which is a purely geometric factor and will play no essential role here. In the path-integral representation, the reduced Hamiltonian $`H_s`$ is expressed as
$`H_s[\widehat{\mathrm{\Omega }},\chi _\alpha ]={\displaystyle \frac{v_f}{4\pi }}{\displaystyle \underset{\alpha =1}{\overset{3}{}}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\chi _\alpha (x,\tau )(i_x)\chi _\alpha (x,\tau )`$ (21)
$`{\displaystyle \frac{iJS_d}{2}}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\delta (x)\text{ }\widehat{\mathrm{\Omega }}(\tau )\stackrel{}{\chi }(x,\tau )\times \stackrel{}{\chi }(x,\tau ),`$ (22)
where $`\chi _\alpha (x,\tau )`$ ($`\alpha =1,2,3`$) are real Grassmann variables corresponding to the three Majorana fermion fields and as far as $`\chi _\alpha `$ are concerned, the path integrals over them are bilinear. The partition function can thus be rewritten in the following form
$`Z`$ $`=`$ $`{\displaystyle D\widehat{\mathrm{\Omega }}\mathrm{exp}\left(iS_d\omega \left(\widehat{\mathrm{\Omega }}\right)\right)\underset{\alpha =1}{\overset{3}{}}D\chi _\alpha }`$ (24)
$`\mathrm{exp}\left\{{\displaystyle \frac{1}{2}}{\displaystyle _0^\beta }𝑑\tau {\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑x\text{ }\mathrm{\Psi }^{}(x,\tau )\widehat{M}\text{ }\mathrm{\Psi }(x,\tau )\right\},`$
with a three-component vector $`\mathrm{\Psi }^{}(x,\tau )=(\chi _1(x,\tau ),`$ $`\chi _2(x,\tau ),`$ $`\chi _3(x,\tau ))`$ and its transposition $`\mathrm{\Psi }(x,\tau ).`$ The matrix is denoted by $`\widehat{M}=(_\tau \overline{v}_fi_x)I+iJS\delta (x)\widehat{M^{}}(\tau ),`$ where
$$\widehat{M^{}}(\tau )=\left[\begin{array}{ccc}0,& \mathrm{\Omega }^z(\tau ),& \mathrm{\Omega }^y(\tau )\\ \mathrm{\Omega }^z(\tau ),& 0,& \mathrm{\Omega }^x(\tau )\\ \mathrm{\Omega }^y(\tau ),& \mathrm{\Omega }^x(\tau ),& 0\end{array}\right],$$
(25)
$`\overline{v}_f=v_f/2\pi ,`$ and $`I`$ is a $`3\times 3`$ unit matrix. Then we integrate out the variables $`\chi _\alpha `$ and obtain an effective action which only contains the spin vector variables:
$`Z`$ $`=`$ $`Z_0{\displaystyle D\widehat{\mathrm{\Omega }}\mathrm{exp}\left(iS_d\omega \left(\widehat{\mathrm{\Omega }}\right)S_{eff}\right)},`$ (26)
$`S_{eff}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\text{Tr }\mathrm{ln}\left[1+iJS_d\delta (x)\widehat{G}\widehat{M^{}}(\tau )\right],`$ (27)
where $`Z_0=\frac{1}{2}det\left[(_\tau \overline{v}_fi_x)I\right]`$ is the partition function of the non-interacting limit of $`H_s,`$ and its free Majorana fermion propagator is given by $`\widehat{G}=(_\tau \overline{v}_fi_x)^1`$. Tracing here is taken over space, imaginary time, and the matrix indices. Using the identity
$`\text{Tr }\mathrm{ln}(1+\widehat{A})={\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(1)^n}{n}}\text{Tr }(\widehat{A})^n,`$we obtain the general expression for the effective action
$`S_{eff}`$ $`=`$ $`{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(iJS_d)^n}{2n}}{\displaystyle _0^\beta }𝑑\tau _1\mathrm{}{\displaystyle _0^\beta }𝑑\tau _nG(\tau _{12})G(\tau _{23})\mathrm{}`$ (29)
$`\text{ }G(\tau _{n1})\text{Tr}\left[\widehat{M^{}}(\tau _1)\widehat{M^{}}(\tau _2)\mathrm{}\mathrm{}\widehat{M^{}}(\tau _n)\right],`$
where the space integration has been easily carried out due to the presence of the delta function so that the free Majorana fermion propagators are replaced by the local one, $`G(\tau )=`$ $`\frac{1}{\overline{v}_f}\frac{\pi /\beta }{\mathrm{sin}(\pi \tau /\beta )}`$ , and $`\tau _{ij}=\tau _i\tau _j`$. Note that $`\tau _1,\tau _2,\mathrm{}.,\tau _n`$ are unequal imaginary times. In fact $`S_{eff}^{(n)}`$ is the contribution of a one-loop diagram made of $`n`$ local propagators $`G(\tau _{ij})`$ and $`n`$ local vertices $`iJS\widehat{M^{}}(\tau _i)`$. The factor $`1/(2n)`$ in front of each term is due to symmetry of the corresponding diagram. In the above derivation, the impurity spin is not specified to be $`1/2`$ and only the spin of the conduction electrons has been assumed to $`1/2`$. However, the following discussion will focus on the overscreening case $`\left(S_d=1/2\right)`$.
Up to second order of $`JS_d,`$ the effective action is worked out as
$$S_{eff}^{(2)}=\frac{1}{2}(JS_d)^2_0^\beta 𝑑\tau _i_0^\beta 𝑑\tau _jG^2(\tau _{ij})\widehat{\mathrm{\Omega }}(\tau _i)\widehat{\mathrm{\Omega }}(\tau _j),$$
(30)
which describes a Heisenberg spin chain with an inverse-square long-range antiferromagnetic interaction between the impurity spins at two different imaginary times, and the sign of the original Kondo exchange coupling (ferromagnetic or antiferromagnetic) is not distinguishable here. Note that the spin coherent state path integral has assumed a periodic boundary condition for the spin variable so that the Heisenberg spins in fact sit on a circle of length $`\beta `$ with exchange inversely proportional to the square of the distance between spins, which has the same form of the path-integral functional as the one-dimensional Heisenberg spin chain with an inverse-square long-range interaction, the HS model in the semiclassical limit. Therefore, all the static properties of the HS model can be readily translated to the present model. (a) The low-energy states of $`S_{eff}^{(2)}`$ in the large-$`\beta `$ limit are described by the chiral-SU(2) invariant $`k=1`$ Wess-Zumino-Witten model, which is a conformally invariant Gaussian field theory. $`S_{eff}^{(2)}`$ thus can represent a fixed point action of the two-channel Kondo model, and the elementary excitations are spinons, i.e., the $`S=1/2`$ particles instead of spin waves which are the elementary excitations of an ordered antiferromagnet. The spinons satisfy semion statistics intermediate between bosons and fermions, being an example of the exclusion statistics interpretation of fractional statistics .(b) $`S_{eff}^{(2)}`$ describes a free gas of spinons, being a fundamental model for gapless spin-$`1/2`$ antiferromagnetic spin model, and the dominant asymptotic spin-spin correlation function of $`S_{eff}^{(2)}`$ is algebraic with an universal exponent $`\eta =1`$ without logarithmic corrections .
$$<T_\tau \widehat{\mathrm{\Omega }}(\tau _i)\widehat{\mathrm{\Omega }}(\tau _j)>\frac{1}{\tau _i\tau _j}$$
(31)
and the impurity spin variable $`\widehat{\mathrm{\Omega }}(\tau )`$ thus acquires a dynamic scaling dimension $`1/2.`$ Such a behavior is also the universal spin-spin correlation function of the low-energy non-FL behavior of the two-channel Kondo model , leading to a marginal FL form of the impurity spin spectrum: Im$`\chi _d(\omega +i0^+)\mathrm{tanh}(\frac{\omega }{2T}).`$ (c). We can thus conclude that $`S_{eff}^{(2)}`$ represents the low-energy non-FL fixed point action of the spin part of the two-channel Kondo model.
The third order of $`JS_d`$ can be viewed as a correction to the fixed point action, which has been derived as
$`S_{eff}^{(3)}`$ $`=`$ $`{\displaystyle \frac{i}{6}}(JS_d)^3{\displaystyle _0^\beta }𝑑\tau _i{\displaystyle _0^\beta }𝑑\tau _j{\displaystyle _0^\beta }𝑑\tau _k`$ (33)
$`G(\tau _{ij})G(\tau _{jk})G(\tau _{ki})\widehat{\mathrm{\Omega }}(\tau _i)\left(\widehat{\mathrm{\Omega }}(\tau _j)\times \widehat{\mathrm{\Omega }}(\tau _k)\right).`$
This term describes a Heisenberg spin chain with a long-range interaction of the impurity spins at three different imaginary times and is completely antisymmetric in its indices ($`ijk`$). In accordance with the low-energy non-FL fixed point action $`S_{eff}^{(2)}`$, this interaction term is irrelevant because it has a dynamic scaling dimension $`3/2.`$ However, it is this interaction that distinguishes the sign of the original Kondo exchange coupling, so that it is a dangerous irrelevant operator. It is conceivable that the non-FL thermodynamic properties of the two-channel Kondo model around the low-energy fixed point are derived from a perturbation theory of $`S_{eff}^{(3)}`$. For instance, the second order perturbation calculation of $`S_{eff}^{(3)}`$ gives rise to the extra low-temperature specific heat due to the presence of the impurity spin as $`T\mathrm{ln}T`$ . The detailed calculations of the thermodynamic properties will be given in the future publication.
When the expansion is carried out to the fourth order, we obtain a four-spin long-range interaction of the impurity spins at four different times
$`S_{eff}^{(4)}={\displaystyle \frac{1}{4}}(JS_d)^4{\displaystyle _0^\beta }𝑑\tau _i{\displaystyle _0^\beta }𝑑\tau _j{\displaystyle _0^\beta }𝑑\tau _k{\displaystyle _0^\beta }𝑑\tau _l`$ (34)
$`\text{ }G(\tau _{ij})G(\tau _{jk})G(\tau _{kl})G(\tau _{li})\widehat{\mathrm{\Omega }}(\tau _i)\widehat{\mathrm{\Omega }}(\tau _j)\text{ }\widehat{\mathrm{\Omega }}(\tau _k)\widehat{\mathrm{\Omega }}(\tau _l),`$ (35)
which is clearly irrelevant as far as the low-energy fixed point is concerned, as its dynamic scaling dimension is $`2.`$ All higher-order terms thus contain no essential physics and can be neglected completely. Therefore, the reduced effective action will be given by $`S_{eff}=S_{eff}^{(2)}+S_{eff}^{(3)}`$, which is a natural generalization of the inverse-square long-range HS spin exchange model.
In summary, we use the Majorana fermion representation of the algebra of spin, charge, and flavor currents to transform the two-channel Kondo model. In a path integral formulation, we derive a reduced effective action of the two-channel Kondo model, which is a one-dimensional Heisenberg spin chains with two-spin, three-spin, etc. long-range interactions, as a natural generalization of the inverse-square long-range HS spin model. It is argued that the nontrivial two-channel Kondo physics in the low-energy regime can be reproduced from the first two terms of the spin action, and the other relevant issues are under investigation. As pointed out in Ref. , the infinite set of multiplicative degeneracy of HS model is due to the hidden $`SU(2)`$ Yangian symmetry. Comparing our present formulation with the CFT treatment of the two-channel Kondo problem, this statement becomes apparent. After completing the present work, we become aware of a general review article on exact results for highly correlated electron systems in one dimension, where some analogies of the inverse square long-range models to other interacting models are discussed within the framework of the Bethe-Ansatz.
Acknowledgment
One of the authors (G. -M. Zhang) would like to thank A. C. Hewson for his helpful discussions on the two-channel Kondo model. |
warning/0001/math0001015.html | ar5iv | text | # References
The Gervais-Neveu-Felder equation for the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra
A. Chakrabarti<sup>1</sup><sup>*</sup><sup>*</sup>*chakra@cpht.polytechnique.fr and R. Chakrabarti<sup>2</sup>
<sup>1</sup> Centre de Physique Théorique Laboratoire Propre du CNRS UPR A.0014, Ecole Polytechnique, 91128 Palaiseau Cedex, France
<sup>2</sup> Department of Theoretical Physics, University of Madras,
Guindy Campus, Madras 600 025, India
## Abstract
Using a contraction procedure, we construct a twist operator that satisfies a shifted cocycle condition, and leads to the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra. The corresponding universal $`_h(y)`$ matrix obeys a Gervais-Neveu-Felder equation associated with the $`U_{h;y}(sl(2))`$ algebra. For a class of representations, the dynamical Yang-Baxter equation may be expressed as a compatibility condition for the algebra of the Lax operators.
Recently a class of invertible maps between the classical $`sl(2)`$ and the non-standard Jordanian $`U_h(sl(2))`$ algebras has been obtained -. The classical and the Jordanian coalgebraic structures may be related - by the twist operators corresponding to these maps. Following the first twist leading from the classical to the Jordanian Hopf structure, it is possible to envisage a second twist leading to a quasi-Hopf quantization of the Jordanian $`U_h(sl(2))`$ algebra. By explicitly constructing the appropriate universal twist operator that satisfies a shifted cocycle condition, we here obtain the Gervais-Neveu-Felder (GNF) equation satisfied by the universal $``$ matrix of a one-parametric quasi-Hopf deformation of the $`U_h(sl(2))`$ algebra.
The GNF equation corresponding to the standard Drinfeld-Jimbo deformed $`U_q(sl(2))`$ algebra was studied in the context of Liouville field theory , quantization of Kniznik-Zamolodchikov-Bernard equation and the quantization of the Calogero-Moser model in the $`R`$ matrix formalism . The general construction of the twist operators leading to the GNF equation corresponding to the quasi-triangular standard Drinfeld-Jimbo deformed $`U_q(𝗀)`$ algebras and superalgebras were obtained in -.
For the sake of completeness, we start by enlisting the general properties of a quasi-Hopf algebra $`𝒜`$ . For all $`a𝒜`$ there exist an invertible element $`\mathrm{\Phi }𝒜𝒜𝒜`$ and the elements $`(\alpha ,\beta )𝒜`$, such that
$`(\text{id})(a)`$ $`=`$ $`\mathrm{\Phi }(\text{id})((a))\mathrm{\Phi }^1,`$
$`(\text{id}\text{id})(\mathrm{\Phi })(\text{id}\text{id})(\mathrm{\Phi })`$ $`=`$ $`(1\mathrm{\Phi })(\text{id}\text{id})(\mathrm{\Phi })(\mathrm{\Phi }1),`$
$`(\epsilon \text{id})`$ $`=`$ $`\text{id},`$
$`(\text{id}\epsilon )`$ $`=`$ $`\text{id},`$
$`{\displaystyle \underset{r}{}}S(a_r^{(1)})\alpha a_r^{(2)}`$ $`=`$ $`\epsilon (a)\alpha ,`$
$`{\displaystyle \underset{r}{}}a_r^{(1)}\beta S(a_r^{(2)})`$ $`=`$ $`\epsilon (a)\beta ,`$
$`{\displaystyle \underset{r}{}}X_r^{(1)}\beta S(X_r^{(2)})\alpha X_r^{(3)}`$ $`=`$ $`1,`$
$`{\displaystyle \underset{r}{}}S(\overline{X}_r^{(1)})\alpha \overline{X}_r^{(2)}\beta S(\overline{X}_r^{(3)})`$ $`=`$ $`1,`$ (1)
where
$$(a)=\underset{r}{}a_r^{(1)}a_r^{(2)},\mathrm{\Phi }=\underset{r}{}X_r^{(1)}X_r^{(2)}X_r^{(3)},\mathrm{\Phi }^1=\underset{r}{}\overline{X}_r^{(1)}\overline{X}_r^{(2)}\overline{X}_r^{(3)}.$$
(2)
A quasi-triangular quasi-Hopf algebra is equipped with a universal $``$ matrix satisfying
$`^{op}(a)`$ $`=`$ $`(a)^1,`$
$`(\text{id})()`$ $`=`$ $`\mathrm{\Phi }_{231}^1_{13}\mathrm{\Phi }_{213}_{12}\mathrm{\Phi }_{123}^1,`$
$`(\text{id})()`$ $`=`$ $`\mathrm{\Phi }_{312}_{13}\mathrm{\Phi }_{132}^1_{23}\mathrm{\Phi }_{123}.`$ (3)
The algebra is known as triangular if the additional relation
$$_{21}=^1$$
(4)
is satisfied. In a quasi-triangular quasi-Hopf algebra, the universal $``$ matrix satisfies quasi-Yang-Baxter equation
$$_{12}\mathrm{\Phi }_{312}_{13}\mathrm{\Phi }_{132}^1_{23}\mathrm{\Phi }_{123}=\mathrm{\Phi }_{321}_{23}\mathrm{\Phi }_{231}^1_{13}\mathrm{\Phi }_{213}_{12}.$$
(5)
An invertible twist operator $`𝒜𝒜`$ satisfying the relation
$$(\epsilon \text{id})()=1=(\text{id}\epsilon )()$$
(6)
performs a gauge transformation as follows:
$`_{}(a)`$ $`=`$ $`(a)^1,`$
$`\mathrm{\Phi }_{}`$ $`=`$ $`_{23}(\text{id})()\mathrm{\Phi }(\text{id})(^1)_{12}^1,`$
$`\alpha _{}`$ $`=`$ $`{\displaystyle \underset{r}{}}S(\overline{f}_r^{(1)})\alpha \overline{f}_r^{(2)},`$
$`\beta _{}`$ $`=`$ $`{\displaystyle \underset{r}{}}f_r^{(1)}\beta S(f_r^{(2)}),`$
$`_{}`$ $`=`$ $`_{21}^1,`$ (7)
where
$$=\underset{r}{}f_r^{(1)}f_r^{(2)},^1=\underset{r}{}\overline{f}_r^{(1)}\overline{f}_r^{(2)}.$$
(8)
The Jordanian Hopf algebra $`U_h(sl(2))`$ is generated by the elements $`(T^{\pm 1}(=e^{\pm hX}),Y,H)`$, satisfying the algebraic relations
$$[H,T^{\pm 1}]=T^{\pm 2}1,[H,Y]=\frac{1}{2}\left(Y(T+T^1)+(T+T^1)Y\right),[X,Y]=H,$$
(9)
whereas the coalgebraic properties are given by
$`(T^{\pm 1})`$ $`=`$ $`T^{\pm 1}T^{\pm 1},(Y)=YT+T^1Y,(H)=HT+T^1H,`$
$`\epsilon (T^{\pm 1})`$ $`=`$ $`1,\epsilon (Y)=\epsilon (H)=0,`$
$`S(T^{\pm 1})`$ $`=`$ $`T^1,S(Y)=TYT^1,S(H)=THT^1.`$ (10)
The universal $`_h`$ matrix of the triangular Hopf algebra $`U_h(sl(2))`$ is given in a convenient form by
$$_h=\text{exp}(hXTH)\text{exp}(hTHX).$$
(11)
An invertible nonlinear map of the generating elements of the $`U_h(sl(2))`$ algebra on the elements of the classical $`U(sl(2))`$ algebra plays a pivotal role in the present work. The map reads
$$T=\stackrel{~}{T},Y=J_{}\frac{1}{4}h^2J_+(J_0^21),H=(1+(hJ_+)^2)^{1/2}J_0,$$
(12)
where $`\stackrel{~}{T}=hJ_++(1+(hJ_+)^2)^{1/2}`$. The elements $`(J_\pm ,J_0)`$ are the generators of the classical $`sl(2)`$ algebra
$$[J_0,J_\pm ]=\pm \mathrm{\hspace{0.17em}2}J_\pm ,[J_+,J_{}]=J_0.$$
(13)
The twist operator specific to the map (12), transforming the trivial classical $`U(sl(2))`$ coproduct structure to the non-cocommuting coproduct properties (10) of the Jordanian $`U_h(sl(2))`$ algebra, has been obtained , as a series expansion in powers of $`h`$. The transforming operator between the two above-mentioned antipode maps has been obtained in a closed form.
Our present derivation of the GNF equation corresponding to the Jordanian $`U_h(sl(2))`$ algebra closely parallels the description in . These authors obtained the solutions of the GNF equation in the case of the standard Drinfeld-Jimbo deformed quasi-Hopf $`U_{q;x}(sl(2))`$ algebra by constructing the universal twist operator depending on a parameter $`x`$ :
$`(x)`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}(1)^k{\displaystyle \frac{(qq^1)^k}{[k]_q!}}x^{2k}q^{k(k+1)/2}\left[{\displaystyle \underset{l=1}{\overset{k}{}}}(11x^2q^{2l}\mathrm{\hspace{0.17em}\hspace{0.17em}1}q^{2𝒥_0})^1\right]\times `$ (14)
$`\times q^{\frac{k}{2}𝒥_0}𝒥_+^kq^{\frac{3k}{2}𝒥_0}𝒥_{}^k,`$
where $`[n]_q=\left(q^nq^n\right)/\left(qq^1\right)`$. The generators of the $`U_q(sl(2))`$ algebra satisfies the relations
$$q^{𝒥_0}𝒥_\pm q^{𝒥_0}=q^{\pm 2}𝒥_\pm ,[𝒥_+,𝒥_{}]=[𝒥_0]_q.$$
(15)
A key ingredient in our method is the contraction technique developed in , where a matrix $`G`$
$$G=E_q(\eta 𝒥_+)E_q(\eta 𝒥_+),\eta =\frac{h}{q1}$$
(16)
performs a similarity transformation on the universal $`_q`$ matrix of the $`U_q(sl(2))`$ algebra . The twisted exponential $`E_q(\chi )`$ reads
$$E_q(\chi )=\underset{n=0}{\overset{\mathrm{}}{}}\frac{\chi ^n}{[n]_q!}.$$
(17)
The transforming matrix $`G`$ is singular in the $`q1`$ limit. The transformed $`R_h^{j_1;j_2}`$ matrix for an arbitrary $`(j_1;j_2)`$ represention
$$R_h^{j_1;j_2}=\underset{q1}{lim}\left[G^1R_q^{j_1;j_2}G\right]$$
(18)
is, however, nonsingular and coincide, on account of the map (12), with the result obtained directly from the expression (11) of the universal $`_h`$ matrix. In the above contraction process the following two identities play a crucial role:
$`(E(\eta 𝒥_+))^1q^{\alpha 𝒥_0/\mathrm{\hspace{0.17em}2}}E(\eta 𝒥_+)`$ $`=`$ $`𝒯_{(\alpha )}q^{\alpha 𝒥_0/\mathrm{\hspace{0.17em}2}},`$
$`(E(\eta 𝒥_+))^1𝒥_{}E(\eta 𝒥_+)`$ $`=`$ $`{\displaystyle \frac{\eta }{qq^1}}\left(𝒯_{(1)}q^{𝒥_0}𝒯_{(1)}q^{𝒥_0}\right)+𝒥_{},`$ (19)
where $`𝒯_{(\alpha )}=(E(\eta 𝒥_+))^1E(q^\alpha \eta 𝒥_+)`$. In the $`q1`$ limit, it may be proved
$$\underset{q1}{lim}𝒯_{(\alpha )}=\stackrel{~}{T}^\alpha =T^\alpha .$$
(20)
The second equality in (20) follows from the map (12).
Using the contraction scheme discussed above we now obtain an one-parametric twist operator $`_h(y)U_h(sl(2))U_h(sl(2))`$, which satisfies a shifted cocycle condition. The twist operator $`_h(y)`$ gauge transforms à la (7) the Jordanian Hopf algebra
$`U_h(sl(2))`$ to a quasi-Hopf $`U_{h;y}(sl(2))`$ algebra and the transformed universal $`_h(y)`$ matrix satisfies the corresponding GNF equation. To this end we first compute
$$\stackrel{~}{}(y)=\underset{q1}{lim}(G^1(x)G)_{x^2=y(q1)},$$
(21)
where $`(x)`$ is given by (14).
A new feature here is the reparametrization described by
$$y=\frac{x^2}{q1},$$
(22)
which is necessary for obtaining nonsingular result in the $`q1`$ limit. In (22) we assume that $`x0`$ in the $`q1`$ in such a way that $`y`$ remains finite. Following the above procedure in the said limit we obtain
$$\stackrel{~}{}(y)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{(hy)^k}{k!}\left(\stackrel{~}{T}J_+\right)^k\left(\stackrel{~}{T}^3(\stackrel{~}{T}\stackrel{~}{T}^1)\right)^k.$$
(23)
The rhs of (23) is interpreted on account of the map (12) as an element of $`U_h(sl(2))U_h(sl(2))`$. Identifying this in the above sense with the twist operator $`_h(y)\left(=\stackrel{~}{}(y)\right)`$ we now obtain the crucial result
$$_h(y)=\text{exp}\left(\frac{y}{2}(1T^2)(T^2T^4)\right).$$
(24)
The above twist operator $`_h(y)`$ satisfies the property (6). Following the arguments in we express $`_h(y)`$ as a shifted coboundary
$$_h(y)=((y))\left(1^1(y)\right)\left(^1(yT_{(2)}^4)\mathrm{\hspace{0.17em}1}\right),$$
(25)
where the expression for the boundary reads
$$(y)=\text{exp}\left(\frac{y}{2}(1T^2)\right).$$
(26)
The operator $`_h(y)`$ given by (24) satisfies the following shifted cocycle condition
$$(1_h(y))\left[(\text{id})_h(y)\right]=(_h\left(yT_{(3)}^4\right)\mathrm{\hspace{0.17em}1})\left[(\text{id})_h(y)\right].$$
(27)
Following (7) the transformed coproduct property may now be read as
$$_y(a)=_h(y)(a)_h^1(y)\text{for all}aU_{h;y}(sl(2)).$$
(28)
It may now be shown that the shifted cocycle condition is a consequence of the following shifted coassociativity property:
$$(\text{id}_y)_y(a)=(_{yT_{(3)}^4}\text{id})_y(a).$$
(29)
Following (7) the gauge-transformed universal $`_h(y)`$ matrix for the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra reads
$$_h(y)=_{h\mathrm{\hspace{0.25em}\hspace{0.17em}21}}(y)_h_h^1(y).$$
(30)
The coassociator $`\mathrm{\Phi }(y)`$ corresponding to the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra may be obtained for the above constuction of the twist operator obeying the shifted cocycle condition (27). Using (7), (24) and (27) we obtain
$`\mathrm{\Phi }(y)`$ $`=`$ $`_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(yT_{(3)}^4)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}^1(y)`$ (31)
$`=`$ $`\text{exp}\left[{\displaystyle \frac{y}{2}}\left(1T^2\right)\left(T^2T^4\right)\left(1T^4\right)\right].`$
The elements $`\alpha (y)`$ and $`\beta (y)`$, characterizing the antipode map of the $`U_{h;y}(sl(2))`$ algebra may be similarly obtained from (7), (10) and (24):
$$\alpha (y)=\text{exp}\left[\frac{y}{2}\left(1T^2\right)^2\right],\beta (y)=\text{exp}\left[\frac{y}{2}\left(1T^2\right)^2\right].$$
(32)
Using the guage transformation property of the universal $``$ matrix in (7) and our construction (24) of the twist operator, we now discuss the GNF equation associated with the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra. The relations (7), (24) and (31) lead to the transformation property
$$_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}\left(yT_{(3)}^4\right)=\mathrm{\Phi }_{213}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(y)\mathrm{\Phi }_{123}^1(y).$$
(33)
Now the quasitriangularity property of $`U_{h;y}(sl(2))`$ algebra implies via (3), (31) and (33) the following relations:
$`(\text{id}_y)_h(y)=_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}^1\left(yT_{(1)}^4\right)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}\left(yT_{(3)}^4\right),`$
$`(_y\text{id})_h(y)=_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}\left(yT_{(2)}^4\right)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}\left(yT_{(3)}^4\right)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}^1(y).`$ (34)
Using the transformation property (33) we may now recast the quasi Yang-Baxter equation (5) as the GNF equation associated with the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra:
$$_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}\left(yT_{(2)}^4\right)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}(y)=_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}\left(yT_{(1)}^4\right)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}(y)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}\left(yT_{(3)}^4\right).$$
(35)
We now briefly consider the solutions of the above GNF equation (35). Using the universal $`_h(y)`$ matrix (30), the twist operator $`_h(y)`$ in (24) and the map (12) of the generators of the $`U_h(sl(2))`$ algebra on the corresponding classical elements, we may construct solutions of the GNF equation (35). As illusrations we describe the representions $`R_h(y)`$ for the $`\frac{1}{2}j`$ and the $`1j`$ cases. A $`(2j+1)`$ dimensional representation of the classical $`sl(2)`$ algebra (13)
$`J_+|jm`$ $`=`$ $`(jm)(j+m+1)|jm+1,J_{}|jm=|jm1,`$
$`J_0|jm`$ $`=`$ $`m|jm,`$ (36)
now, via the map (12), immediately furnishes the corresponding $`(2j+1)`$ dimensional representation of the $`U_h(sl(2))`$ algebra (9). For the $`j=\frac{1}{2}`$ case, the generators remain undeformed. For the $`j=1`$ case, we list the representation of $`U_h(sl(2))`$ below.
$`(j=1)`$ (47)
$`X=\left(\begin{array}{ccc}0& 2& 0\\ 0& 0& 2\\ 0& 0& 0\end{array}\right),Y=\left(\begin{array}{ccc}0& {\scriptscriptstyle \frac{1}{2}}h^2& 0\\ 1& 0& {\scriptscriptstyle \frac{3}{2}}h^2\\ 0& 1& 0\end{array}\right),`$
$`H=\left(\begin{array}{ccc}2& 0& 4h^2\\ 0& 0& 0\\ 0& 0& 2\end{array}\right).`$
Using the above representations in the expression (30) of the universal $`_h(y)`$ matrix,we obtain
$$R_h^{\frac{1}{2};j}(y)=\left(\begin{array}{cc}T& hH+\frac{1}{2}h(TT^1)\left(1+2y(1T^4)\right)\\ 0& T^1\end{array}\right)$$
(48)
and
$$R_h^{1;j}(y)=\left(\begin{array}{ccc}T^2& A& B\\ 0& 1& C\\ 0& 0& T^2\end{array}\right),$$
(49)
where
$`A`$ $`=`$ $`2hTH2hy\left(1T^2\right)\left(1T^4\right),`$
$`B`$ $`=`$ $`2h^2\left[T^2T^22TH\left(1T^2\right)(TH)^2T^2\right]4h^2y\left(1T^2\right)\left(1+4T^2T^4\right)`$
$`4h^2yTH\left(1T^2\right)\left(T^2T^2\right)+2h^2y^2\left(TT^1\right)^2\left(1T^4\right)^2,`$
$`C`$ $`=`$ $`2h\left(1T^2+THT^2\right)+2hy\left(1T^2\right)\left(T^2T^2\right).`$ (50)
From (48) it follows that the $`R_h^{\frac{1}{2};\frac{1}{2}}`$ matrix for the fundamental $`(1/2;1/2)`$ case does not depend on the parameter $`y`$. The $`R_h(y)`$ matrices for the higher representations, however, nontrivially depend on $`y`$. The $`R_h(y)`$ matrices satisfy an “exchange symmetry” between the two sectors of the tensor product spaces:
$$\left(R_h^{j_1;j_2}(y)\right)_{km,ln}=\left(R_h^{j_2;j_1}(y)\right)_{mk,nl}.$$
(51)
In the remaining part of the present work we recast the Jordanian GNF equation (35) as a compatibility condition for the algebra of $`L`$ operators. Using a new parametrization $`y=\text{exp}(z)`$, we perform a translation
$$_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(z)_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(z2hX_{(3)})$$
(52)
to express (35) in a symmetric form
$`_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(z2hX_{(3)})_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}(z+2hX_{(2)})_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}(z2hX_{(1)})`$
$`=_{h\mathrm{\hspace{0.25em}\hspace{0.17em}23}}(z+2hX_{(1)})_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}(z2hX_{(2)})_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}(z+2hX_{(3)}).`$ (53)
This is equivalent to the Jordanian GNF equation (35) for the class of representations $`\varrho _{j_1;j_2}`$ satisfying the property
$$\varrho _{j_1;j_2}\left([\left(X_{(k)}+X_{(l)}\right)_z,_{hkl}(z)]\right)=0.$$
(54)
Adopting the procedure in we here use the following construction of the Lax operator for the $`U_{h;y}(sl(2))`$ algebra
$$L_{13}(z)=\text{exp}\left[2h\left(2X_{(1)}+X_{(3)}\right)_z\right]_{h\mathrm{\hspace{0.25em}\hspace{0.17em}13}}(z)\text{exp}\left[2hX_{(3)}_z\right],$$
(55)
where the subscript $`3`$ denotes the quantum space. For the representations satisfying (54) the relation (53) may be expressed in a Lax martix form
$$R_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}^{j_1;j_2}(z2hX_{(3)})L_{13}(z)L_{23}(z)=L_{23}(z)L_{13}(z)R_{h\mathrm{\hspace{0.25em}\hspace{0.17em}12}}^{j_1;j_2}(z+2hX_{(3)}).$$
(56)
As illustrations we note that the representations $`R_h^{\frac{1}{2};\mathrm{\hspace{0.17em}1}}(z)`$, $`R_h^{1;\frac{1}{2}}(z)`$ and $`R_h^{1;\mathrm{\hspace{0.17em}1}}(z)`$ obtained from (48) and (49) satisfy the requirement (54).
To summarize, here we have constructed the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra by explicitly obtaining the relevant twist operator via a contraction method. In the contraction method used here we start with the standard Drinfeld-Jimbo deformed quasi-Hopf $`U_{q;x}(sl(2))`$ algebra and use a suitable similarity transformation followed by a $`q1`$ limiting process. An important point here is that the reparametrization as obtained in (22) is essential for obtaining a nonsingular twist operator for the $`U_{h;y}(sl(2))`$ algebra in the $`q1`$ limit. Our contraction method has an advantage in that it furnishes the dynamical quantities for the Jordanian quasi-Hopf $`U_{h;y}(sl(2))`$ algebra from the corresponding quantities of the standard Drinfeld-Jimbo deformed quasi-Hopf $`U_{q;x}(sl(2))`$ algebra. The present twist operator associated with the $`U_{h;y}(sl(2))`$ algebra satisfies a shifted cocycle condition. The universal $`_h(y)`$ matrix satisfies the GNF equation associated with the $`U_{h;y}(sl(2))`$ algebra. For a special class of representations, the GNF equation may be recast as a compatibility condition of the $`L`$ operators. As an extension of the present work, a similar formalism may be developed to describe a quasi-Hopf quantization of the coloured Jordanian deformed $`gl(2)`$ algebra considered in , , . A similar construction of the twist operators associated with the quasi-Hopf deformation of the Jordanian $`sl_h(N)`$ algebra may also be attempted following the discussion in .
Acknowledgments:
One of us (RC) wishes to thank A. J. Bracken for a kind invitation to the University of Queensland, where part of this work was done. |
warning/0001/cond-mat0001460.html | ar5iv | text | # Counting Statistics of an Adiabatic Pump
Adiabatic charge pumping has attracted considerable theoretical and experimental interest. It occurs when the Hamiltonian of the system is changed periodically with time. At the end of the pumping cycle a finite charge may be transmitted through the system. Such a charge transfer takes place even in the absence of any dc voltage applied to the system. The idea is originally due to Thouless, who showed that in certain one-dimensional systems the transmitted charge is quantized in the adiabatic limit. This mechanism has been proposed for making electric current standards .
Recently research has focussed on adiabatic pumping through mesoscopic devices . Most of the work concentrated on the average pumping current and its mesoscopic fluctuations . The general expression for the average transmitted charge was derived by Brouwer in terms of the time–dependent $`S`$-matrix based on the formalism of Ref. . The average charge transmitted during a cycle depends on the system’s conductance and is not in general quantized . The absence of quantization is related to the fact that the true Thouless adiabatic conditions can never be achieved for a system with infinite leads and hence a vanishing excitation gap. The possibility to create electron-hole pairs with arbitrarily small energy leads to dissipation and violates the exact quantization . Taken together with the discrete nature of charge carriers this fact suggests the presence of thermal and shot noise in the adiabatic pumps. A better understanding of the noise characteristics of such devices is certainly necessary if the prospect of their use as current standards is to be taken seriously.
In this Letter we address the problem of charge counting statistics in an adiabatic pump. We have greatly benefited from the extensive work on the current statistics of voltage-biased resistors by Levitov et. al. .
We consider a scatterer connected with the left and right leads, each having $`n`$ transverse channels . Such a scatterer is characterized by a $`2n\times 2n`$ unitary scattering matrix $`S`$, which can be written in the $`n\times n`$ block form
$$S=\left(\begin{array}{cc}r\hfill & t^{}\hfill \\ t\hfill & r^{}\hfill \end{array}\right).$$
(1)
Here $`r`$ ($`r^{}`$) and $`t`$ ($`t^{}`$) are left (right) reflection and transmission matrices correspondingly. We shall assume hereafter that the $`S`$-matrix is a periodic function of time $`\tau `$, with the period $`\tau _0`$, $`S(\tau +\tau _0)=S(\tau )`$. This time–dependence is supposed to be adiabatic, i.e. $`\tau _0`$ is much greater than the Wigner delay time of the device. We restrict our attention to open systems, i.e. those with the dwell time shorter than the inverse mean level spacing in the device. Under these conditions one may neglect both the energy relaxation and the Coulomb blockade effects inside the device.
The transmission and reflection matrices in Eq. (1) can be simultaneously diagonalized by the block-diagonal unitary matrices $`U(\tau )`$ and $`V(\tau )`$ (see for example Ref. )
$$S(\tau )=U(\tau )\stackrel{~}{S}(\tau )V^{}(\tau ).$$
(2)
Here $`\stackrel{~}{S}(\tau )`$ is a matrix of the form Eq. (1) with real diagonal reflection and transmission blocks. The eigenvalues of the latter play the role of the instantaneous reflection and transmission amplitudes in the channels. The ambiguity in the definition of matrices $`U(\tau )`$ and $`V(\tau )`$ does not affect our final results.
In each cycle a number of electrons $`Q`$ may pass through the scattering region in a direction which depends on the detailed form of $`S(\tau )`$. The average charge transmitted in $`N`$ pumping cycles was recently given by Brouwer
$$Q=\frac{1}{2i}\underset{0}{\overset{N\tau _0}{}}\frac{d\tau }{2\pi }\text{Tr}\left\{\frac{S}{\tau }S^{}\sigma _3\right\},$$
(3)
where charge $`Q`$ is measured in units of the electron charge. We have used the representation where the transmitted charge is written as half the sum of that through the left and through the right leads , thus $`\sigma _3=\mathrm{diag}\{1,1\}`$ is a Pauli matrix with the $`n\times n`$ block structure. Both quantum and thermal noise lead to fluctuations in the transmitted charge $`Q`$. As a result the charge transmitted in one cycle can be described by a certain probability distribution.
Here we obtain the distribution function of transmitted charge for an adiabatic pump characterized by a periodic time-dependent $`S`$-matrix, $`S(\tau )`$. More precisely, we calculate the probability $`P_N(Q)`$ to transmit the charge $`Q`$ upon completion of $`N`$ pumping cycles. It is convenient to formulate the results in terms of the generating function $`F_N(\lambda )`$ of the moments of the transmitted charge defined as
$$F_N(\lambda )=e^{iQ\lambda }=𝑑QP_N(Q)e^{iQ\lambda }.$$
(4)
For the $`S`$-matrix of the form Eq. (2) we obtain
$$F_N(\lambda )=e^{i\widehat{N}\lambda }\text{det}\left[1+\stackrel{~}{S}_\lambda (\tau )\stackrel{~}{n}(\tau ,\tau ^{})(\stackrel{~}{S}_\lambda ^{}(\tau ^{})\stackrel{~}{S}_\lambda ^{}(\tau ^{}))\right],$$
(5)
where
$$\stackrel{~}{S}_\lambda (\tau )\mathrm{exp}\{i\sigma _3\lambda /4\}\stackrel{~}{S}(\tau )\mathrm{exp}\{i\sigma _3\lambda /4\}.$$
(6)
The matrix $`\stackrel{~}{n}(\tau ,\tau ^{})`$ is defined as
$$\stackrel{~}{n}(\tau ,\tau ^{})=V^{}(\tau )\widehat{n}(\tau \tau ^{})V(\tau ^{}).$$
(7)
Here the diagonal matrix $`\widehat{n}(\tau \tau ^{})`$ is the time Fourier transform of $`\widehat{n}(ϵ)=\text{diag}\{n_L(ϵ),n_R(ϵ)\}`$, where $`n_{L(R)}(ϵ)`$ is the energy distribution function of the left (right) lead. The operator in the determinant in Eq. (5) should be understood as an operator in the time space as well as a matrix in the space of channels. Finally, the integer number $`\widehat{N}`$ defined as
$$\widehat{N}=\frac{1}{2i}\underset{0}{\overset{N\tau _0}{}}\frac{d\tau }{2\pi }\text{Tr}\left\{U^{}\frac{U}{\tau }\sigma _3V^{}\frac{V}{\tau }\sigma _3\right\},$$
(8)
is the contribution to the generating function arising from the chiral anomaly (see below). Equation (5) is the main result of this Letter. Before presenting the derivation, we shall illustrate its applications on a few examples.
Expanding $`\mathrm{ln}F_N(\lambda )`$ to the first power of $`i\lambda `$ one finds for the average transmitted charge, $`Q_QQP_N(Q)`$,
$$Q=\underset{0}{\overset{N\tau _0}{}}\frac{d\tau }{4\pi }\text{Tr}\left\{\stackrel{~}{S}(\tau )\stackrel{~}{n}(\tau ,\tau ^{})[\sigma _3,\stackrel{~}{S}^{}(\tau ^{})]\right\}|_{\tau ^{}\tau }+\widehat{N}.$$
(9)
In the absence of an external voltage $`\stackrel{~}{n}(\tau ,\tau ^{})i/[2\pi (\tau \tau ^{}+i\eta )]+i/(2\pi )V^{}(\tau )V(\tau )/\tau `$ for $`\tau ^{}\tau `$. Expanding $`\stackrel{~}{S}^{}(\tau ^{})`$ to the first power in $`\tau \tau ^{}`$ and using Eq. (8) one obtains Brouwer’s result , Eq. (3).
A particularly instructive example is given by a single channel case ($`n=1`$) with the $`S`$-matrix depending on time as
$$S(\tau )=e^{i\sigma _3\theta (\tau )/2}\stackrel{~}{S}e^{i\sigma _3\theta (\tau )/2},$$
(10)
where $`\theta (\tau +\tau _0)\theta (\tau )=2\pi `$, and $`\stackrel{~}{S}`$ is time-independent. This form of the $`S`$-matrix is equivalent to the phase winding of the reflection amplitudes, $`rr\mathrm{exp}\{i\theta (\tau )\}`$ ($`r^{}r^{}\mathrm{exp}\{i\theta (\tau )\}`$). The contribution from the chiral anomaly in Eq. (8) coincides with the number of cycles, $`\widehat{N}=N`$. We concentrate first on a particularly simple time–dependence, $`\theta (\tau )=2\pi \tau /\tau _0`$. In this case the determinant in Eq. (5) may be easily calculated in the Fourier basis leading to
$`F_N(\lambda )`$ $`=`$ $`{\displaystyle \underset{k}{}}[1+n_L(ϵ_{k+N})(1n_R(ϵ_k))(e^{i\lambda }1)|t|^2`$ (11)
$`+`$ $`(1n_L(ϵ_{k+N}))n_R(ϵ_k)(e^{i\lambda }1)|t|^2]e^{iN\lambda },`$ (12)
where $`ϵ_k=\pi (2k+1)/(N\tau _0)`$ are fermionic frequencies. At small temperature, $`T1/\tau _0,V`$, this expression simplifies substantially, leading to
$$F_N(\lambda )=(e^{i\lambda }|r|^2+|t|^2)^N(|r|^2+e^{i\lambda }|t|^2)^{\frac{N\tau _0V}{2\pi }},$$
(13)
where $`V`$ is a voltage applied between the left and right leads, which is chosen to be such that $`N\tau _0V/(2\pi )`$ is an integer . In the absence of external voltage Eqs. (4) and (13) lead to the binomial distribution function of the charge transmitted through the adiabatic pump
$$P_N(Q)=C_N^Q|r|^{2Q}(1|r|^2)^{NQ}.$$
(14)
Note that only the integer values of the transmitted charge have non–zero probability to be detected. The physical meaning of this expression is that each pumping cycle is associated with an attempt to transfer one electron. The success probability of such an attempt is given by the reflection probability, $`|r|^2`$, whereas the probability of failure is $`1|r|^2`$. The above statistics should be compared with the case of the dc voltage applied across the scattering region (no pumping). This case may be obtained from Eq. (13) in the limit $`N0`$, whereas $`N\tau _0V/(2\pi )\stackrel{~}{N}`$ – integer. We immediately recover the familiar results
$$P_{\stackrel{~}{N}}(Q)=C_{\stackrel{~}{N}}^Q|t|^{2Q}(1|t|^2)^{\stackrel{~}{N}Q}.$$
(15)
The distribution is also binomial, however the probability of success is given by the transmission probability, $`|t|^2`$. The second cumulant of the transmitted charge coincide for the adiabatic pump and applied dc voltage and is given by
$$Q^2=|t|^2(1|t|^2)N,$$
(16)
leading to the maximal noise power in both cases at $`|t|^2=|r|^2=1/2`$.
The binomial distribution, Eq. (14) was derived above for the simplest time dependence of the form $`r(\tau )=r\mathrm{exp}\{2\pi i\tau /\tau _0\}`$. The same result may be shown to be valid for a more general class of pumping strategies, which we call the coherent pumping, following the terminology of Ref. :
$$e^{i\theta (\tau )}=\frac{e^{2\pi i\tau /\tau _0}z}{1z^{}e^{2\pi i\tau /\tau _0}},$$
(17)
where $`z`$ is a complex number with $`|z|<1`$. Indeed, at zero temperature $`n(ϵ_k)=\theta (2k1)`$ has infinitely degenerate eigenvalues. One can therefore diagonalize the operator in Eq. (5), in the $`z`$-dependent basis which does not mix positive and negative frequencies. The value of the determinant does not depend on $`z`$ and one recovers the binomial distribution (14). Such coherent pumping strategy guarantees the minimal possible value of the noise in the transmitted charge. The general expression for the noise of the adiabatic pump may be derived by expanding $`\mathrm{ln}F_N(\lambda )`$, Eq. (5), to the second order in $`i\lambda `$ and is given by
$$Q^2=\underset{0}{\overset{N\tau _0}{}}\frac{d\tau d\tau ^{}}{(4N\tau _0)^2}\frac{\text{Tr}\left\{1(S^{}\sigma _3S)_\tau (S^{}\sigma _3S)_\tau ^{}\right\}}{\mathrm{sin}^2[\pi (\tau \tau ^{})/(N\tau _0)]}.$$
(18)
Substituting the $`S`$–matrix of the form Eq. (10) and minimizing the charge fluctuations, $`Q^2`$, with respect to $`e^{i\theta (t)}`$ , one finds that the coherent pumping, Eqs. (10), (17), lead to the minimal possible noise. The value of this minimal noise is given by Eq. (16). Note also that the second moment, Eq. (18) (as well as the first one, Eq. (3), and the higher ones) may be expressed through the $`S`$–matrix only, rather than through auxiliary matrices defined in Eq. (2).
Next we consider a $`2\times 2`$ scattering matrix with real reflection and transmission amplitudes given by $`r=r^{}=\mathrm{cos}(2\pi \tau /\tau _0)`$ and $`t=t^{}=\mathrm{sin}(2\pi \tau /\tau _0)`$ respectively. In this case $`U(\tau )=V(\tau )=1`$ leading to $`\widehat{N}=0`$, and $`\stackrel{~}{n}`$ is the equilibrium distribution function. One can show that $`F_N(\lambda )`$ in Eq. (5) is an even function of $`\lambda `$, and is therefore real. As a result one may employ the method of Ref. to compute the determinant of the operator in Eq. (5): one multiplies this operator by its Hermitian conjugate and takes the square root. The resulting operator may be written as $`1+(1\stackrel{~}{n})\stackrel{~}{S}_\lambda ^{}\stackrel{~}{S}_\lambda \stackrel{~}{n}+\stackrel{~}{n}S_\lambda ^{}\stackrel{~}{S}_\lambda (1\stackrel{~}{n})`$. In the energy representation it has a finite number of the off-diagonal matrix elements and its determinant can be straightforwardly evaluated. This way one obtains
$$F_N(\lambda )=\left(\frac{1+\mathrm{cos}\lambda }{2}\right)^N.$$
(19)
Therefore such pump is a realization of the “random walk motion” for charge. Indeed, there are three possible values of the transmitted charge in each cycle: $`Q=0`$ with the probability $`1/2`$, and $`Q=\pm 1`$ with the probability $`1/4`$ each.
We note that the logarithmic derivative $`iV^{}(\tau )V(\tau )/\tau `$ is analogous to the instantaneous matrix of “voltages” applied to the incoming channels. The integral of this quantity can be interpreted as the number of transmission attempts . The different outcomes of such attempts lead to the noise of the pumping current. In general the probability distribution of the transmitted charge is not binomial. If the “voltage” matrix can not be diagonalized simultaneously with the reflection and transmission matrices then the distribution function of the transmitted charge does not factorize into binomial distributions of elementary transmission processes.
In contrast, the matrix $`U(\tau )`$ corresponds to the outgoing channels and enters the final expression (5) only through the chiral anomaly term (8) and therefore contributes to the average current but not to the noise. For example, the pumping cycle of the form $`S(\tau )=U(\tau )\stackrel{~}{S}`$ at zero temperature would produce a noiseless quantized pumping current.
We turn now to the derivation of Eq. (5). To this end we model the leads by a $`2n`$-component vector of chiral incoming fermions $`(\psi _L(x,\tau ),\psi _R(x,\tau ))`$ and $`2n`$-component vector of chiral outgoing fermions $`(\xi _L(x,\tau ),\xi _R(x,\tau ))`$. The action for e.g. left lead is written as
$$S_L=\underset{𝒞}{}𝑑\tau \underset{\mathrm{}}{\overset{0}{}}𝑑x\overline{\psi }_L(_t+\widehat{v}_L_x)\psi _L+\overline{\xi }_L(_t\widehat{v}_L_x)\xi _L,$$
(20)
where $`\widehat{v}_L`$ is a diagonal $`n\times n`$ matrix of the left lead channel velocities. In this expression the time integral runs along the Keldysh contour, $`𝒞`$, from $`\tau =0`$ to $`\tau =N\tau _0`$ and then back to $`\tau =0`$. The right lead is described by the similar action with the space integral running from $`x=0`$ to $`x=+\mathrm{}`$, and the velocity matrix $`\widehat{v}_R`$. Finally the incoming and outgoing channels at $`x=0`$ are related by the time–dependent $`S`$-matrix operator
$$\xi (0,\tau )=\widehat{v}^{1/2}S(\tau )\widehat{v}^{1/2}\psi (0,\tau ).$$
(21)
The current operator has a form $`I=(I_L+I_R)/2`$, where
$$I_L(\tau )=\left[\overline{\psi }_L(0^{},\tau )\widehat{v}_L\psi _L(0^{},\tau )\overline{\xi }_L(0^{},\tau )\widehat{v}_L\xi _L(0^{},\tau )\right].$$
(22)
The operator of the charge transmitted in $`N`$ cycles is given by $`Q=\underset{0}{\overset{N\tau _0}{}}𝑑\tau I(\tau )`$. Finally, the generating function may be written as
$$F_N(\lambda )=D[\psi ,\xi ]e^{S_LS_R+\frac{i}{2}\underset{𝒞}{}𝑑\tau \widehat{\lambda }(\tau )I(\tau )},$$
(23)
where $`\widehat{\lambda }(\tau )`$ is equal to $`\lambda `$ on the forward and $`\lambda `$ on the backward part of the Keldysh contour. The fermion fields in this integral obey the boundary condition, Eq. (21). One has to specify the initial, $`\tau =0`$, density matrix, which implicitly defines the Green functions. We fix the occupation numbers in the incoming channels of the left and right leads to be $`n_L(ϵ)`$ and $`n_R(ϵ)`$ correspondingly, whereas the outgoing channels are supposed to be initially empty in accord with the scattering setup.
The subsequent calculations amount to the evaluation of the Gaussian integral in Eq. (23). To this end we first make the chiral gauge transformation of the fermionic fields: $`\psi (x,\tau )V(\tau )\psi (x,\tau )`$ and $`\xi (x,\tau )U(t)\xi (x,\tau )`$. As a result, the boundary condition for the new fermions contains the $`\stackrel{~}{S}(\tau )`$ matrix only and the action acquires an additional time–dependent (matrix) chemical potential term
$$\delta S=\underset{𝒞}{}𝑑\tau 𝑑x\overline{\psi }\left[V^{}\frac{V}{\tau }\right]\psi +\overline{\xi }\left[U^{}\frac{U}{\tau }\right]\xi .$$
(24)
Such potential term results in the redefinition of the density matrix according to Eq. (7). Importantly, upon the chiral gauge transformation the expression for the current acquires an extra term $`II+1/(4\pi i)\mathrm{Tr}\left\{U^{}(\tau )\sigma _3U(\tau )/\tau V^{}(\tau )\sigma _3V(\tau )/\tau \right\}`$ arising from the chiral anomaly .
Since the source field, $`\widehat{\lambda }(\tau )`$, is a constant on the both branches of the Keldysh contour, one may eliminate the $`\widehat{\lambda }I`$ term from the action by the time–independent gauge transformation , e.g. $`\psi _Le^{i\theta (x0^{})\lambda /2}\psi _L`$ on the forward branch of the contour and $`\psi _Le^{i\theta (x0^{})\lambda /2}\psi _L`$ on the backward branch. Such transformation leads to the change in the phase of the forward scattering amplitude and can be taken into account by a redefinition of the $`\stackrel{~}{S}`$–matrix in the boundary condition, Eq. (21), $`\stackrel{~}{S}\stackrel{~}{S}_{\pm \lambda }`$ on the forward (backward) branches, with $`\stackrel{~}{S}_\lambda `$ defined in Eq. (6).
The subsequent steps are straightforward. One integrates out all degrees of freedom except for those which reside directly at the scatterer, $`x=0`$. Using the boundary condition, Eq. (21), with $`\stackrel{~}{S}_{\pm \lambda }`$–matrix, one eliminates the incoming degrees of freedom, $`\psi (x=0,\tau )`$. The remaining Gaussian integral over the outgoing fermions, $`\xi (x=0,\tau )`$, can be straightforwardly evaluated, resulting in the determinant written in Eq. (5). The remaining term, $`\mathrm{exp}\{i\widehat{N}\lambda \}`$, is the contribution from the chiral anomaly, as explained above.
To conclude, we have derived a general expression for the counting statistics of the charge transmitted through a system described by a time–dependent $`S`$–matrix. The only limitations of our result are the requirements of adiabaticity and the absence of inelastic processes in the scattering region. The absolute minimum of the noise power may be achieved by the coherent pumping strategy, in which case the charge distribution is given by the product of binomial distributions. We point out the major role played by the chiral anomaly contribution to the average transmitted charge. Such anomalous term did not arise in the context of voltage–biased systems , but is extremely important for the adiabatic pumping setup.
It is our pleasure to acknowledge helpful discussions with I. Aleiner, Y. Avron, B. Spivak and L. Sadun. We appreciate the warm hospitality of the Norwegian Centre for Advanced Studies, where part of this work was performed. This work was partly supported through the grants BSF-9800338 and DMR-9984002. A. A. is an A. P. Sloan and Packard Fellow. |
warning/0001/hep-ph0001205.html | ar5iv | text | # The Ginzburg regime and its effects on topological defect formation
## I Introduction
The Ginzburg temperature $`T_G`$ was thought historically to be the determining energy scale at which topological defects are formed in the aftermath of a second order symmetry breaking phase transition . More recently theoretical and experimental progress has pointed in the direction that it is the critical dynamics of the second order transition, induced by a change in some external parameter such as temperature or pressure, that determines the numbers (and configuration) of topological defects emerging below the critical point .
Nevertheless the role of large thermal fluctuations within the Ginzburg regime in defect formation mechanisms remains relatively poorly understood. In particular it is not clear how a density of defects created, presumably by the critical dynamics of the system, could evade alteration when exposed extensively to thermal fluctuations in the Ginzburg regime.
This issue has been rekindled recently due to the possibility of new experimental tests and, in particular, by the negative results of a pressure quench experiment in $`{}_{}{}^{4}He`$ , a system in which (because of strong interactions) the Ginzburg regime is particularly extensive. This experiment has improved on an apparatus used earlier by McClintock et al. to implement a superfluid transition in $`{}_{}{}^{4}He`$ through a sudden pressure quench. The results show no evidence for the formation of topological defects at the anticipated levels, contrary to expectations based both on the old experiment , the theory<sup>*</sup><sup>*</sup>*Although a factor $`f\stackrel{>}{}10`$ in the formula for the string density $`n1/(f\widehat{\xi })^2`$ could explain the new results and seems consistent with recent numerical studies . and the $`{}_{}{}^{3}He`$ data .
The discrepancy with the earlier $`{}_{}{}^{4}He`$ quench data is now seen as the evidence of mechanical stirring in the first version of the experiment. Nevertheless to address this discrepancy with $`{}_{}{}^{3}He`$ it was suggested that because the Ginzburg regime in $`{}_{}{}^{4}He`$ extends over a broad range of temperatures around the $`\lambda `$-line, large scale fluctuations may be able to unwind and alter the configuration of the order parameter (in contrast to $`{}_{}{}^{3}He`$) while the quench proceeds. The Ginzburg temperature is defined, somewhat qualitatively, through the loss of ability of the order parameter to hop, through thermal activation, over the potential barrier between broken symmetry minima. Thus one might worry with Karra and Rivers that when the defect densities are eventually measured, at a much later time, little or no string would have survived unwinding through thermal activation.
In in this paper we study in detail the role of the Ginzburg regime in vortex string formation. In section II we discuss our model and its properties. We show in particular that it transcends the more usual time dependent Ginzburg-Landau (TDGL) dynamics in generality and reduces to it in particular cases. In section III we describe the traditional arguments for the existence of a well defined Ginzburg temperature and critically analyze their underlying assumptions in the light of known results on the thermodynamics of vortex strings. We also establish a quantitative definition of the Ginzburg temperature and discuss its uncertainties. In section IV we investigate the role of the Ginzburg regime in the dynamics of defect formation. This is achieved by exposing field configurations created at criticality to intermediate temperatures within the Ginzburg regime and analyze the effect upon the final density of defects measured at late times. We also study the memory of the order parameter when subjected to reheating. This constitutes a direct test on the theory of defect formation as a consequence of the critical dynamics of the theory. Finally we draw our conclusions and discuss in the light of our results the possible relevance of the Ginzburg regime in explaining recent experimental results in $`{}_{}{}^{4}He`$ pressure quench experiments.
## II Langevin and Fokker-Planck Field dynamics
As a working model we consider a $`U(1)`$ symmetric $`\lambda \varphi ^4`$ global field theory in 3 spatial dimensions (3D), in the canonical ensemble, i.e. in contact with a heat bath at a given temperature $`T`$. This model has the advantage of having been extensively studied in thermal equilibrium and moreover of describing the thermodynamics of $`{}_{}{}^{4}He`$ at criticality by permitting the computation of relevant critical exponents with extreme accuracy.
As such it supplies us with a controlled realistic environment in which the role of thermal fluctuations within the Ginzburg regime in changing the density of topological defects can be studied. The equations of motion for the fields are
$`\left[_t^2+\eta _t^2m^2\right]\varphi _i(x)`$ (1)
$`+\lambda \left({\displaystyle \underset{j=1}{\overset{2}{}}}\varphi _j^2(x)1\right)\varphi _i(x)=\xi _i(x,t),`$ (2)
$`\xi _i(x,t)=0,`$ (3)
$`\xi _i(x,t)\xi _j(x^{},t^{})=\mathrm{\Omega }\delta (xx^{})\delta (tt^{})\delta _{ij}.`$ (4)
where $`i,j\{1,2\}`$ and the heat bath fields $`\xi _i(x,t)`$ obey the fluctuation dissipation relation in equilibrium. Thus, for long times, the system equilibrates to its canonical distribution at temperature $`T`$. This can be shown most conveniently by writing the Fokker-Planck equation corresponding to the Langevin dynamics of Eq. (4) ,
$`_tP_{FP}[\pi ,\varphi ,t]=_{\mathrm{FP}}P_{FP}[\pi ,\varphi ,t].`$ (5)
where
$`_{\mathrm{FP}}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Omega }}{2}}{\displaystyle \frac{\delta ^2}{\delta \pi _i^2}}+\pi _i{\displaystyle \frac{\delta }{\delta \varphi _i}}`$ (7)
$`{\displaystyle \frac{\delta }{\delta \pi _i}}\left(\eta \pi _i^2\varphi _i+{\displaystyle \frac{\delta V(\varphi )}{\delta \varphi _i}}\right),`$
where sum over $`i\{1,2\}`$ is implied here and below. If, as in most applications, the potential $`V(\varphi )`$ is explicitly time independent we can invoke a separation ansatz for $`P_{FP}`$ such that
$`P_{FP}[\pi ,\varphi ,t]=𝒫[\pi ,\varphi ]T(t)`$ (8)
Thus we can regard Eq. (5) as a functional Schrödinger equation, in imaginary time. Then we can write the time independent and dependent equations
$`_{\mathrm{FP}}𝒫_n=E_n𝒫_n,_tT(t)=E_nT(t).`$ (9)
The functional dependence on the fields is now limited to the static probability eigenfunctionals $`𝒫_n`$. The time evolution of the Fokker-Planck distribution is completely characterized by the spectrum of eigenvalues of $`_{\mathrm{FP}}`$, $`E_n`$.
Formally, we can then project the evolution of $`P_{FP}`$ in terms of its eigenvalues $`E_n`$ and eigenfunctionals $`𝒫_n`$ as:
$`P_{FP}[\pi ,\varphi ,t]={\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}C_n𝒫_n[\pi ,\varphi ]e^{E_nt}.`$ (10)
where the $`C_i`$’s are the projections of the initial time $`P_{FP}`$ onto the basis of eigenfunctionals $`𝒫_n`$.
The equilibrium solution must be static. It corresponds to a zero eigenvalue in Eq. (9), which implies the solution
$`P_{\mathrm{eq}}[\pi ,\varphi ]=N\mathrm{exp}\left[\beta {\displaystyle d^Dx\frac{\pi _i^2}{2}}+{\displaystyle \frac{(\varphi _i)^2}{2}}+V[\varphi ]\right],`$ (11)
where we took $`\mathrm{\Omega }=2\eta /\beta `$, which is the Einstein relation enforcing equilibrium between fluctuation and dissipation at long times. Summation over $`i`$ is implied. On general grounds we expect the canonical equilibrium distribution to be approached at long times, i.e. we expect the excited time-dependent states $`P_n`$, $`n0`$ to decay with time.
The full spectrum of excited states and their corresponding eigenvalues can only be found for specific forms of the field potential $`V(\varphi )`$. In particular this is possible in closed analytic form for harmonic potentials $`V=\frac{1}{2}m^2\varphi ^2`$. For each mode in k-space the excited states are given in terms of Hermite polynomials of functions of the field modes and those of their conjugate momenta. More importantly the corresponding eigenvalues are given by
$`E_n`$ $`=n{\displaystyle \frac{\eta }{2}}\left[1\pm \sqrt{14(k^2+m^2)/\eta ^2}\right].`$ (12)
Close to the phase transition the leading effect of the $`\frac{\lambda }{4}\varphi ^4`$ interactions is to make the effective mass temperature dependent as
$`m^2(T)=m_0^2|{\displaystyle \frac{TT_c}{T_c}}|^{2\nu }`$ (13)
where $`T_c`$ is the critical temperature and $`\nu `$ a universal critical exponent, which depends only on the dimensionality of space and the internal symmetries of the theory. To 1-loop in perturbation theory we have
$`m^2(T)=m^2+\mathrm{\Delta }m^2(T),`$ (14)
$`\mathrm{\Delta }m^2(T)=\lambda {\displaystyle \frac{d^Dk}{(2\pi )^D}\varphi _k\varphi _k}.`$ (15)
i.e. the temperature correction to $`m^2`$ is given by the (classical) tadpole diagram. The values of the $`O(2)`$ symmetric thermal average $`\varphi _k\varphi _k`$ depends on the specific form of the thermal distribution, classical or quantum. Under these approximations one obtains the mean-field value of $`\nu =1/2`$.
In the critical domain we can thus obtain an approximate analytical description of the non-linear field dynamics by taking the mass in the harmonic potential to be of the form (13). Although only approximately true we will show below that this assumption leads to a good qualitative understanding of the full non-linear field dynamics in the critical domain and the effects of the Ginzburg regime.
The present scheme, Eqs. (4), is therefore convenient both as a thermalization algorithm, if the system is started at any given configuration and run for long times, or as a means of performing effective non-equilibrium dynamics. The equilibration time itself $`t_{eq}1/E_1`$, is dependent on spatial scale (or wave length) and on temperature. Qualitatively large spatial scales equilibrate more slowly and in particular, in the critical domain, exhibit critical slowing down. This can be seen explicitly by considering long wave-length modes ($`k^20`$) in the vicinity of $`T_c`$. Then, for half of the eigenvalues the termalization time is inversely proportional to $`n`$ times $`t_{eq}`$ with
$`t_{eq}=`$ $`\left[{\displaystyle \frac{\eta }{2}}\left(1\sqrt{14(k^2+m^2)/\eta ^2}\right)\right]^1`$ (17)
$`{\displaystyle \frac{\eta }{m^2(T)}}\mathrm{},`$
as $`TT_c`$. This is the result for overdamped dynamics $`\eta >>m(T)`$, and could have been obtained by a Langevin equation with a single (dissipative) time derivative. In this sense the evolution of the long-wave length modes in the vicinity of $`T_c`$ is always overdamped, which is the essence of the perhaps more familiar TDGL evolution, to which Eqs. (4) reduce to in the appropriate regime. Note that the TDGL dynamics is an effective equation for long-wave length field modes in the critical domain while our model holds more generally.
In the converse limit the decay of short wave-length modes is dictated by $`\eta `$ and is thus scale invariant in this approximation. The appropriate physical value of $`\eta `$ can be computed in perturbation theory given a microscopic model. Particularly interesting are situations for which $`m(T)<\eta `$ as happens in the critical domain. Then there is true time-scale separation in the sense that short wavelength modes thermalize much faster than long-wave length modes.
This kind of considerations will help us understand the behavior of the fully non-linear dynamics in the Ginzburg regime. Before we do this we need to develop a clear picture of equilibrium to which we now turn.
## III Equilibrium results and the definition the Ginzburg temperature
The rationale behind the original proposal according to which the Ginzburg temperature $`T_G`$ is the energy scale for the formation of topological defects , was that, at lower temperatures, thermal fluctuations would be unable to overcome the potential energy barrier associated with the defect’s topological stability. Thus, field configurations with non-trivial topology, below this temperature would necessarily acquire stability on the average.
It is clear that such an appealingly simple physical picture assumes implicitly a separation of physical scales and associated time evolution or equivalently, as we discuss below, that at least part of the system is out of thermal equilibrium. Indeed one must assume that field configurations can be separated in two populations - one of topological defects and another of thermal fluctuations. The former, at least in the sense of the definition of the Ginzburg regime (see below) live on a zero temperature background. This situation is at best an idealization.
Strictly in thermal equilibrium at temperatures not too low, field thermal fluctuations with non-trivial topology will always exist. The density of vortex string thermal fluctuations in our model is shown in Fig. 1. It is, however, remarkable that long strings can only exist in equilibrium strictly above $`T_c`$ . This phenomenon is the analog of vortex pair unbinding in the well known Kosterlitz-Thouless transition in this very same model in 2D . In 3D, however, long strings appear abruptly as we are dealing with a true critical phenomenon instead of a crossover.
The appearance of long strings exactly at $`T_c`$ can be understood, in turn, in terms of their thermal statistical properties namely their tension $`\sigma _{\mathrm{eff}}(T)`$ (free energy per unit length) and other statistical properties like their fractal or Hausdorf dimension . The dependence of the string tension on temperature is shown in Fig. 2. As seen the string tension diminishes continuously with increasing temperature until the critical point where it vanishes. This permits the creation of strings of all lengths above $`T_c`$. Below $`T_c`$, on the other hand strings are exponentially suppressed and only those smaller than the temperature dependent length $`lT/\sigma _{\mathrm{eff}}(T)`$ are likely as thermal fluctuations. It is the existence of long strings as thermal fluctuations that will lead to defect formation if the system is suddenly cooled .
Although some of the above comments may seem somewhat marginal they establish that the thermodynamics of the $`U(1)`$ theory under consideration is much richer than the assumptions on which the traditional role of the Ginzburg regime is based. They show in particular that the vortex strings themselves as a subset of the theory’s thermal fluctuations have a very non-trivial thermodynamics and cannot be taken as their cold classical solutions over a nontrivial thermal background.
The thermodynamics of vortex strings in more complex field theories, with gauge fields and larger symmetry groups, remains to date largely unstudied, although some work has been done in the Abelian case . We expect nevertheless that most of the features of our $`U(1)`$ global theory may persist albeit characterized by different critical exponents. This statement is supported by analytic studies of the statistics of free strings , which, thanks to their large configurational entropy, exhibit a similar (Hagedorn) transition but display eg. a different dependence of the string tension on temperature. The interactions thus change particular temperature dependences of certain quantities but not their qualitative behavior. In theories where defects are not topologically stable however (as in the case of embedded or semilocal defects) the role of these configurations may be potentially different and possibly marginal.
The cumulative results from the equilibrium study of the thermodynamics of vortex strings in our model cast considerable doubt upon the assumptions underlying the traditional role of the Ginzburg temperature in defect formation. It remains unclear however what the role may be of large thermal fluctuations in the critical domain (above $`T_G`$) in changing defect densities produced by the critical dynamics of the fields.
In order to investigate this issue we need a quantitative definition of $`T_G`$. In tune with the arguments given above consider a volume of characteristic size $`\xi (T)`$, the correlation length, and a theory with two energetically degenerate minima of an effective potential $`V(\varphi )`$, separated by a potential barrier $`\mathrm{\Delta }V`$. This applies also for theories with a general $`O(N)`$ symmetry, since we will be interested in the ’radial’ direction only. The effective potential is obtained by singling out an arbitrary direction in field space , which we denote by $`\phi `$. The rate for the field to change coherently from one minimum to the other per unit volume due to thermal activation is
$`T^4\mathrm{exp}\left(\mathrm{\Delta }V/k_BT\right).`$ (18)
For an effective potential of the form (obtained, eg. perturbatively at 1-loop)
$`V(\varphi )={\displaystyle \frac{1}{2}}m^2(T)\phi ^2+{\displaystyle \frac{\lambda }{4}}\phi ^4`$ (19)
$`\mathrm{\Delta }V=\frac{m(T)^4}{4\lambda }`$. For a volume $`\xi ^3`$, we define $`T_G`$ such that the probability of overcoming the potential barrier is of order unity:
$`T_G:{\displaystyle \frac{\mathrm{\Delta }V(T_G)}{T_G}}.\xi ^3(T_G)=1\lambda T_G/m(T_G)={\displaystyle \frac{1}{4}}.`$ (20)
This definition however has some caveats, for instance, an effective potential of the form Eq. (19) is only valid for the mean field and not on smaller scales. A more careful accounting of scales leads to different results , which show an enhancement of the hoping probability. Thus, the factor of $`1/4`$ in Eq. (20) should not be taken at face value.
A perhaps more rigorous definition arises from the range of temperatures below $`T_c`$ for which fluctuations are large and consequently where perturbative finite temperature field theory fails to be useful. In order to set up a perturbative scheme at finite temperature from an initial 3+1 dimensional quantum field theory one implements dimensional reduction which is valid provided the temperature is high compared to all mass scales. As a consequence the coupling of the dimensionally reduced 3D field theory becomes dimensionful, i.e. $`\lambda \lambda T=\lambda _3`$. In order to proceed one has to identify an appropriate dimensionless coupling. This is done by taking $`\lambda T/m(T)`$. The Ginzburg regime is entered when this 3D effective coupling becomes strong, in the vicinity of the critical point, namely
$`T_G:\lambda T_G/m(T_G)=1.`$ (21)
To compute $`T_G`$ one needs the scaling of $`m(T)`$ in the critical domain. We write $`m^2(T)=m_0^2|ϵ|^\nu `$, with $`ϵ`$ being the reduced temperature $`ϵ=\frac{TT_c}{Tc}`$.
Thus $`ϵ_G=0.18`$ for $`\nu =0.5`$. This mean-field estimate produces an upper bound in $`T`$ for $`T_G`$ (and lower bound for $`\beta =1/T`$). For realistic 3D exponents, $`\nu =0.67`$, we obtain $`ϵ_G=0.25`$. The first criterion, based on the hopping of a correlation sized volume, results in higher values of $`T_G`$. This brings about a relatively large uncertainty in the value of $`T_G`$, which is $`1825\%`$ below $`T_c`$.
## IV The role of the Ginzburg regime in the dynamics of defect formation
In order to investigate the role of the Ginzburg temperature in the dynamics of defect formation we perform a series of tests both directly over the evolution of string densities and the evolution of the order parameter, when exposed to thermal fluctuations in the Ginzburg regime.
### A Strings Survive the Ginzburg Regime
To investigate the effects of thermal fluctuations directly upon strings we deliberately expose the system to a heat bath at temperature $`ϵ_i`$, within the Ginzburg regime and below.
We are attempting to emulate the worst case scenario of an experimental quench where the temperature or pressure are dropped monotonically but where the system makes a long stopover within the Ginzburg regime. We repeat this procedure for a range of time intervals $`\mathrm{\Delta }t`$, after which the bath temperature is taken to zero. This set of temperature trajectories is shown in Fig. 3.
We would expect that, if the Ginzburg regime indeed produced enhanced decay of strings, then the string densities measured at later times should be smaller the longer the time the system spent within the range $`T_cTT_G`$.
Our results for the final string densities as a function of intermediate temperature $`ϵ_i`$ and $`\mathrm{\Delta }t`$ are summarized in Fig. 4. There is no apparent effect of the Ginzburg regime in reducing string densities at formation.
If any trend is visible from Fig. 4 it is the opposite, namely a monotonic dependence of the final string densities on $`ϵ_i`$ \- the lower $`ϵ_i`$, the less string is measured at later times.
The knowledge of the vortex string thermodynamics and of the time response of the fields in the critical domain again helps us understand this result. Strings and in particular long strings are inherited from high temperature (higher than $`T_c`$) topological fluctuations .
We can now use our knowledge of the Fokker-Planck solution to understand the observations of Fig. 4. As we discussed above the small scales in the system equilibrate faster on a characteristic timescale $`t\eta ^1`$. Small scale fluctuations dominate the thermal average in (15), which then allows us to take the effective value of $`m^2m^2(T_i)`$.
On the other hand, upon cooling through the critical point the large scales in the system display critical slowing down. This includes in particular the long strings in the sample which will be responsible for the signal at the time of measurement later. This slowing down leads to an imbalance in the string population out of equilibrium relative to their equilibrium counterpart, given by the existence of many more long strings.
This constitutes an excited state (described by $`𝒫_{n0}`$) relative to the true equilibrium of the system at intermediate temperatures below $`T_c`$. These states will then decay on a timescale $`t_{eq}=E_n^1\eta /m^2(T_i)`$. The value of $`m^2(T_i)`$ is smaller the closer $`T_i`$ is to $`T_c`$ and thus leads to a longer time scale for the decay of long strings.
We can then predict a monotonic behavior for the string densities as observed in Fig. 4. At $`T_G`$ in particular $`m^2=\lambda T_G^2`$, by definition and $`t_{eq}\lambda T_G^2/\eta `$.
Thus the conclusion is inescapable: The longer the time the system spends further from $`T_c`$ the less string it will display at later times where formation rates are measured.
### B Memory of the Order Parameter Configuration near $`T_c`$
An independent test on the possible role of thermal fluctuations in affecting string densities consists in reheating a quenched system to a temperature around $`T_c`$ (both below and above) and cooling it again at the same rate. This process tests the memory of the order parameter as well as that of other related quantities (see also ), such as defects.
The importance of this test is directly related to the canonical theory of defect formation as due to the critical dynamics of the fields. The final density of strings formed at the transition are then infered indirectly through the correlation length associated with the two-point correlator of the fields.
An example of the temperature ($`ϵ(t)`$) trajectories used in testing the memory of the order parameter are shown in Fig. 5a.
We are particularly interested in investigating under what circumstances thermal fluctuations can affect the large scale configuration of the order parameter. In order to produce a quantitative test we define the unequal-time two-point correlation function
$`\varphi (x,t_{\mathrm{rh}})\varphi (x,t+t_{\mathrm{rh}})`$ (22)
$`=N{\displaystyle \underset{j=1}{\overset{2}{}}}{\displaystyle \underset{i}{}}\varphi _j(x_i,t_{\mathrm{rh}})\varphi _j(x_i,t+t_{\mathrm{rh}}),`$ (23)
where $`N`$ is an irrelevant normalization factor. This correlator has several interesting properties. For short times it displays a characteristic time, which describes the decay of correlations over very small spatial scales. This is the initial transient in Fig. 5b. We verified that this time and the form of the correlation function is in agreement with the forms predicted from a Boltzmann distribution for the fields. For later times the residual correlation comes from the motion of the order parameter (the field volume average). This average can be either positive or negative but, if thermal, will converge to zero at and above $`T_c`$.
Now, we are interested in determining whether the final field configuration over large spatial scales is correlated to the configuration prior to reheating. Fig. 5 shows that only if one crosses $`T_c`$, by more than $`+\widehat{ϵ}`$, is the memory of the initial quenched configuration erased (see in particular the two trajectories reaching higher temperatures in comparison to the others). For these trajectories the field correlations reach zero and after reheating evolve to a value manifestly different from that prior to reheating.
For trajectories within the Ginzburg regime, that do not cross $`T_c`$, the change in the configuration of the order parameter as measured by Eq. (23) is small. In particular the field configuration existing before reheating is approximately recovered as the fields are cooled. The same is true for the string densities, including those of long strings.
Again we can understand these results using the tool developed in section II. The thermalization time, which is nothing else than the response time $`\tau `$ in the Kibble-Zurek scenario, for long-wave length modes in the system is given by
$`t_{eq}{\displaystyle \frac{\eta }{m^2(T)}}.`$ (24)
We argued that it is a reasonable approximation to take $`T`$ to be the final temperature since the small scales in the system equilibrate much faster (provided $`m(T)<\eta `$).
Now $`\widehat{t}`$, or equivalently $`\widehat{ϵ}`$, is defined as the time when the long wave-length modes in the system can first respond to a change in bath temperature linear in time. It is computed by equating the scaling of the response time $`\tau (T)`$
$`\tau (ϵ(T))={\displaystyle \frac{\eta }{m^2ϵ^{\nu z}(T)}}`$ (25)
to a linear change in time, imposed externally, i.e.
$`\tau (ϵ(\widehat{t}))={\displaystyle \frac{\eta }{m^2ϵ^{\nu z}(\widehat{t})}}=\widehat{t}`$ (26)
$`z`$ is another critical exponent whose mean-field value is 2, see . This relation is usually solved by assuming a linear dependence in time for $`ϵ(t)=t/\tau _Q`$, where $`\tau _q`$ is the rate of the external quench.
Explicit calculation of $`\mathrm{\Delta }m^2(T)`$ in (15) for our classical Boltzmann ensemble shows that
$`\mathrm{\Delta }m^2(T){\displaystyle \frac{\lambda }{\pi ^2}}\mathrm{\Lambda }T,`$ (27)
where $`\mathrm{\Lambda }`$ is the ultraviolet cut off. This cutoff has physical meaning and is related to the breakdown of our scalar field model for high energy excitations, eg. fermionic quasi-particles in $`{}_{}{}^{3}He`$. If indeed the external bath temperature is changed linearly then $`T`$ in Eq. (27) can be taken, over time scales larger than $`\eta ^1`$, to be linear in time. It then results trivially that $`ϵ(t)`$ is also linear, which validates our assumption.
The significance of $`\widehat{t}`$ is that only when $`|ϵ(t)|>\widehat{ϵ}`$ can the long wave-length modes in our system thermalize under an externally changing temperature at a rate $`\tau _Q`$, i.e. keep pace with the externally imposed changes. Due to theoretical uncertainties the value of $`\widehat{ϵ}`$ adopted in Fig. 5a was measured by monitoring the response of the system directly. Details are described elsewhere .
At the initial time, for temperature trajectories of Fig. 5, the system is in the process of breaking the $`U(1)`$ symmetry spontaneously, i.e. the expectation value of the $`k=0`$ mode of $`\varphi `$, $`\varphi `$ is non-zero. Then, as the system is heated towards $`T_c`$ equilibration of the long wave-length modes means that $`\varphi 0`$ and upon cooling show zero correlation in Fig. 5b to its initial state. Since thermalization of $`\varphi `$ can only occur for $`ϵ\widehat{ϵ}`$, only the temperature trajectories crossing $`+\widehat{ϵ}`$ can attain zero correlations.
It is then clear that the Ginzburg regime cannot change the symmetry breaking process of the system, including its associated long string configurations, unless a long amount of time is allowed. The Ginzburg regime is therefore less efficient at destroying topological defects (in the sense of requiring a longer time) than any other temperature range outside the critical domain.
## V Discussion and conclusions
In this paper we have performed the most extensive analysis to date of the effects of large thermal fluctuations, within the Ginzburg regime, on the formation of topological defects. Our model field theory has already been studied extensively both in equilibrium and in tests of the theory of defect formation at temperature quenches, as predicted by the critical dynamics of the theory.
Under these controlled circumstances we analyzed critically the assumptions underlying the traditional argument for the Ginzburg temperature as the energy scale at which topological defects are formed. We then proceeded to show that the effects of thermal fluctuations in the Ginzburg regime upon a population of topological defects formed by the critical dynamics of the theory carries no particular signature and leads mostly to small qualitative changes in the defect densities predicted by the theory of defect formation.
We have also shown that even prolonged exposure of a quenched field configuration to the Ginzburg regime has little consequences in changing the order parameter configurations emerging at $`\widehat{ϵ}`$, and associated string densities. In addition we established that to truly destroy a quenched field configuration existing below $`\widehat{ϵ}`$, one has to expose the system to temperatures well above $`T_c`$. In particular for a linear quench trajectory, a temperature of $`ϵ+\widehat{ϵ}`$, must be reached in order to erase memory of the initial configuration.
These results were confirmed by analytical arguments based on the solutions of the associated Fokker-Plank equation. This analysis supports the conclusion that given the same amount of time of exposure to a thermal bath at a given temperature, the Ginzburg regime is actually the least efficient range of temperatures at destroying the pattern of symmetry breaking inherited from criticallity. This includes topological defect configurations.
Our results fully support the theory of defect formation resulting from the critical dynamics of second order transitions and all known thermodynamic results for vortex strings in $`O(N)`$ theories . In face of this evidence we are lead to conclude that arguments singling out a special energy scale $`T_GT_c`$, which would play an important role in defect formation rely on assumptions that are not realized in the true (thermo)dynamics of our model and are thus invalid.
Thus we expect the results of this paper to carry over to the new Lancaster $`{}_{}{}^{4}He`$ experiments. The results of reported in in these experiments cannot therefore be attributed to the effects of Ginzburg regime in $`{}_{}{}^{4}He`$.
## Acknowledgments
We thank T. Kibble, P. Laguna and R. Rivers for useful discussions. Numerical work was done on the T-division/CNLS Avalon Beowulf cluster, LANL. This research was supported by the U.S. Department of Energy, under contract W-7405-ENG-36. |
warning/0001/math0001055.html | ar5iv | text | # Left-modular Elements
## 1 Left-modular elements
Throughout this paper $`L`$ is a finite lattice where $`\widehat{0}=\widehat{0}_L`$ and $`\widehat{1}=\widehat{1}_L`$ are the minimal and maximal elements, respectively. We say that $`x`$ is covered by $`y`$, and write $`xy`$, if $`x<y`$ and there is no element $`zL`$ such that $`x<z<y`$.
We use $``$ for the meet (greatest lower bound) and $``$ for the join (least upper bound) in $`L`$. Given any $`x,y,zL`$ with $`z<y`$, the *modular inequality*
$$z(xy)(zx)y$$
(1)
is always true and equality holds whenever $`y`$ or $`z`$ is comparable to $`x`$. We say that $`x`$ and $`y`$ form a *modular pair* $`(x,y)`$ if (1) is an equality for any $`z<y`$. Note that this relation is not symmetric, in general. Two kinds of elements are associated to the modular pair:
###### Definition 1.1
1. An element $`x`$ is called a *left-modular element* if $`(x,y)`$ is a modular pair for every $`yL`$.
2. An element $`x`$ is called a *modular element* if both $`(x,y)`$ and $`(y,x)`$ are modular pairs for every $`yL`$.
In a semimodular lattice with rank function $`\rho `$, the pair $`(x,y)`$ is modular if and only if $`\rho (xy)+\rho (xy)=\rho (x)+\rho (y)`$ \[2, p. 83\]; so in this case the relation of being a modular pair is symmetric, and then there is no difference between modularity and left-modularity. However, there are examples such as the non-crossing partition lattices (see Sec. 3) and the Tamari lattices where the two concepts do not coincide.
Let $`L`$ be a graded lattice of rank $`n`$ with rank function $`\rho `$. Then the *characteristic polynomial* of $`L`$ is defined by
$$\chi (L,t)=\underset{xL}{}\mu (x)t^{n\rho (x)}$$
where $`t`$ is an indeterminate, $`\mu :L\times L`$ is the Möbius function of $`L`$, and $`\mu (x)=\mu (\widehat{0},x)`$. There are two important factorization theorems for $`\chi `$ given by R. Stanley:
###### Theorem 1.2 (Partial Factorization Theorem )
Let $`L`$ be an atomic, semimodular lattice (i.e., a geometric lattice) of rank $`n`$. If $`x`$ is a modular element of $`L`$, then
$$\chi (L,t)=\chi ([\widehat{0},x],t)\underset{b:bx=\widehat{0}}{}\mu (b)t^{n\rho (x)\rho (b)}.\text{ }$$
###### Theorem 1.3 (Total Factorization Theorem )
Let $`(L,\mathrm{\Delta })`$ be a supersolvable, semimodular lattice of rank $`n`$ with $`\mathrm{\Delta }:\widehat{0}=x_0x_1\mathrm{}x_n=\widehat{1}`$. Then
$$\chi (L,t)=(ta_1)(ta_2)\mathrm{}(ta_n)$$
(2)
where $`a_i`$ is the number of atoms of $`L`$ that are below $`x_i`$ but not below $`x_{i1}`$
Note that all elements in the maximal chain $`\mathrm{\Delta }`$ of a supersolvable lattice are left-modular (see ). So the hypotheses of Theorem 1.3 imply that they are modular. In recent work , A. Blass and B. Sagan generalized the Total Factorization Theorem to LL lattices where the first “L” stands for the fact that the lattice has a maximal chain all of whose elements are all left-modular. The purpose of this paper is to generalize the Partial Factorization Theorem by replacing the modular element with a left-modular one and relaxing the hypotheses requiring that the lattice be atomic and semimodular. To do so, we will derive a general characterization of left-modular elements in this section. In the next section, we introduce a generalized rank function for a lattice which might not be graded in the usual sense, and then develop a general formula for the characteristic polynomial of a lattice with a left-modular element in Theorem 2.3. Under an extra rank-preserving hypothesis we obtain our generalization of the Partial Factorization Theorem (Theorem 2.6). In Sections 3 and 4, we calculate the characteristic polynomials and the Möbius functions of the non-crossing partition lattices and the shuffle posets by using these two formulae, respectively. The last section contains two inductive proofs for Blass and Sagan’s Total Factorization Theorem for LL lattices using our two main theorems. Consequently, our factorization theorem generalizes the three others.
We say that $`y`$ is a *complement* of $`x`$ if $`xy=\widehat{0}`$ and $`xy=\widehat{1}`$. Stanley showed that, in an atomic and semimodular lattice, $`x`$ is modular if and only if no two complements of $`x`$ are comparable. The next theorem provides an analog for left-modular elements.
###### Theorem 1.4
Let $`x`$ be an element of any lattice $`L`$. The following statements are equivalent:
1. The element $`x`$ is left-modular.
2. For any $`y`$, $`zL`$ with $`z<y`$, we have $`xzxy`$ or $`xzxy.`$
3. For any $`y`$, $`zL`$ with $`zy`$, we have $`xz=xy`$ or $`xz=xy`$ but not both.
4. For every interval $`[a,b]`$ containing $`x`$, no two complements of $`x`$ with respect to the sublattice $`[a,b]`$ are comparable.
Proof. We will prove the implications (i) $``$ (ii) $``$ (iii) $``$ (i). The proof of (ii) $``$ (iv) is immediate.
First we make some preliminary observations. Suppose $`z<y`$. We claim that $`xy=xz`$ if and only if $`y=(zx)y`$. The forward direction is trivial since $`(xy)y=y`$. For the reverse, note that $`y=(zx)y`$ implies $`yxz`$. Now $`z<yxz`$, and joining all sides with $`x`$ gives $`xy=xz`$. Dually $`xy=xz`$ if and only if $`z=z(xy)`$.
For any $`z<y`$ the inequalities
$$zz(xy)(zx)yy$$
(3)
are true by the modular inequality (1). Since $`zy`$, at least one of the $``$’s in (3) should be $`<`$. Therefore (i) $``$ (ii). If $`zy`$, then exactly two of the $``$’s should be $`=`$ and the remaining one must be $``$. Thus (ii) $``$ (iii).
To show (iii) $``$ (i), let us consider the contrapositive: assume that there are $`u`$, $`vL`$ with $`u<v`$ such that $`u(xv)<(ux)v`$. Given any $`y`$, $`z[u(xv),(ux)v]`$ with $`zy`$, we have $`y(ux)vv`$. This implies $`u(xy)u(xv)z`$, so that $`xyz`$. It follows that $`xz=xy`$. Similarly, we can get $`xz=xy`$.
The existence of a left-modular element in $`L`$ implies that such elements are also present in certain sublattices as the next proposition shows.
###### Proposition 1.5
Let $`x`$ be a left-modular element in lattice $`L`$. Then for any $`yL`$
1. the meet $`xy`$ is a left-modular element in $`[\widehat{0},y]`$, and
2. the join $`xy`$ is a left-modular element in $`[y,\widehat{1}]`$.
Proof. Let $`a`$, $`b[\widehat{0},y]`$ with $`b<a`$. By left-modularity of $`x`$, we have
$`b((xy)a)`$ $`=`$ $`b(x(ya))=(bx)(ya)`$
$`=`$ $`((bx)y)a=(b(xy))a.`$
So $`xy`$ is a left-modular element in $`[\widehat{0},y]`$. The proof for join is similar.
## 2 The characteristic polynomial
We begin with a general lemma.
###### Lemma 2.1
Let $`L`$ be a lattice with an arbitrary function $`r:L`$ and let $`n`$. If $`xL`$ is a left-modular element, then
$$\underset{yL}{}\mu (y)t^{nr(y)}=\underset{bx=\widehat{0}}{}\mu (b)\underset{y[b,bx]}{}\mu (b,y)t^{nr(y)}.$$
Proof. We will mimic Stanley’s proof in . By Crapo’s Complementation Theorem , for any given $`a[\widehat{0},y]`$
$$\mu (y)=\underset{a^{},a^{\prime \prime }}{}\mu (\widehat{0},a^{})\zeta (a^{},a^{\prime \prime })\mu (a^{\prime \prime },y),$$
where $`a^{}`$ and $`a^{\prime \prime }`$ are complements of $`a`$ in $`[\widehat{0},y]`$, and $`\zeta `$ is the zeta function defined by $`\zeta (u,v)=1`$ if $`uv`$ and $`\zeta (u,v)=0`$ otherwise. Let us choose $`a=xy`$. The element $`a`$ is left-modular in $`[\widehat{0},y]`$ by Proposition 1.5. But no two complements of $`a`$ in $`[\widehat{0},y]`$ are comparable by Theorem 1.4. Thus
$$\mu (y)=\underset{b}{}\mu (\widehat{0},b)\mu (b,y),$$
(4)
where the sum is over all complements $`b`$ of $`a`$ in $`[\widehat{0},y]`$, i.e., over all $`b`$ satisfying $`by`$, $`b(xy)=\widehat{0}`$ and $`b(xy)=y`$. Since $`x`$ is left-modular, it is equivalent to say that the sum in (4) is over all $`bL`$ satisfying $`bx=\widehat{0}`$ and $`y[b,bx]`$. Thus we have
$`{\displaystyle \underset{yL}{}}\mu (y)t^{nr(y)}`$ $`=`$ $`{\displaystyle \underset{yL}{}}{\displaystyle \underset{\genfrac{}{}{0pt}{}{bx=\widehat{0}}{\text{}y[b,bx]}}{}}\mu (\widehat{0},b)\mu (b,y)t^{nr(y)}`$
$`=`$ $`{\displaystyle \underset{bx=\widehat{0}}{}}\text{ }\mu (b){\displaystyle \underset{y[b,bx]}{}}\mu (b,y)t^{nr(y)}.\text{ }`$
Obviously, the previous lemma is true for the ordinary rank function if $`L`$ is graded. To apply this result to more general lattices we make the following definition.
###### Definition 2.2
A *generalized rank function* of a lattice $`L`$ is a function $`\rho :\{(x,y)L\times Lxy\}`$ such that for any $`abc`$
$$\rho (a,c)=\rho (a,b)+\rho (b,c).$$
In this case, we say $`L`$ is *generalized graded* by $`\rho `$.
For short we write $`\rho (x)=\rho (\widehat{0},x)`$. Conversely, if we take any function $`\rho :L`$ such that $`\rho (\widehat{0})=0`$, then we can easily construct a generalized rank function, namely $`\rho (x,y)=\rho (y)\rho (x)`$. So the ordinary rank function is a special case.
If $`L`$ is generalized graded by $`\rho `$, we now define a generalized characteristic polynomial of $`L`$ by
$$\chi (L,t)=\underset{xL}{}\mu (x)t^{\rho (x,\widehat{1})}=\underset{xL}{}\mu (x)t^{\rho (\widehat{1})\rho (x)}.$$
(5)
Note that $`\chi `$ will depend on which generalized rank function we pick. Since the restriction of a generalized rank function to an interval $`[a,b]`$ still satisfies Definition 2.2 with $`L=[a,b]`$, the characteristic polynomial of the interval is defined in the same manner.
The following theorem, which follows easily from Lemma 2.1, is one of our main results. In it, the *support* of $`\mu `$ is defined by
$$H(L)=\{xL\mu (x)0\}.$$
###### Theorem 2.3
Let $`L`$ be generalized graded by $`\rho `$. If $`xL`$ is a left-modular element, then
$$\chi (L,t)=\underset{\genfrac{}{}{0pt}{}{bH\left(L\right)}{\text{}bx=\widehat{0}}}{}\left[\mu (b)t^{\rho (\widehat{1})\rho (bx)}\chi ([b,bx],t)\right].\text{ }$$
(6)
In the sum (6), the term $`\chi ([b,bx],t)`$ depends on $`b`$. To get a factorization formula, we will remove the dependency by applying certain restrictions so that $`\chi ([b,bx],t)=\chi ([\widehat{0},x],t)`$ for all $`b`$ in the sum.
First, we will obtain a general condition under which two lattices have the same characteristic polynomial. In the following discussion, let $`L`$ and $`L^{}`$ be lattices and let $`\tau :LL^{}`$ be any map. For convenience, we also denote $`\widehat{0}=\widehat{0}_L`$, $`\widehat{0}^{}=\widehat{0}_L^{}`$ and similarly for $`\widehat{1}`$, $`\widehat{1}^{}`$, $`\mu `$, $`\mu ^{}`$, etc.
We say $`\tau `$ is a *join-preserving* map if
$$\tau (uv)=\tau (u)\tau (v)$$
for any $`u`$, $`vL`$. Note that from this definition $`\tau `$ is also order-preserving since
$$xyy=xy\tau (y)=\tau (xy)=\tau (x)\tau (y)\tau (x)\tau (y).$$
If $`\tau `$ is join-preserving, then given any $`x^{}\tau (L)`$, we claim that the subset $`\tau ^1(x^{})`$ has a unique maximal element in $`L`$. Suppose that $`\tau (u)=\tau (v)=x^{}`$ for some $`u`$, $`vL`$. We have $`\tau (uv)=\tau (u)\tau (v)=x^{}`$. Thus $`uv\tau ^1(x^{})`$ and the claim follows.
If, in addition, $`\tau `$ is surjective then we can define a map $`\sigma :L^{}L`$ by
$$\sigma (x^{})=\text{ the maximal element of }\tau ^1(x^{}).$$
(7)
The map $`\sigma `$ must also be order preserving. To see this, suppose $`x^{}y^{}`$ in $`L^{}`$ and consider $`x=\sigma (x^{}),y=\sigma (y^{})`$. Then
$$\tau (xy)=\tau (x)\tau (y)=x^{}y^{}=y^{}.$$
So $`xy\tau ^1(y^{})`$ which forces $`xyy`$ by definition of $`\sigma `$. Thus $`xy`$ as desired.
###### Lemma 2.4
Using the previous notation, suppose that $`\tau `$ is surjective and join-preserving and that $`\sigma `$ satisfies $`\sigma (\widehat{0}^{})=\widehat{0}`$. Then for any $`x^{}L^{}`$ we have
$$\mu ^{}(x^{})=\underset{y\tau ^1(x^{})}{}\mu (y).$$
Proof. This is trivial when $`x^{}=\widehat{0}^{}`$. Let $`x=\sigma (x^{})`$. From the assumptions on $`\tau `$ and $`\sigma `$ it is easy to see that
$$[\widehat{0},x]=\underset{y^{}[\widehat{0}^{},x^{}]}{}\tau ^1(y^{}).$$
(8)
Now, by surjectivity of $`\tau `$ and induction, we get
$$\mu ^{}(x^{})=\underset{y^{}<x^{}}{}\mu ^{}(y^{})=\underset{\genfrac{}{}{0pt}{}{y\tau ^1\left(y^{}\right)}{\text{}y^{}<x^{}}}{}\mu (y)=\underset{y\tau ^1(x^{})}{}\mu (y).\text{ }$$
Let $`L`$ and $`L^{}`$ be generalized graded by $`\rho `$ and $`\rho ^{}`$, respectively. We say an order-preserving map $`\tau :LL^{}`$ is *rank-preserving* on a subset $`SL`$ if $`\rho (x,y)=\rho ^{}(\tau (x),\tau (y))`$ for any $`x`$, $`yS`$, $`xy`$.
###### Lemma 2.5
If, in addition to the hypotheses of Lemma 2.4, the map $`\tau `$ is rank-preserving on $`H(L)\{\widehat{1}\}`$ then
$$\chi (L,t)=\chi (L^{},t).$$
Proof. From (8) in the proof of Lemma 2.4, we know $`L=_{x^{}L^{}}\tau ^1(x^{})`$. Then by Lemma 2.4 and the rank-preserving nature of $`\tau `$, we have
$`\chi (L^{},t)`$ $`=`$ $`{\displaystyle \underset{x^{}L^{}}{}}\mu ^{}(x^{})t^{\rho ^{}(x^{},\widehat{1}^{})}`$
$`=`$ $`{\displaystyle \underset{x^{}L^{}}{}}{\displaystyle \underset{y\tau ^1(x^{})}{}}\mu (y)t^{\rho ^{}(x^{},\widehat{1}^{})}`$
$`=`$ $`{\displaystyle \underset{x^{}L^{}}{}}{\displaystyle \underset{y\tau ^1(x^{})H(L)}{}}\mu (y)t^{\rho ^{}(\tau (y),\tau (\widehat{1}))}`$
$`=`$ $`{\displaystyle \underset{yH(L)}{}}\mu (y)t^{\rho (y,\widehat{1})}`$
$`=`$ $`\chi (L,t).\text{ }`$
It is easy to generalize the previous lemma to arbitrary posets as long as the map $`\sigma `$ is well defined. However, we know of no application of the result in this level of generality.
Returning to our factorization theorem, we still need one more tool. For any given $`a`$, $`b`$ in a lattice, we define
$$\sigma _a:[b,ab][ab,a]\text{by}\sigma _a(u)=ua,$$
$$\tau _b:[ab,a][b,ab]\text{by}\tau _b(v)=vb.$$
The map $`\tau _b`$ is the one we need to achieve $`\chi ([b,bx],t)=\chi ([\widehat{0},x],t)`$. In the following, we write $`H(x,y)`$ for $`H([x,y])`$ which is the support of $`\mu `$ defined on the sublattice $`[x,y]`$. We can now prove our second main result.
###### Theorem 2.6
Let $`L`$ be generalized graded by $`\rho `$ and let $`xL`$ be an left-modular element. If the map $`\tau _b`$ is rank-preserving on $`H(\widehat{0},x)\{x\}`$ for every $`bH(L)`$ satisfying $`bx=\widehat{0}`$. Then
$$\chi (L,t)=\chi ([\widehat{0},x],t)\underset{\genfrac{}{}{0pt}{}{bH\left(L\right)}{\text{}bx=\widehat{0}}}{}\mu (b)t^{\rho (\widehat{1})\rho (x)\rho (b)}.$$
(9)
Proof. First, we will show that $`\chi ([b,bx],t)=\chi ([\widehat{0},x],t)`$ for any $`bH(L)`$ with $`bx=\widehat{0}`$ by verifying the hypotheses of Lemma 2.5. By left-modularity of $`x`$, we have
$$\tau _b\sigma _x(y)=b(xy)=(bx)y=y$$
(10)
for any $`y[b,bx]`$. So $`\tau _b`$ is surjective. And it is easy to check that $`\tau _b`$ is join-preserving. As for $`\sigma _x`$, we must check that it satisfies the definition (7). Given $`z\tau _b^1(y)`$ we have $`y=\tau _b(z)=zb`$. So by the modular inequality (1) we get
$$\sigma _x(y)=yx=(zb)xz(bx)z.$$
Since this is true for any such $`z`$, we have $`\sigma _x(y)\mathrm{max}\tau _b^1(y)`$. But equation (10) implies $`\sigma _x(y)\tau _b^1(y)`$, so we have equality. Finally $`\widehat{0}_{[b,bx]}=b`$ so $`\sigma _x(b)=bx=\widehat{0}`$ as desired.
Now we need only worry about the exponent on $`t`$ in Theorem 2.3. But since $`\tau _b`$ is rank-preserving on $`H(\widehat{0},x)\{x\}`$, we get
$$\rho (bx)=\rho (\widehat{0},b)+\rho (b,bx)=\rho (\widehat{0},b)+\rho (\widehat{0},x)=\rho (b)+\rho (x).\text{ }$$
Here we state a corollary which relaxes the hypothesis in Stanley’s Partial Factorization Theorem.
###### Corollary 2.7
Equation (9) holds when $`L`$ is a semimodular lattice (graded by the ordinary rank function) with a modular element $`x`$.
Proof. To apply Theorem 2.6, it suffices to show that $`\rho (\widehat{0},z)=\rho (b,zb)`$ for every $`z[\widehat{0},x]`$. Since $`(b,x)`$ is a modular pair, we have $`(zb)x=z(bx)=z\widehat{0}=z`$. By Proposition 1.5, $`z=(zb)x`$ is left-modular in $`[\widehat{0},zb]`$, so $`(z,b)`$ is a modular pair in this lattice. Thus $`\rho (zb)+\rho (zb)=\rho (z)+\rho (b)`$, because $`[\widehat{0},zb]`$ is a semimodular lattice. Since $`zb=\widehat{0}`$ we are done.
We take the divisor lattice $`D_n`$ as an example. It is semimodular, but not atomic in general, so Stanley’s theorem does not apply. However, Corollary 2.7 can be used for any $`xD_n`$, since all elements are modular.
We will now present a couple of applications of the previous results in the following two sections.
## 3 Non-crossing Partition Lattices
The non-crossing partition lattice was first studied by Kreweras who showed its Möbius function is related to the Catalan numbers. By using NBB sets (see Sec. 5 for the definition), Blass and Sagan combinatorially explained this fact. In this section we will calculate the characteristic polynomial for a non-crossing partition lattice and then offer another explanation for the value of its Möbius function.
If it causes no confusion, we will not explicitly write out any blocks of a partition that are singletons. Let $`n1`$. We say that a partition $`\pi [n]`$ is *non-crossing* if there do not exist two distinct blocks $`B,C`$ of $`\pi `$ with $`i`$, $`kB`$ and $`j`$, $`lC`$ such that $`i<j<k<l`$. Otherwise $`\pi `$ is *crossing*.
Another way to view non-crossing partitions will be useful. Let $`G=(V,E)`$ be a graph with vertex set $`V=[n]`$ and edge set $`E`$. We say that $`G`$ is *non-crossing* if, when the vertices are arranged in their natural order clockwise around a circle and the edges are drawn as straight line segments, no two edges of $`G`$ cross geometrically. Given a partition $`\pi `$ we can form a graph $`G_\pi `$ by representing each block $`B=\{i_1<i_2<\mathrm{}<i_l\}`$ by a cycle with edges $`i_1i_2,i_2i_3,\mathrm{},i_li_1`$. (If $`|B|=1`$ or 2 then $`B`$ is represented by an isolated vertex or edge, respectively.) Then it is easy to see that $`\pi `$ is non-crossing as a partition if and only if $`G_\pi `$ is non-crossing as a graph.
The set of non-crossing partitions of $`[n]`$, denoted by $`NC_n`$, forms a meet-sublattice of partition lattice $`\mathrm{\Pi }_n`$ with the same rank function. However unlike $`\mathrm{\Pi }_n`$, the non-crossing partition lattice is not semimodular in general, since if $`\pi =13`$ and $`\sigma =24`$ then $`\pi \sigma =\widehat{0}`$ and $`\pi \sigma =1234`$. So we have
$$\rho (\pi )+\rho (\sigma )=2<3=\rho (\pi \sigma )+\rho (\pi \sigma ).$$
The $`\mathrm{\Pi }_n`$-join $`\pi \sigma =13/24`$ also explains why $`NC_n`$ is not a sublattice of $`\mathrm{\Pi }_n`$.
Let $`n2`$ and $`\pi =12\mathrm{}(n1)`$. It is well-known that $`\pi `$ is modular in $`\mathrm{\Pi }_n`$ and so left-modular there. Given any $`\alpha `$, $`\beta NC_n`$ with $`\alpha <\beta `$ and both incomparable to $`\pi `$. It is clear that $`\alpha \pi =\beta \pi =\widehat{1}`$ in $`\mathrm{\Pi }_n`$ as well as in $`NC_n`$. By Theorem 1.4 we get $`\alpha \pi <\beta \pi `$ in $`\mathrm{\Pi }_n`$. Since $`NC_n`$ is a meet-sublattice of $`\mathrm{\Pi }_n`$, this inequality for the two meets still holds in $`NC_n`$. This fact implies that $`\pi `$ is left-modular in $`NC_n`$. In general, $`\pi `$ is not modular in $`NC_n`$. If $`n4`$, let $`\sigma =2n`$ and $`\varphi =1(n1)/23\mathrm{}(n2)`$. Clearly $`\varphi <\pi `$, $`\pi \sigma =\varphi \sigma =\widehat{0}`$ and $`\pi \sigma =\varphi \sigma =\widehat{1}`$ in $`NC_n`$, so that $`(\sigma ,\pi )`$ is not a modular pair.
###### Proposition 3.1
The characteristic polynomial of the non-crossing partition lattice $`NC_n`$ satisfies
$$\chi (NC_n,t)=t\chi (NC_{n1},t)\underset{i=1}{\overset{n1}{}}\chi (NC_i,t)\chi (NC_{ni},t)$$
with the initial condition $`\chi (NC_1,t)=1`$.
Proof. The initial condition is trivial. Let $`n2`$ and $`\pi =12\mathrm{}(n1)`$. We will apply Theorem 2.3. Note that $`b\pi =\widehat{0}`$ if and only if any two numbers of $`[n1]`$ are in different blocks of $`b`$, so either $`b=\widehat{0}`$ or $`b=mn`$ with $`1mn1`$.
If $`b=\widehat{0}`$, then $`\chi ([b,b\pi ],t)=\chi ([\widehat{0},\pi ],t)=\chi (NC_{n1},t)`$. Thus we get the first term of the formula. Now let $`b=mn`$. It is clear that $`b\pi =\widehat{1}`$, so we need to consider the sublattice $`[b,\widehat{1}]`$. Given any $`\omega [b,\widehat{1}]`$, the edge $`mn`$ (which may not be in $`E(G_\omega )`$) geometrically separates the graph $`G_\omega `$ into two parts, $`G_{\omega ,1}`$ and $`G_{\omega ,2}`$, which are induced by vertex sets $`\{1,2,\mathrm{},m,n\}`$ and $`\{m,m+1,\mathrm{},n1,n\}`$, respectively. By contracting the vertices $`m`$ and $`n`$ in both $`G_{\omega ,1}`$ and $`G_{\omega ,2}`$, we get two non-crossing graphs $`\overline{G}_{\omega ,1}`$ and $`\overline{G}_{\omega ,2}`$. It is easy to check that the map $`f:[b,\widehat{1}]NC_m\times NC_{nm}`$ defined by $`f(G_\omega )=(\overline{G}_{\omega ,1},\overline{G}_{\omega ,2})`$ is an isomorphism between these two lattices. Therefore
$$\chi ([b,b\pi ],t)=\chi (NC_m,t)\chi (NC_{nm},t),$$
and the proof is complete.
For any $`\omega =B_1/B_2/\mathrm{}/B_kNC_n`$, the interval $`[\widehat{0},\omega ]_iNC_{|B_i|}`$. Hence to compute the Möbius function of $`NC_n`$, it suffices to do this only for $`\widehat{1}`$. By Proposition 3.1 we have the recurrence relation
$`\mu (NC_n)`$ $`=`$ $`\chi (NC_n,0)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n1}{}}}\chi (NC_i,0)\chi (NC_{ni},0)`$
$`=`$ $`{\displaystyle \underset{i=1}{\overset{n1}{}}}\mu (NC_i)\mu (NC_{ni})`$
with the initial condition $`\mu (NC_1)=1`$. Recall that the Catalan numbers $`C_n=\frac{1}{n+1}\left(\genfrac{}{}{0pt}{}{2n}{n}\right)`$ satisfy the recurrence relation
$$C_n=\underset{i=0}{\overset{n1}{}}C_iC_{n1i}$$
with the initial condition $`C_0=1`$. Therefore, by induction, we obtain Kreweras’ result that
$$\mu (NC_n)=(1)^{n1}C_{n1}.$$
## 4 Shuffle Posets
The poset of shuffles was introduced by Greene , and he obtained a formula for its characteristic polynomial
$$\chi (𝒲_{m,n},t)=(t1)^{m+n}\underset{i0}{}\left(\genfrac{}{}{0pt}{}{m}{i}\right)\left(\genfrac{}{}{0pt}{}{n}{i}\right)\frac{1}{(1t)^i}.$$
In this section we will derive an equivalent formula by using Theorem 2.6. Before doing this, we need to recall some definitions and results of Greene. Let $`𝒜`$ be a set, called the *alphabet of letters*. A *word* over $`𝒜`$ is a sequence $`𝐮=u_1u_2\mathrm{}u_n`$ of distinct letters of $`𝒜`$. We will sometimes also use $`𝐮`$ to stand for the set of letters in the word, depending upon the context. A *subword* of $`𝐮`$ is $`𝐰=u_{i_1}\mathrm{}u_{i_l}`$ where $`i_1<\mathrm{}<i_l`$. If $`𝐮`$, $`𝐯`$ are any two words then the *restriction* of $`𝐮`$ to $`𝐯`$ is the subword $`𝐮_𝐯`$ of $`𝐮`$ whose letters are exactly those of $`𝐮𝐯`$. A *shuffle* of $`𝐮`$ and $`𝐯`$ is any word $`𝐬`$ such that $`𝐬=𝐮𝐯`$ as sets and $`𝐬_𝐮=𝐮`$, $`𝐬_𝐯=𝐯`$ as words.
Given nonnegative integers $`m`$ and $`n`$, fix disjoint words $`𝐱=x_1\mathrm{}x_m`$ and $`𝐲=y_1\mathrm{}y_n`$. The *poset of shuffles* $`𝒲_{m,n}`$ consists all shuffles $`𝐰`$ of a subword of $`𝐱`$ with a subword of $`𝐲`$ while the partial order is that $`𝐯𝐰`$ if $`𝐯_𝐱𝐰_𝐱`$, $`𝐯_𝐲𝐰_𝐲`$ as sets and $`𝐯_𝐰=𝐰_𝐯`$ as words. The covering relation is more intuitive: $`𝐯𝐰`$ if $`𝐰`$ can be obtained from $`𝐯`$ by either adding a single $`y_i`$ or deleting a single $`x_j`$. It is easy to see that $`𝒲_{m,n}`$ has $`\widehat{0}=𝐱`$, $`\widehat{1}=𝐲`$, and is graded by the rank function
$$\rho (𝐰)=(m|𝐰_𝐱|)+|𝐰_𝐲|.$$
For example, $`𝒲_{2,1}`$ is shown in Figure 1 where $`𝐱=de`$ and $`𝐲=D`$.
Every shuffle poset is actually a lattice. To describe the join operation in $`𝒲_{m,n}`$, Greene defined crossed letters as follows. Given $`𝐮`$, $`𝐯𝒲_{m,n}`$ then $`x𝐮𝐯𝐱`$ is *crossed* in $`𝐮`$ and $`𝐯`$ if there exist letters $`y_i`$, $`y_j𝐲`$ with $`ij`$ and $`x`$ appears before $`y_i`$ in one of the two words but after $`y_j`$ in the other. For example, let $`𝐱=def`$ and $`𝐲=DEF`$. Then in the two shuffles $`𝐮=dDEe`$, $`𝐯=Fdef`$, the only crossed letter is $`d`$. The join of $`𝐮`$, $`𝐯`$ is then the unique word $`𝐰`$ greater than both $`𝐮`$, $`𝐯`$ such that
$$\begin{array}{ccc}\hfill 𝐰_𝐱& =& \{x𝐮_𝐱𝐯_𝐱x\text{ is not crossed}\}\hfill \\ \hfill 𝐰_𝐲& =& 𝐮_𝐲𝐯_𝐲.\hfill \end{array}$$
In the previous example, $`𝐮𝐯=DEFe`$. This join also shows that $`𝒲_{m,n}`$ is not semimodular in general, because $`\rho (𝐮)+\rho (𝐯)=3+1<5=\rho (𝐮𝐯)\rho (𝐮𝐯)+\rho (𝐮𝐯).`$ Since $`(𝒲_{n,m})^{}=𝒲_{m,n}`$, the meet operation in $`𝒲_{m,n}`$ is as same as the join operation in $`(𝒲_{n,m})^{}`$. So to find the meet in the analogous way we need to consider those letter $`y𝐮𝐯𝐲`$ crossed in $`𝐮`$ and $`𝐯`$.
Greene also showed that subwords of $`𝐱`$ and subwords of $`𝐲`$ are modular elements of $`𝒲_{m,n}`$. In particular, the empty set $`\mathrm{}`$ is modular. Also note that $`[\widehat{0},\mathrm{}]B_m`$. We now give our formula for the characteristic polynomial of $`𝒲_{m,n}`$.
###### Proposition 4.1
The characteristic polynomial of the shuffle poset is
$$\chi (𝒲_{m,n},t)=(t1)^m\underset{i=0}{\overset{n}{}}(1)^i\left(\genfrac{}{}{0pt}{}{n}{i}\right)\left(\genfrac{}{}{0pt}{}{m+i}{i}\right)t^{ni}.$$
(11)
Proof. Consider any $`𝐮`$ with $`𝐮\mathrm{}=\widehat{0}`$. In general, if $`𝐮\mathrm{}=𝐰`$ then $`𝐰_𝐱=𝐮_𝐱\mathrm{}_𝐱=𝐮_𝐱`$. So $`𝐮\mathrm{}=\widehat{0}`$ if and only if $`𝐱`$ is a subword of $`𝐮`$, i.e., the element $`𝐮`$ is a shuffle of $`𝐱`$ with a subword of $`𝐲`$. Furthermore, for any $`𝐯[\widehat{0},\mathrm{}]`$, there is no crossed letter $`x`$ in $`𝐮`$ and $`𝐯`$ since $`𝐯_𝐲=\mathrm{}`$. It follows that $`(𝐮𝐯)_𝐱=𝐮_𝐱𝐯_𝐱=𝐯`$ and $`(𝐮𝐯)_𝐲=𝐮_𝐲𝐯_𝐲=𝐮_𝐲`$ as sets. Then we get
$`\rho (𝐮𝐯)\rho (𝐮)`$ $`=`$ $`[(m|𝐯|)+|𝐮_𝐲|][(mm)+|𝐮_𝐲|]`$
$`=`$ $`m|𝐯|=\rho (𝐯)\rho (\widehat{0}).`$
Thus the map $`\tau _𝐮:[\widehat{0},\mathrm{}][𝐮,\mathrm{}𝐮]`$ is rank-preserving.
Since $`[\widehat{0},\mathrm{}]B_m`$ we get, by Theorem 2.6,
$$\chi (𝒲_{m,n},t)=(t1)^m\underset{𝐮\mathrm{}=\widehat{0}}{}\mu (𝐮)t^{(m+n)\rho (𝐮)m}.$$
It is easy to see that the interval $`[\widehat{0},𝐮]`$ is isomorphic to $`B_i`$ where $`i=|𝐮_𝐲|`$. So $`\mu (𝐮)=(1)^{|𝐮_𝐲|}=(1)^{\rho (𝐮)}`$. Now we conclude that
$`\chi (𝒲_{m,n},t)`$ $`=`$ $`(t1)^m{\displaystyle \underset{i=0}{\overset{n}{}}}\left[{\displaystyle \genfrac{}{}{0pt}{}{\text{the number of ways to}}{\text{shuffle }𝐱\text{ with }i\text{ letters of }𝐲}}\right](1)^it^{ni}`$
$`=`$ $`(t1)^m{\displaystyle \underset{i=0}{\overset{n}{}}}(1)^i\left({\displaystyle \genfrac{}{}{0pt}{}{n}{i}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{m+i}{i}}\right)t^{ni}.\text{ }`$
To determine the Möbius function of $`𝒲_{m,n}`$, it suffices to compute $`\mu (\widehat{1})`$ since for any $`𝐰𝒲_{m,n}`$ the interval $`[\widehat{0},𝐰]`$ is isomorphic to a product of $`𝒲_{p,q}`$’s for certain $`pm`$ and $`qn`$. Simply plugging $`t=0`$ into formula (11) gives us the Möbius function $`\mu (𝒲_{m,n})`$.
###### Corollary 4.2 (Greene, )
We have
$$\mu (𝒲_{m,n})=(1)^{m+n}\left(\genfrac{}{}{0pt}{}{m+n}{n}\right).\text{ }$$
## 5 NBB Sets and Factorization Theorems
Blass and Sagan derived a Total Factorization Theorem for LL lattices which generalizes Theorem 1.3. Applying Theorem 2.3 and 2.6, respectively, we will offer two inductive proofs for their theorem. First of all, we would like to outline their work.
Given a lattice $`L`$, let $`A=A(L)`$ is the set of atoms of $`L`$. Let $``$ be an arbitrary partial order on $`A`$. A nonempty set $`DA`$ is *bounded below* or *BB* if, for every $`dD`$ there is an $`aA`$ such that
$$ad\text{and}a<D.$$
A set $`BA`$ is called *NBB* (*no bounded below* subset) if it does not contain any $`D`$ which is bounded below. An NBB set is said to be a base for its join. One of the main results of Blass and Sagan’s paper is the following theorem which is a simultaneous generalization of both Rota’s NBC and Crosscut Theorems (for the crosscut $`A(L)`$).
###### Theorem 5.1 (Blass and Sagan, )
Let $`L`$ be a finite lattice and let $``$ be any partial order on $`A`$. Then for all $`xL`$ we have
$$\mu (x)=\underset{B}{}(1)^{|B|}$$
where the sum is over all NBB bases $`B`$ of $`x`$.
Given an arbitrary lattice $`L`$, let $`\mathrm{\Delta }:\widehat{0}=x_0x_1\mathrm{}x_n=\widehat{1}`$ be a maximal chain of $`L`$. The $`i^{\mathrm{𝑡ℎ}}`$ *level* of $`A`$ is defined by
$$A_i=\{aAax_i\text{ but }ax_{i1}\},$$
and we partially order $`A`$ by setting $`ab`$ if and only if $`aA_i`$ and $`bA_j`$ with $`i<j`$. We say $`a`$ is in *lower level* than $`b`$ or $`b`$ is in *higher level* than $`a`$ if $`ab`$. Note that the level $`A_i`$ is an empty set if and only if $`x_i`$ is not an atomic element. A pair $`(L,\mathrm{\Delta })`$ is said to satisfy the *level condition* if this partial order $``$ of $`A`$ has the following property.
$$\text{If }ab_1b_2\mathrm{}b_k\text{ then }a\underset{i=1}{\overset{k}{}}b_i.$$
If all elements of $`\mathrm{\Delta }`$ are left-modular, then we say $`(L,\mathrm{\Delta })`$ is a *left-modular* lattice. A pair $`(L,\mathrm{\Delta })`$ is called an *LL lattice* if it is left-modular and satisfies the level condition.
A generalized rank function $`\rho :L`$ is defined by
$$\rho (x)=\text{ number of }A_i\text{ containing atoms less than or equal to }x\text{.}$$
Note that, for any $`xL`$, we have $`\rho (x)=\rho (\delta (x))`$ where $`\delta (x)`$ is the maximum atomic element in $`[\widehat{0},x]`$. So $`\rho (\widehat{1})`$ is not necessary equal to $`n`$, the length of $`\mathrm{\Delta }`$.
In the following we list several properties in that we need.
1. If $`a`$ and $`b`$ are distinct atoms from the same level $`A_i`$ in a left-modular lattice, then $`ab`$ is above some atom $`cA_j`$ with $`j<i`$.
2. In an LL lattice, a set $`BA`$ is NBB if and only if $`|BA_i|1`$ for every $`i`$.
3. Let $`B`$ be an NBB set in an LL lattice. Then every atom $`aB`$ is in the same level as some element of $`B`$. In particular, any NBB base for $`x`$ has exactly $`\rho (x)`$ atoms.
Blass and Sagan generalized Stanley’s Total Factorization Theorem to LL lattices using their theory of NBB sets. Here we present two inductive proofs for their theorem. In the first proof we will apply Theorems 2.6 as well as the theory of NBB sets.
###### Theorem 5.2 (Blass and Sagan, )
If $`(L,\mathrm{\Delta })`$ is an LL lattice then its characteristic polynomial factors as
$$\chi (L,t)=(t|A_i|)$$
where the product is over all non-empty levels $`A_i`$.
Proof of Theorem 5.2 I. We will induct on $`n`$, the length of $`\mathrm{\Delta }`$. The theorem is trivial when $`n1`$. If $`A_n=\mathrm{}`$, then $`\rho (x_n)=\rho (x_{n1})`$ and $`\mu (x)=0`$ for $`xx_{n1}`$. Thus $`\chi (L,t)=\chi ([\widehat{0},x_{n1}],t)`$, so we are done by induction.
If $`A_n\mathrm{}`$, consider $`bH(L)`$. Then, by Theorem 5.1, $`b`$ must have an NBB base, say $`B`$. In addition, if $`bx_{n1}=\widehat{0}`$ then $`BA_n`$ and also $`|BA_n|1`$ by (B). So $`b=\widehat{0}`$ or $`bA_n`$. Now it suffices to check that $`\tau _b`$ is rank-preserving on $`H(\widehat{0},x_{n1})\{x_{n1}\}`$ for every $`bA_n`$ since then we get $`\chi (L,t)=\chi ([\widehat{0},x_{n1}],t)(t|A_n|)`$ by Theorem 2.6. Because $`A_n\mathrm{}`$ and $`\rho (b)=1`$, $`\tau _b`$ is rank-preserving on $`\{x_{n1}\}`$. Given any $`yH(\widehat{0},x_{n1})`$, suppose $`B`$ be an NBB base for $`y`$. By (B), $`B^{}=B\{b\}`$ is an NBB base for $`\tau _b(y)`$. Now $`\rho (\tau _b(y))=|B^{}|=|B|+1=\rho (y)+\rho (b)`$ by (C). Hence $`\rho (b,\tau _b(y))=\rho (\tau _b(y))\rho (b)=\rho (y)=\rho (\widehat{0},y)`$
In a similar way, Corollary 2.7 provides us with an inductive proof for Theorem 1.3. Note that the lattice in Theorem 1.3 is graded, so $`\rho (\widehat{1})`$ equals the length of $`\mathrm{\Delta }`$. Therefore the product (2) is over all levels $`A_i`$ (including empty ones).
We will use Theorem 2.3 for the second proof. This demonstration sidesteps the machinery of NBB sets and reveals some properties of LL lattices in the process. To prepare, we need the following two lemmas.
###### Lemma 5.3
If $`w`$ is a left-modular element in $`L`$ and $`vw`$, then $`vuwu`$ for any $`uL`$.
Proof. Suppose not and then there exists $`sL`$ such that $`vu<s<wu`$. Taking the join with $`w`$ and using $`vw=w`$, we get $`w(vu)=ws=w(wu)`$. So we should have $`w(vu)<ws<w(wu)=w`$ by Theorem 1.4. Combining this with $`vw(vu)`$, we have a contradiction to $`vw`$
###### Lemma 5.4
If $`(L,\mathrm{\Delta })`$ is an LL lattice with $`\mathrm{\Delta }:\widehat{0}=x_0x_1\mathrm{}x_n=\widehat{1}`$ and $`A_n\mathrm{}`$, then $`([b,\widehat{1}],\mathrm{\Delta }^{})`$ is also an LL lattice for any $`bA_n`$ where $`\mathrm{\Delta }^{}`$ consists of the distinct elements of the multichain
$$b=x_0^{}x_1^{}x_2^{}\mathrm{}x_{n2}^{}x_{n1}^{}=\widehat{1}$$
where $`x_i^{}=x_ib,0in1`$. Furthermore we have $`|A_i|=|A_i^{}|`$ for such $`i`$, where
$$A_i^{}=\{aA(b,\widehat{1})ax_i^{}\text{ but }ax_{i1}^{}\}.$$
Proof. By Lemma 5.3, the chain $`\mathrm{\Delta }^{}`$ is indeed saturated. So $`\mathrm{\Delta }^{}`$ is a left-modular maximal chain by Proposition 1.5.
Let $`\tau (x)=\tau _b(x)=xb`$. This map is surjective (see the proof of Theorem 2.6) and order-preserving from $`[\widehat{0},x_{n1}]`$ to $`[b,\widehat{1}]`$. Also let $`A=A(\widehat{0},x_{n1})`$ and $`A^{}=A(b,\widehat{1})`$. First, We prove that the map $`\tau :AA^{}`$ is well-defined and bijective. Suppose that there is an $`aA_i`$ such that $`bx<\tau (a)=ab`$ for some $`x`$. By the level condition, any atom $`cab`$ is in a level at least as high as $`a`$; furthermore, if $`cA_i`$ we must have $`c=a`$ because of (A). Since $`x<ab`$ and $`ax`$, any atom $`dx`$ is in a higher level than $`a`$. It follows that $`x_ix=\widehat{0}`$. Now $`b(x_ix)=b`$ and $`(bx_i)x(ba)x=x`$ contradicts the left-modularity of $`x_i`$. We conclude that $`\tau :AA^{}`$ is well-defined.
The restriction $`\tau |_A`$ is surjective since $`\tau `$ is surjective and order-preserving. To show injectivity of $`\tau |_A`$, let us suppose there are two distinct atoms $`u`$ and $`v`$ such that $`\tau (u)=\tau (v)`$. If $`u`$ and $`v`$ are from two different levels then this contradicts the level condition. If $`u`$ and $`v`$ are from the same level, by (A), there exists an atom $`c`$ in a lower level such that $`cuv\tau (u)\tau (v)=\tau (u)`$, contradicting the level condition again.
Now let us prove $`|A_i|=|A_i^{}|`$. This is trivial for $`i=1`$. Let $`uA_i`$ for some nonempty $`A_i`$ with $`2in1`$. It is clear that $`\tau (u)x_i^{}`$. Suppose that $`\tau (u)x_{i1}^{}`$, i.e., $`ubx_{i1}b`$. By the level condition, $`b(x_{i1}(ub))=b\widehat{0}=b`$. But $`(bx_{i1})(ub)=ub>b`$ contradicts the modularity of $`x_{i1}`$. Thus $`\tau (A_i)A_i^{}`$ and then the bijectivity of $`\tau |_A`$ implies that $`|A_i|=|A_i^{}|`$ for all $`in1`$.
Since $`\tau |_A`$ is bijective and level-preserving, if $`\tau (a)_{i=1}^k\tau (b_i)`$ for some $`\tau (a)\tau (b_1)\tau (b_2)\mathrm{}\tau (b_k)`$ in $`[b,\widehat{1}]`$, then $`a<ab(_{i=1}^kb_i)b`$ with $`ab_1b_2\mathrm{}b_kb`$ in $`L`$. Therefore $`([b,\widehat{1}],\mathrm{\Delta }^{})`$ satisfies the level condition.
Proof of Theorem 5.2 II. We will induct on $`n=\mathrm{}(\mathrm{\Delta })`$. The cases $`n1`$ and $`A_n=\mathrm{}`$ are handled as before.
If $`A_n\mathrm{}`$, consider $`bH(L)`$ with $`bx_{n1}=\widehat{0}`$. Then $`b`$ is atomic and can only be above atoms in $`A_n`$. So by (A), $`b`$ must be the join of at most one atom, i.e., either $`b=\widehat{0}`$ or $`bA_n`$. Thus by Lemma 5.4 and induction we get, for any $`bA_n`$,
$$\chi ([b,\widehat{1}],t)=\underset{in1}{}(t|A_i^{}|)=\underset{in1}{}(t|A_i|)=\chi ([\widehat{0},x_{n1}],t)$$
where the product is over all non-empty $`A_i`$. Applying Theorem 2.3 gives $`\chi (L,t)=\chi ([\widehat{0},x_{n1}],t)(t|A_n|)`$, so again we are done. |
warning/0001/math0001013.html | ar5iv | text | # 1 Adeles, Ideles and Zeros
## 1 Adeles, Ideles and Zeros
We express the Riemann Hypothesis for abelian L–functions as a Hilbert space closure property (theorem 1.11 below). This takes place within the adelic set-up used by Tate (1950) and Iwasawa (1952) to establish the functional equations of these L–functions. We treat simultaneously the number field and function field cases (the Tate–Iwasawa ideas have been adapted by Weil to the function field case in ). Our approach is Hilbert space-theoretical. We take our hint from Nyman’s equivalent formulation of the original Riemann Hypothesis (1950). Beurling (1949, for the disc), and Lax (1959, for the half-plane), described the invariant subspaces of the Hardy spaces and, as is explained in (see also and for Beurling’s $`L^p`$–extension (1955)), this description is the conceptual element behind Nyman’s thorem. We devote a section to explain (without mention of adeles and ideles) what our construction amounts to for the Riemann zeta function. It is technically of a very straightforward nature, its only deeper aspects being embedded in the Beurling–Lax theory.
We associate to the global field an adelic Lax–Phillips scattering (1967). All axioms (where the idele class group replaces the more usual $``$ or $``$) are satisfied, except possibly the *causality axiom* which we show to be equivalent to the Riemann Hypothesis (this is our main result, theorem 1.7). The validity of one of the axioms is related to an observation of Connes \[9, proof of VIII.5\]. The study of connections between the Riemann zeta function and scattering theory is at least thirty years old. In particular the Faddeev–Pavlov study of scattering for automorphic functions (1972), further developped by Lax and Phillips in their book (1976), has attracted widespread attention. In their approach the scattering matrix is directly related to the values taken by the Riemann zeta function on the line $`Re(s)=1`$, and the Riemann Hypothesis itself is equivalent to some decay properties of scattering waves. Another well-known instance is the approach of De Branges ( 1986, 1994) within the theory of Hilbert spaces of entire functions, also related to scattering (Conrey and Li have recently pointed out some difficulties of this approach (, 1998)). The connection between our scattering process and the Riemann zeta function (or more generally an abelian L–function) is the following: each ‘bad’ zero ($`Re(\rho )>\frac{1}{2}`$) appears as a pole of the scattering operator, where there should be none, if the process was causal. But if the Riemann Hypothesis holds, then the scattering itself is of a trivial nature, and says absolutely nothing on the zeros on the critical line. We point out that the same holds with the positivity criterion of Weil ( 1952, 1972): the Weil distribution if of positive type if and only if the Riemann Hypothesis holds, but beyond that, positivity tells nothing on the location of the zeros except that they are indeed on the critical line. Our formulation applies equally well to function fields and number fields: this is as in Weil’s positivity approach (especially when formulated as in ), and as in the work of Connes (, 1999). The infinite places cause us less trouble than in and . Our sole motivation in formulating the Riemann Hypothesis in a novel manner is the hope that creators of other tools, of a deeper nature than those used here, would incorporate the gained insight in their design constraints. An obvious deficiency of this paper is its inability to achieve an alternative proof of the Riemann Hypothesis in the function field case, where it is not an hypothesis but a well-known theorem.
Let $`K`$ be a global field (an A–field in the terminology of Weil ): either an algebraic number field or a field finitely generated and of transcendence degree $`1`$ over a finite field. We briefly review some normalizations. The adele ring $`𝔸_K`$ is its own Pontrjagin dual. The set of characters (additive, unitary) for which $`K`$ (diagonally embedded) is its own annihilator is non-empty (and a single orbit under the action of $`K^\times `$). We pick one such good character and let the additive Fourier transform $``$ be defined with respect to it (and the corresponding self–dual Haar measure, which is in fact independent of the choice made). On each local multiplicative group $`K_\nu ^\times `$ we write $`d^{}v_\nu `$ for the multiplicative measure which assigns volume $`1`$ to the units (finite place) or is $`\frac{dx}{2|x|}`$ (real place) or $`\frac{drd\theta }{\pi r}`$ (complex place). On the idele group $`𝔸_K^\times `$ (also seen as a subset of $`𝔸_K`$) we use $`d^{}v=_\nu d^{}v_\nu `$, and on the idele class group $`𝒞_K=𝔸_K^\times /K^\times `$ we use the Haar measure $`d^{}u`$ which (function field case) assigns volume $`1`$ to the units or (number field case) is pushed down to $`\frac{dt}{t}`$ under $`t=|u|=|v|=_\nu |v_\nu |_\nu `$ ($`v𝔸_K^\times ,u=\overline{v}`$).
Let $`𝒮(𝔸_K)`$ be the vector space of Bruhat–Schwartz functions.
###### Definition 1.1
$`E:𝒮(𝔸_K)`$ $``$ $`(𝒞_K)`$
$`\phi (x)`$ $``$ $`f(\overline{v})=\sqrt{|v|}{\displaystyle \underset{qK^\times }{}}\phi (qv){\displaystyle \frac{_{𝔸_K}\phi (x)𝑑x}{\sqrt{|v|}}}`$
For functions satisfying the additional conditions $`\phi (0)=_{𝔸_K}\phi (x)𝑑x=0`$, $`E`$ is a tool at the heart of the constructions of Connes in . For technical, class-field theoretical, reasons, we do not impose any vanishing condition. The map $`E`$ is related to the ideas of Tate and Iwasawa , and is especially tuned for Hilbert space matters, as expressed in the following lemma:
###### Lemma 1.2
$`E(𝒮(𝔸_K))L^2(𝒞_K,d^{}u)`$ and is dense in it. The Fourier–Mellin transform of $`E(\phi )`$, as a function of the unitary characters of $`𝒞_K`$, is, up to a multiplicative constant depending only on $`K`$, equal to the Tate L–functions associated to $`\phi `$ (restricted to the critical line).
###### Note 1.3
As has already been noted by Connes \[9, proof of VIII.5\], $`E(𝒮_{00})`$ is dense in $`L^2(𝒞_K,d^{}u)`$, where $`𝒮_{00}=\{\phi 𝒮(𝔸_K)|\phi (0)=_{𝔸_K}\phi (x)𝑑x=0\}`$.
The idele group acts on $`𝒮(𝔸_K)`$ ($`U(v)\phi (x)=\frac{1}{\sqrt{|v|}}\phi (\frac{x}{v})`$) and on $`L^2(𝒞_K,d^{}u)`$ ($`U(v)f(u)=f(\frac{u}{\overline{v}})`$), and $`E`$ intertwines the two actions. Furthermore the Poisson-Tate summation formula shows that $`E`$ intertwines the Fourier transform $``$ on $`𝔸_K`$ with the inversion $`I`$ ($`f(u)f(\frac{1}{u})`$) on $`𝒞_K`$. Each idele $`v`$ defines an adelic parallelepiped
$$P(v)=\left\{x=(x_\nu )𝔸_K\right|\nu |x_\nu |_\nu |v_\nu |_\nu \}$$
whose volume is proportional to $`|v|`$.
###### Definition 1.4
$$𝒮_1=\{\phi 𝒮(𝔸_K)|v𝔸_K^\times :|v|=1\text{ and supp}(\phi )P(v)\}$$
$$\stackrel{~}{𝒮_1}=\{\phi 𝒮(𝔸_K)|v𝔸_K^\times :|v|=1\text{ and supp}((\phi ))P(v)\}$$
$$𝒟_+=E(𝒮_1)^{}$$
$$𝒟_{}=E(\stackrel{~}{𝒮_1})^{}$$
###### Lemma 1.5
The Lax–Phillips scattering axioms (, with $``$ or $``$ replaced with $`𝒞_K`$) are satisfied for the “incoming” subspace $`𝒟_{}`$
$$|\lambda |1U(\lambda )𝒟_{}𝒟_{}U(\lambda )𝒟_{}=\{0\}\overline{U(\lambda )𝒟_{}}=L^2(𝒞_K,d^{}u)$$
and for the “outgoing” subspace $`𝒟_+`$
$$|\lambda |1U(\lambda )𝒟_+𝒟_+U(\lambda )𝒟_+=\{0\}\overline{U(\lambda )𝒟_+}=L^2(𝒞_K,d^{}u)$$
###### Note 1.6
The property $`U(\lambda )𝒟_+=\{0\}`$ is cousin to the density property $`\overline{E(𝒮_{00})}=L^2(𝒞_K,d^{}u)`$ noted by Connes. The property $`\overline{U(\lambda )𝒟_+}=L^2(𝒞_K,d^{}u)`$ is an easy corollary of the Artin–Whaples product formula. As $`𝒟_{}=I(𝒟_+)`$ and as $`I`$ is an isometry which interchanges dilations and contractions, the axioms for $`𝒟_{}`$ and $`𝒟_+`$ are equivalent.
Our main result is:
###### Theorem 1.7 (A causality criterion)
The Riemann Hypothesis holds for all abelian L–functions of $`K`$ if and only if $`𝒟_{}𝒟_+`$.
We also express the Riemann Hypothesis as a closure property. We need a slightly technical definition first:
###### Definition 1.8
Let $`A`$ be the convolution operator
$$(Af)(u_0)=_{𝒞_K}a(\frac{u_0}{u})f(u)d^{}u$$
where, in the number field case
$$a(w)=\sqrt{|w|}\mathrm{𝟏}_{|w|1}$$
and in the function field case ($`q`$ the cardinality of the field of constants)
$$a(w)=\sqrt{|w|}(\sqrt{q}\frac{1}{\sqrt{q}})\mathrm{𝟏}_{|w|<1}+(1\frac{1}{\sqrt{q}})\mathrm{𝟏}_{|w|=1}$$
###### Definition 1.9
$$^2=\left\{fL^2(𝒞_K,d^{}u)\right|\text{ess-supp}(f)\{|u|1\}\}$$
###### Lemma 1.10
The operator $`V=1A`$ is a unitary operator on $`L^2(𝒞_K,d^{}u)`$, commuting with the regular action of $`𝒞_K`$, and sending $`^2`$ to (a subspace of) itself.
###### Theorem 1.11 (A closure criterion)
$`V(\overline{E(𝒮_1)})^2`$ with equality if and only if the Riemann Hypothesis holds for all abelian L–functions of $`K`$.
## 2 The criterion for the Riemann zeta function
When considering only the Riemann zeta function, Theorem 1.11 boils down to a variant of Nyman’s criterion . Let us recall this criterion (see also , , ):
Let $`\rho _\alpha (u)=\left\{\frac{\alpha }{u}\right\}\alpha \left\{\frac{1}{u}\right\}`$, for $`0<\alpha <1`$, and $`u(0,1)`$ (with $`\{\}`$ the fractional part). Let $`N`$ be the closed span in $`L^2((0,1),du)`$ of the functions $`\rho _\alpha `$. We also consider both $`N`$ and $`L^2((0,1),du)`$ as closed subspaces of $`L^2((0,\mathrm{}),du)`$.
###### Theorem 2.1 (Nyman, 1950 )
The constant function 1 on $`(0,1)`$ belongs to $`N`$ if and only if the Riemann Hypothesis holds.
Note that $`N`$ is invariant under the semi-group of unitary contractions $`U(\lambda ):f(u)\sqrt{\frac{1}{\lambda }}f(\frac{u}{\lambda })`$, $`\lambda 1`$, $`u>0`$ as $`U(\lambda )\rho _\alpha =\sqrt{\frac{1}{\lambda }}(\rho _{\alpha \lambda }\alpha \rho _\lambda )`$. So, it will contain the constant function $`\mathrm{𝟏}`$ (hence all step functions) if and only if it actually coincides with all of $`L^2((0,1),du)`$.
For the proof one considers the Mellin transform:
$$f(u)\widehat{f}(s)=_0^1f(u)u^{s1}𝑑u$$
which by a Paley–Wiener theorem establishes an isometry between $`L^2((0,1),du)`$ and the Hardy space $`^2(Re(s)>\frac{1}{2})`$ of analytic functions with bounded norm
$$F^2=\underset{\sigma >\frac{1}{2}}{sup}_{Re(s)=\sigma }|F(s)|^2\frac{|ds|}{2\pi }$$
Such functions $`\widehat{f}(s)`$ have (a.e.) boundary values also obtained as
$$\widehat{f}(\frac{1}{2}+i\tau )=\underset{ϵ0}{\mathrm{l}\mathrm{i}\mathrm{m}}_ϵ^1\sqrt{u}f(u)u^{i\tau }\frac{du}{u}$$
equivalently as the Fourier–Plancherel transform of $`e^{t/2}f(e^t)`$, $`t0`$.
The unitary semi-group considered above acts on $`^2(Re(s)>\frac{1}{2})`$ as $`F(s)\lambda ^{s\frac{1}{2}}F(s)`$, and Lax has described the closed subspaces invariant under this action. It can be directly shown (see ) that the conformal representation
$$w=\frac{s1}{s}$$
$$g(w)=sF(s)$$
establishes an isometry between $`^2(Re(s)>\frac{1}{2})`$ and $`^2(|w|<1)`$ which identifies the invariant subspaces of the former with closed subspaces invariant under the shift $`g(w)wg(w)`$ for the latter. These were described by Beurling and we learn that the “continuous” case (Lax) and “discrete” case (Beurling) are completely equivalent (this equivalence is also a corollary to the conformal invariance of Brownian motion on the complex numbers).
The Beurling–Lax recipe to determine an invariant closed subspace such as $`N`$ is to look at the Mellin transforms of the functions $`\rho _\alpha (u)`$’s:
$$\widehat{\rho _\alpha }(s)=\frac{\alpha \alpha ^s}{s}\zeta (s)$$
and at the “greatest lower bound of their inner factors”: first there will be the Blaschke product
$$B(s)=\underset{\zeta (\rho )=0,Re(\rho )>\frac{1}{2}}{}\frac{s\rho }{s(1\overline{\rho })}\frac{1\overline{\rho }}{\rho }\left|\frac{\rho }{1\rho }\right|$$
where the zeros appear according to their multiplicities, then an inner factor associated to a singular measure on the critical line (the analytic continuation of $`\zeta (s)`$ implies its non-existence), and a final inner factor $`\lambda ^{s\frac{1}{2}}`$ ($`0<\lambda 1`$). We argue that $`\lambda =1`$ as follows: $`\lambda ^{s\frac{1}{2}}^2`$ is the Mellin transform of $`L^2((0,\lambda ),du)`$ which contains $`N`$ only if $`\lambda =1`$ (obviously).
Bercovici and Foias \[3, 2.1\] prove $`\lambda =1`$ in the following manner: if $`\widehat{\rho _\alpha }(s)=\lambda ^{s\frac{1}{2}}f(s)`$ for some $`f(s)^2(Re(s)>\frac{1}{2})`$ then $`\widehat{\rho _\alpha }(\sigma )=O(\lambda ^\sigma )`$ for $`\sigma +\mathrm{}`$. Indeed <sup>1</sup><sup>1</sup>1I thank the referee for correcting my incomplete understanding of the Bercovici–Foias proof at this point. $`f(s)`$ is $`O(1)`$ in any half-plane $`Re(s)\frac{1}{2}+\epsilon `$, $`\epsilon >0`$ (this follows from its Cauchy integral representation or from $`\widehat{f}(s)=_0^1f(u)u^{s1}𝑑u`$ and Cauchy-Schwarz). But obviously $`lim_{\sigma +\mathrm{}}\sigma \widehat{\rho _\alpha }(\sigma )0`$, thus giving a contradiction if $`\lambda <1`$ . The following lemma, of independent interest, could also have been used:
###### Lemma 2.2
If $`F(s)^2`$ is $`O(|s|^K)`$ on the critical line, then its outer factor $`F\mathrm{out}(s)`$ is $`O(|s|^K)`$ on the entire closed half-plane.
###### Proof 2.3
One has
$$\mathrm{log}(|sF\mathrm{out}(s)|)=_{Re(s_0)=\frac{1}{2}}\mathrm{log}(|s_0F(s_0)|)\frac{2Re(s)1}{|ss_0|^2}\frac{|ds_0|}{2\pi }$$
and
$$\mathrm{log}(|s|)=_{Re(s_0)=\frac{1}{2}}\mathrm{log}(|s_0|)\frac{2Re(s)1}{|ss_0|^2}\frac{|ds_0|}{2\pi }$$
hence the result $``$
Let us add a few more words to this discussion of Nyman’s theorem. As
$$_0^1\left\{\frac{1}{u}\right\}u^{s1}𝑑u=\frac{1}{s1}\frac{\zeta (s)}{s}$$
(for $`Re(s)>0`$) and $`\frac{s1}{s}\frac{1}{s1}=\frac{1}{s}=_0^1u^{s1}𝑑u`$ we see that $`\frac{s1}{s}\frac{\zeta (s)}{s}`$ belongs to $`^2(Re(s)>\frac{1}{2})`$. The unitary operator $`V`$ on $`L^2((0,\mathrm{}),du)`$ given by the multiplier $`\frac{s1}{s}`$ in the spectral representation acts as
$$f(u)f(u)_u^{\mathrm{}}\frac{1}{t}f(t)𝑑t$$
As $`\frac{\zeta (s)}{s}=_0^{\mathrm{}}\{\frac{1}{u}\}u^{s1}𝑑u`$ (for $`0<Re(s)<1`$) we obtain after a straightforward computation:
$$\frac{s1}{s}\frac{\zeta (s)}{s}=_0^1A(u)u^{s1}𝑑u$$
$$A(u)=[\frac{1}{u}]\mathrm{log}(u)+\mathrm{log}([\frac{1}{u}]!)+[\frac{1}{u}]$$
Stirling’s formula implies $`A(u)=\frac{1}{2}\mathrm{log}(\frac{1}{u})+O(1)`$ so this integral representation is valid for $`Re(s)>0`$. As $`A(u)=1+\mathrm{log}(u)`$ for $`\frac{1}{2}<u1`$ there is no inner factor of the type $`\lambda ^{s\frac{1}{2}}`$ with $`\lambda <1`$. There is no other singular factor thanks to the analytic continuation, so $`\frac{s1}{s}\frac{\zeta (s)}{s}`$ is the product of an outer factor with the Blaschke product $`B(s)`$. Hence
###### Theorem 2.4
The Riemann Hypothesis holds if and only if $`\frac{s1}{s}\frac{\zeta (s)}{s}`$ is an outer function, or equivalently if the functions $`U(\lambda )A(u)`$ ($`0<\lambda 1`$) span $`L^2((0,1),du)`$.
The generalized Jensen’s formula (see ) then implies a formula first derived by Balazard, Saias and Yor:
###### Theorem 2.5 ()
$$\frac{1}{2\pi }_{Re(s)=\frac{1}{2}}\frac{\mathrm{log}|\zeta (s)|}{|s|^2}|ds|=\underset{\zeta (\rho )=0,Re(\rho )>\frac{1}{2}}{}\mathrm{log}\left|\frac{\rho }{1\rho }\right|$$
The only difference with the proof of Balazard, Saias and Yor is that we do not need the general theory of Hardy spaces beyond that of $`^2`$, which is of a more elementary nature. This concludes our discussion of Nyman’s theorem. We now turn to some variations on this theme (other variations have been considered by Bercovici and Foias in and ).
Let $`\varphi (x)`$ be a smooth function on the real line with compact support in $`[0,1]`$, and $`_0^1\varphi (x)𝑑x=0`$. The Mellin transform
$$\widehat{\varphi }(s)=_0^{\mathrm{}}\varphi (u)u^{s1}𝑑u$$
is an entire function, vanishing at $`1`$. We consider:
$$T(\varphi )(u)=\underset{n1}{}\varphi (nu)(u>0)$$
which is a smooth function of $`u`$ on $`(0,\mathrm{})`$ with support in $`(0,1]`$. Its behavior when $`u0`$ is governed by the Poisson summation formula:
$$T(\varphi )(u)=\frac{1}{|u|}\underset{n}{}\psi (\frac{n}{u})$$
where $`\psi `$ is the Fourier transform $`\varphi (y)e^{2\pi ixy}𝑑y`$ of $`\varphi `$ (hence belongs to the Schwartz space of rapidly decreasing functions). So
$$KT(\varphi )(u)=_{u0}O(u^K)$$
and the Mellin transform
$$\widehat{T(\varphi )}(s)=_0^1T(\varphi )(u)u^{s1}𝑑u$$
is an entire function. For $`Re(s)>1`$
$$\widehat{T(\varphi )}(s)=\zeta (s)\widehat{\varphi }(s)$$
hence by analytic continuation this holds true for all $`s`$.
Let $`𝒮_1^0`$ be the vector space consisting of these functions $`\varphi `$, $`\overline{𝒮_1^0}`$ its closure in $`L^2((0,1),du)`$ and $`K`$ the closure of the vector space of functions $`T(\varphi )`$. Both $`\overline{𝒮_1^0}`$ and $`K`$ are invariant under contractions, hence described by the Beurling–Lax theory. One just has to take the “greatest lower bound” of the inner factors of the $`\widehat{\varphi }(s)`$’s (resp. the $`\widehat{T(\varphi )}(s)`$’s). Obviously $`\overline{𝒮_1^0}`$ is the subspace perpendicular to the constant 1 and this shows that the “greatest lower bound” for the zeros of the $`\widehat{\varphi }(s)`$’s is simply $`s=1`$ with multiplicity $`1`$. This cancels exactly the pole of the zeta function. For the $`\widehat{T(\varphi )}(s)`$’s the analytic continuation across the critical line implies that the only possible singular factor is of the type $`\lambda ^{s\frac{1}{2}}`$ with $`\lambda 1`$. For a suitably chosen $`\varphi `$, $`T(\varphi )`$ does not vanish in $`(\frac{1}{2},1)`$ so necessarily $`\lambda =1`$. The conclusion is that $`K`$ coincides with the space $`N`$ considered by Nyman. Thus:
###### Theorem 2.6
The Riemann Hypothesis holds if and only if the constant function 1 belongs to the closure of $`\{T(\phi ):\phi 𝒮_1^0\}`$
We describe one more variation. Let $`𝒮^{ev}`$ be the vector space of even Schwartz functions on $``$. Let, for $`u>0`$:
$$E(\phi )(u)=\underset{n1}{}\phi (nu)\frac{_0^{\mathrm{}}\phi (x)𝑑x}{u}$$
The Poisson summation formula gives
$$E(\phi )(u)=\frac{1}{u}\underset{n1}{}(\phi )(\frac{n}{u})\frac{1}{2}\phi (0)$$
so that $`E(\phi )(u)`$ is $`0(1)`$ when $`u0`$ and is $`O(\frac{1}{u})`$ when $`u\mathrm{}`$ and belongs to $`L^2(_+,du)`$. Its Mellin transform
$$\widehat{E(\phi )}(s)=_0^{\mathrm{}}E(\phi )(u)u^{s1}𝑑u$$
is absolutely convergent and analytic for $`0<Re(s)<1`$. It can be rewritten as
$$_0^1E(\phi )(u)u^{s1}𝑑u+_1^{\mathrm{}}\underset{n1}{}\phi (nu)u^{s1}du+\frac{_0^{\mathrm{}}\phi (x)𝑑x}{s1}$$
which is then valid in the half-plane $`Re(s)>0`$. Then, for $`Re(s)>1`$, as
$$_0^{\mathrm{}}\underset{n1}{}\phi (nu)u^{s1}du_0^1\frac{_0^{\mathrm{}}\phi (x)𝑑x}{u}u^{s1}𝑑u+\frac{_0^{\mathrm{}}\phi (x)𝑑x}{s1}$$
hence simply as
$$\underset{n1}{}n^s_0^{\mathrm{}}\phi (u)u^{s1}𝑑u=\zeta (s)\widehat{\phi }(s)$$
which remains valid for $`Re(s)>0`$.
We now need to get rid of the pole of $`\zeta (s)`$ with the help of the operator $`V`$ (which on $`L^2(_+,du)`$ acts as $`\frac{s1}{s}`$ in the spectral representation):
$$Vf(u)=f(u)_u^{\mathrm{}}\frac{1}{v}f(v)𝑑v$$
One checks $`V\frac{1}{u}=0`$ so
$$VE(\phi )(u)=\underset{n1}{}\phi (nu)_u^{\mathrm{}}\underset{n1}{}\phi (nv)\frac{dv}{v}=\underset{n1}{}\phi (nu)_0^{\mathrm{}}[\frac{v}{u}]\phi (v)\frac{dv}{v}$$
Let $`𝒮_1`$ be the vector space of smooth even functions with support in $`[1,1]`$. For $`\phi 𝒮_1`$, $`VE(\phi )`$ has support in $`(0,1]`$ and its Mellin transform $`\frac{s1}{s}\zeta (s)\widehat{\phi }(s)`$ thus belongs to $`^2`$. As in the previous discussions, the Mellin transform of the (closure of) $`VE(𝒮_1)`$ is the space of multiples of the Blaschke product $`B(s)^2`$. Hence:
###### Theorem 2.7
$`\overline{VE(𝒮_1)}^2`$ with equality if and only if the Riemann Hypothesis holds.
Let $`B`$ be the unitary operator on $`L^2(_+,du)`$ which acts in the spectral representation as multiplication with $`B(s)`$. Let
$$𝒟_+=E(𝒮_1)^{}=V^1B(^2)^{}=V^1BI^2$$
(where $`I`$ is the inversion $`f(u)\frac{1}{u}f(\frac{1}{u})`$, or spectrally $`s1s`$). Let
$$𝒟_{}=E((𝒮_1))^{}=I(𝒟_+)=IV^1BI^2=VB^1^2$$
Then, in the terminology of Lax and Phillips , $`𝒟_+`$ (resp. $`𝒟_{}`$) is an “outgoing” (resp. “incoming”) space for the action of $`_+^\times `$ on $`L^2(_+,du)`$. The scattering operator associated to them is
$$S=(V^1B)^1VB^1=V^2B^2$$
It is an invariant operator whose spectral multiplier is $`(\frac{s1}{s})^2B(s)^2`$ and is an inner function if and only if $`B(s)`$ has no zero in $`Re(s)>\frac{1}{2}`$, that is if the Riemann Hypothesis holds. The scattering multiplier is inner if and only if $`𝒟_+𝒟_{}`$. So:
###### Theorem 2.8
$`E(𝒮_1)^{}\overline{E((𝒮_1))}`$ if and only if the Riemann Hypothesis holds.
## 3 An adelic scattering
We now prove theorems 1.7 and 1.11. Let $`𝒞_K^1`$ be the (compact) subgroup of idele classes of unit modulus. There is some (non-canonical) isomorphism $`𝒞_K=𝒞_K^1\times N`$, $`N=\{|u|:u𝒞_K\}_+^\times `$. If $`K`$ has positive characteristic we let $`q`$ be the cardinality of the field of constants. It is known that the module group $`N`$ is $`q^{}`$. Each character $`\chi `$ of $`𝒞_K^1`$ extends to a character of $`𝒞_K`$ trivial on $`N`$, which we still denote by $`\chi `$. At each place $`\nu `$ there is a local character $`\chi _\nu `$ from the embedding $`K_\nu ^\times 𝒞_K`$. And $`\chi `$ is said to be ramified at $`\nu `$ if the restriction of $`\chi _\nu `$ to the unit subgroup is non-trivial.
We start with the properties of
$`E:𝒮(𝔸_K)`$ $``$ $`(𝒞_K)`$
$`\phi (x)`$ $``$ $`f(\overline{v})=\sqrt{|v|}{\displaystyle \underset{qK^\times }{}}\phi (qv){\displaystyle \frac{_{𝔸_K}\phi (x)𝑑x}{\sqrt{|v|}}}`$
From the definition one has $`E(\phi )(u)=O(\frac{1}{\sqrt{|u|}})`$ when $`|u|\mathrm{}`$, and as the Poisson-Tate formula gives
$$E=IE$$
one also has $`E(\phi )(u)=O(\sqrt{|u|})`$ when $`|u|0`$. So indeed
$$E(𝒮(𝔸_K))L^2(𝒞_K,d^{}u)$$
Let $`\chi `$ be a unitary character on $`𝒞_K`$ (trivial on $`N`$). The Fourier-Mellin transform (for $`Re(s)=\frac{1}{2}`$)
$$\widehat{E(\phi )}(\chi ,s)=_{𝒞_K}E(\phi (u))\chi (u)|u|^{s\frac{1}{2}}d^{}u$$
is in fact absolutely convergent and analytic for $`0<Re(s)<1`$. It can be rewritten (with $`u=\overline{v}`$, $`v𝔸_K^\times `$) as
$$_{|u|1}E(\phi )(u)\chi (u)|u|^{s\frac{1}{2}}d^{}u+_{|u|>1}\underset{qK^\times }{}\phi (qv)\chi (u)|u|^sd^{}u$$
$$_{𝔸_K}\phi (x)𝑑x_{|u|>1}\chi (u)|u|^{s1}d^{}u$$
The integral $`_{|u|>1}\chi (u)|u|^{s1}d^{}u`$ (which vanishes if $`\chi \mathrm{𝟏}`$) is a meromorphic function $`F_\chi (s)`$, which can be evaluated explicitely. One obtains (both in the number field and function field cases)
$$(Re(s)>1)F_\chi (s)=_{|u|1}\chi (u)|u|^{s1}d^{}u$$
So $`\widehat{E(\phi )}(\chi ,s)`$ has a meromorphic continuation to $`Re(s)>0`$ which, for $`Re(s)>1`$, coincides with
$$_{|u|1}E(\phi )(u)\chi (u)|u|^{s\frac{1}{2}}d^{}u+_{|u|>1}\underset{qK^\times }{}\phi (qv)\chi (u)|u|^sd^{}u$$
$$+_{𝔸_K}\phi (x)𝑑x_{|u|1}\chi (u)|u|^{s1}d^{}u$$
$$=_{𝒞_K}\underset{qK^\times }{}\phi (qv)\chi (u)|u|^sd^{}u$$
$$=C(K)_{𝔸_K^\times }\phi (v)\chi (v)|v|^sd^{}v$$
The constant $`C(K)`$ being as in Tate’s thesis related to the way the measures $`d^{}u`$ on $`𝒞_K`$ and $`d^{}v`$ on $`𝔸_K^\times `$ differ. We recognize in the last integral the Tate L–function $`L(\phi ,\chi ,s)`$. The identity
$$\widehat{E(\phi )}(\chi ,s)=C(K)L(\phi ,\chi ,s)$$
for $`Re(s)=\frac{1}{2}`$ holds by analytic continuation. With this lemma 1.2 is proven.
We turn to the description of $`\mathrm{\Delta }=\overline{E(𝒮_1)}`$. The crucial thing is that it is invariant (obviously) under the (unitary) action of the semi-group of contractions $`\{|u|1\}`$, in particular under the action of the compact group $`𝒞_K^1`$. It thus decomposes as a Hilbert space sum of isotypical components $`\mathrm{\Delta }_\chi `$, which we wish to compare to the isotypical components of $`^2=\left\{fL^2(𝒞_K,d^{}u)\right|\text{ess-supp}(f)\{|u|1\}\}`$. We do this in the spectral representation using the Fourier–Mellin transform (in the function field case we write $`z=q^{(s\frac{1}{2})}`$).
Firstly it is a straightforward check that the $`A`$-operator (1.8) is an invariant operator whose action on $`L^2`$ is given by the following spectral multipliers $`A(\chi ,s)`$:
$$\chi 1A(\chi ,s)=0$$
$$A(1,s)=\frac{1}{s}\text{(number field case)}$$
$$A(1,z)=1\frac{1\sqrt{q}z}{\sqrt{q}z}\text{(function field case)}$$
so that $`V=1A`$ is indeed a unitary (on $`L^2`$) invariant operator with multipliers
$$\chi 1V(\chi ,s)=1$$
$$V(1,s)=\frac{s1}{s}\text{(number field case)}$$
$$V(1,z)=\frac{1\sqrt{q}z}{\sqrt{q}z}\text{(function field case)}$$
From this spectral representation or with a direct computation we also find the important identity
$$V(\frac{1}{\sqrt{|u|}}\mathrm{𝟏}_{|u|>1})=\alpha (K)\sqrt{|u|}\mathrm{𝟏}_{|u|1}$$
with $`\alpha (K)=1`$ (resp. $`\frac{1}{\sqrt{q}}`$) in the number field case (resp. function field case). From the Artin–Whaples product formula we obtain $`E(\phi )(u)=\frac{_{𝔸_K}\phi (x)𝑑x}{\sqrt{|u|}}`$ for $`|u|>1`$ and $`\phi 𝒮_1`$. So we see that $`V(\mathrm{\Delta })`$ is a subspace of $`^2`$. We now describe it exactly with the help of the Beurling–Lax theory.
Let $`S_f`$ be the set of finite places of $`K`$, and $`S_{\mathrm{}}`$ the (possibly empty) set of infinite places. Let $`q_\nu `$ be the cardinality of the residue field at the finite place $`\nu `$, $`\pi _\nu `$ a uniformizer element of $`K_\nu ^\times `$, which we also consider as an element of $`𝔸_K^\times `$. The value $`\chi (\pi _\nu )`$ is independent of the choice of $`\pi _\nu `$ if the character $`\chi `$ is unramified at $`\nu `$. The (“incomplete” in the number field case) L–function associated to $`\chi `$ is
$$L(\chi ,s)=\underset{\nu S_f,\text{unramified}}{}\frac{1}{1\chi (\pi _\nu )q_\nu ^s}$$
The Bruhat-Schwartz function $`\phi `$ is built from local components, all of them except finitely many being equal to the characteristic function of the local integers, so its Tate L–function $`L(\phi ,\chi ,s)`$ is a multiple of $`L(\chi ,s)`$ by a function holomorphic in $`Re(s)>0`$. By lemma 1.2 this implies that the Paley–Wiener transform $`\widehat{E(\phi )}(\chi ,s)`$ ($`Re(s)>\frac{1}{2}`$) vanishes at each bad zero with at least the same multiplicity as $`L(\chi ,s)`$.
###### Definition 3.1
Let $`B`$ be the unitary invariant operator whose spectral multiplier in the $`\chi `$-isotypical component of $`L^2(𝒞_K,d^{}u)`$ is the Blaschke product on the zeros (with multiplicity) of the L–function $`L(\chi ,s)`$ in the half-plane $`Re(s)>\frac{1}{2}`$ (number field case) or the open disc $`|z|<1`$ ($`z=q^{(s\frac{1}{2})}`$, function field case).
We will soon show that one can indeed build a convergent Blaschke product with the bad zeros so that $`B`$ exists! (the function field case is trivial as there are only finitely many). This being temporarily admitted we have obtained $`V(\mathrm{\Delta })B^2`$. And we prove
###### Theorem 3.2
$$V(\mathrm{\Delta })=B^2$$
We treat the function-field case first. We choose $`\phi _\nu `$ to be $`\mathrm{𝟏}_{|x|_\nu 1}`$ at a non-ramified place, and $`\overline{\chi _\nu (x)}\mathrm{𝟏}_{|x|_\nu =1}`$ at a ramified place. With these choices we obtain $`\phi =_\nu \phi _\nu `$ which belongs to $`𝒮_1`$ and for which (at first for $`Re(s)>1`$):
$$L(\phi ,\chi ,s)=L(\chi ,s)$$
We do not claim that $`E(\phi )`$ is $`\chi `$–equivariant, nevertheless this identity combined with lemma 1.2 and the inclusion $`V(\mathrm{\Delta })^2`$ shows that $`V(\chi ,s)L(\chi ,s)`$ belongs to $`^2(|z|<1)`$. It is clear from the product representation that it does not vanish at $`z=0`$, and it is known for $`\chi =1`$ that the pole at $`s=1`$ of the zeta function $`Z_K(s)`$ is of order $`1`$. Analytic continuation across $`|z|=1`$ implies the non-existence of a singular inner factor. So the smallest closed subspace of $`^2(|z|<1)`$ containing $`V(\chi ,s)L(\chi ,s)`$, and invariant under shifts, is exactly $`B(\chi ,s)^2`$. The conclusion follows.
Let us now consider the case where $`K`$ is an algebraic number field. We define $`\phi _\nu (x_\nu )`$ exactly as in the function field case when $`\nu `$ is finite and as $`\overline{\chi _\nu (x)}g_\nu (|x|_\nu )`$ at each infinite place, with $`g_\nu `$ a smooth function on $`_+^\times `$ with compact support in $`(0,1)`$. The product function $`\phi (x)=_\nu \phi _\nu (x_\nu )`$ then belongs to $`𝒮_1`$ and $`E(\phi )`$ has a Paley–Wiener transform
$`{\displaystyle _{𝒞_K}}E(\phi )(u)\chi (u)|u|^{s\frac{1}{2}}d^{}u`$
$`=`$ $`C(K){\displaystyle _{𝔸_K^\times }}\phi (v)\chi (v)|v|^sd^{}v`$
$`=`$ $`C(K)L(\chi ,s){\displaystyle \underset{\nu S_{\mathrm{}}}{}}\widehat{g_\nu }(s)`$
From this and the inclusion $`V(\mathrm{\Delta })^2`$ follows the existence of the Blaschke product $`B(\chi ,s)`$ as promised above. Furthermore it is clearly possible to choose the $`g_\nu `$ in such a manner that $`\widehat{g_\nu }(s)`$ does not vanish at any $`s`$ prescribed in advance, and the existence of analytic continuation accross the critical line then reduces the possibility of an inner factor to $`\lambda ^{s\frac{1}{2}}`$ with $`\lambda 1`$. The Bercovici–Foias argument implies as in our discussion of Nyman’s theorem that $`\lambda =1`$. Finally it is known that the pole of the zeta function ($`\chi =1`$) has exact order $`1`$. With all this the identity $`V(\mathrm{\Delta })=B^2`$ is proven. This completes the proof of the closure criterion 1.11.
Let $`𝒟_+=E(𝒮_1)^{}=\mathrm{\Delta }^{}=V^1B(^2)^{}`$. Let $`Z`$ be the unitary operator which is just $`1`$ in the number field case and $`z`$ (in each isotypical component) in the function field case. Then $`(^2)^{}=Z^1I^2`$ and $`𝒟_+=V^1BZ^1I^2`$. From this follows
$$U(\lambda )𝒟_+=\{0\}\overline{U(\lambda )𝒟_+}=L^2(𝒞_K,d^{}u)$$
so that $`𝒟_+`$ indeed qualifies as an outgoing subspace and $`𝒟_{}`$ as an incoming subspace. One has $`𝒟_{}=IV^1BZ^1I^2=VB^1Z^2`$. The Lax–Phillips scattering operator associated to the pair $`(𝒟_+,𝒟_{})`$ is an invariant unitary operator, unique up to a multiplicative constant in each isotypical component. It is:
$$S=(V^1B)^1VB^1Z=ZV^2B^2$$
With the help of $`S`$ the pair $`(𝒟_+,𝒟_{})`$ is unitarily equivalent to $`((^2)^{},S^2)`$. So it is an orthogonal pair if and only if $`S^2^2`$, if and only if $`B=1`$, if and only if the Riemann Hypothesis holds for all abelian L–functions of $`K`$. With this the proof of the causality criterion 1.7 is complete.
###### Note 3.3
The reader of the monograph of Lax and Phillips \[17, chapter 2\] will perhaps be perplexed by the fact that “causal” means there “inner with respect to the exterior domain $`|z|>1`$” (in the discrete case). But this is because they represent the semi-group leaving invariant the outgoing space with the help of the non-negative powers of $`z`$. In our case we represent it with the help of the non-negative powers of $`\frac{1}{z}`$. So “causal” is to be understood to mean “inner with respect to the domain $`|\frac{1}{z}|>1`$” (that is $`|z|<1`$).
###### Note 3.4
We have used $`IBI=B^1`$. This follows from $`\overline{L(\chi ,\overline{s})}=L(\overline{\chi },s)`$ which implies $`B(\overline{\chi },\overline{s})=\overline{B(\chi ,s)}`$ ($`=B(\chi ,s)^1`$ for $`Re(s)=\frac{1}{2}`$).
Université de Nice – Sophia Antipolis
Laboratoire J.-A. Dieudonné
Parc Valrose
F-06108 Nice Cédex 02
France
burnol@math.unice.fr |
warning/0001/hep-th0001082.html | ar5iv | text | # Untitled Document
hep-th/0001082 PUPT-1909 ITEP-TH-83/99 CALT-68-2252 CITUSC/00-001
Domain Walls and Superpotentials from M Theory on Calabi-Yau Three-Folds
Klaus Behrndt and Sergei Gukov
behrndt@theory.caltech.edu , gukov@theory.caltech.edu
Department of Physics
California Institute of Technology
Pasadena, CA 91125, USA
CIT-USC Center For Theoretical Physics
University of Southern California
Los Angeles, CA 90089-2536, USA
Abstract
Compactification of M theory in the presence of $`G`$-fluxes yields $`𝒩=2`$ five-dimensional gauged supergravity with a potential that lifts all supersymmetric vacua. We derive the effective superpotential directly from the Kaluza-Klein reduction of the eleven-dimensional action on a Calabi-Yau three-fold and compare it with the superpotential obtained by means of calibrations. We discuss an explicit domain wall solution, which represents five-branes wrapped over holomorphic cycles. This solution has a “running volume” and we comment on the possibility that quantum corrections provide a lower bound allowing for an $`AdS_5`$ vacuum of the 5-dimensional supergravity.
1. Introduction
In this paper we study compactification of M theory on Calabi-Yau three-folds in the presence of background $`G`$-fluxes. If there were no $`G`$-fluxes, the effective field theory would be $`𝒩=2`$ five-dimensional supergravity interacting with some number of hypermultiplets and vector multiplets whose scalar fields parametrize a manifold $``$. Turning on a non-trivial $`G`$-flux generates effective superpotential in the five-dimensional theory, which is related to gauging of global isometries of the scalar manifold $``$. If the potential in the five-dimensional theory allows for isolated extrema, the vacuum is given by a space-time of constant negative curvature (i.e. an AdS space) and such a theory is relevant to the AdS/CFT correspondence . On the other hand, we find that the potential is a monotonic function of volume scalar, and this “run-away” case is relevant to the generalization of the AdS/CFT correspondence, the so-called domain wall/QFT correspondence \[2,,3\]. However, if the Calabi-Yau space has positive Euler number then the running of the volume is bounded by quantum corrections , so that five-dimensional supergravity has an AdS<sub>5</sub> vacuum. Below we list various applications which motivated our work.
Over the past year, domain walls as solutions of 5-dimensional gauged supergravity that interpolate between different vacua has been a subject of intensive research; for earlier work on domain wall solution of 4-dimensional supergravity see . Most of them are dealing with the maximal supersymmetric case like in \[6,,7\] and many subsequent papers, but also the least supersymmetric case has been discussed \[8,,9,,10,,11\]. For a recent review see and a discussion that of a running breathing mode is given in \[13\].
A model of domain wall universe was used by Randall and Sundrum to address the hierarchy problem in a novel way, alternative to compactification. An interesting feature of Randall-Sundrum construction is that gravity is localized on the domain wall (or D-brane) by a suitable gravitational potential. The original construction of is purely classical and is based on a non-supersymmetric example of a domain wall which interpolates between two regions of five-dimensional space-time with negative cosmological constant. However, locally AdS form of the vacua on each side of the wall suggests that there must be a corresponding supersymmetric solution. For a related recent work see . Furthermore, motivated by the celebrated D-brane construction of MQCD , one would like to embed a model a la Randall-Sundrum in string theory or M theory to learn about non-perturbative effects in the theory on the domain wall. Even though we will not be able to solve this problem in the full generality, we hope that our study of domain walls constructed from M5-branes wrapped on holomorphic curves inside a Calabi-Yau space will be a useful step in this direction. This configuration has been discussed in the heterotic M-theory compactification in .
Another line of research which motivated this paper is a quest for new supersymmetric vacua in compactifications of string theory or M theory on Calabi-Yau manifolds with background fluxes. Since the flux has to be quantized, its different values correspond to distinct disconnected components in the space of supersymmetric vacua. Therefore, if we call $``$ the background flux and $`_{}`$ the corresponding component of the moduli space, the total space of vacua looks like:
$$=\underset{}{}_{}$$
The component $`_0`$ is equivalent, at least locally, to the moduli space, $`(Y)`$, of the Calabi-Yau space $`Y`$. The other components are isomorphic to some subspaces in the Calabi-Yau moduli space, $`_0(Y)`$, such that all points in $`_{}(Y)`$ correspond to the values of Calabi-Yau moduli which lead to supersymmetric compactifications on $`Y`$ with a given flux $``$. For example, when $`Y`$ is a Calabi-Yau three-fold new supersymmetric vacua can be found at some special (conifold) points of the moduli space \[17,,18\]. For a Calabi-Yau four-fold there is usually more possibility to turn on background fluxes which do not break supersymmetry further \[19,,20,,21,,22\]. Since the value of the flux $``$ jumps across a brane of the appropriate dimension, this brane wrapped over a supersymmetric cycle in $`Y`$ can be identified with a BPS domain wall interpolating between different components in (1.1). This interpretation was used in \[20,,22\] to deduce the effective superpotential $`W()`$ generated by a flux $``$ in compactification on a Calabi-Yau four-fold $`Y`$, such that its minima over $`(Y)`$ reproduce the space of vacua $`_{}`$. Using a more general argument which also applies to compactifications on $`G_2`$ and $`Spin(7)`$ manifolds, one finds the following universal formula for the effective superpotential in terms of calibrations of $`Y`$ :
$$W=_Y(\mathrm{fluxes})(\mathrm{calibrations})$$
The paper is organized as follows. In the next section we perform a Kaluza-Klein reduction of the eleven-dimensional supergravity action on a Calabi-Yau three-fold with a $`G`$-flux. Among other things we find that all supersymmetric vacua of the five-dimensional theory are lifted by the effective superpotential which does not have stable minima. In section 3 we rederive the same result identifying BPS domain walls with five-branes wrapped over holomorphic curves, and argue that the formula (1.1) can be also applied to compactifications on Calabi-Yau three-folds. In section 4 we explicitly construct domain wall solutions in the effective $`D=5`$ $`𝒩=2`$ gauged supergravity which correspond to M5-branes wrapped over holomorphic curves in the Calabi-Yau space. The discussion in section 2 and 4 is in part parallel to the work , which we extend by the inclusion of quantum corrections yielding an AdS vacuum solution.
2. Compactification of M Theory on Calabi-Yau Three-Folds with $`G`$-Fluxes
In this section we perform the compactification of M theory on a Calabi-Yau three-fold $`Y`$ with a $`G`$-flux. For the Kaluza-Klein reduction of the eleven-dimensional action:
$$S_{11}=\frac{1}{2}d^{11}x\sqrt{g}R\frac{1}{2}[\frac{1}{2}GG+\frac{1}{6}CGG]$$
we follow the standard procedure, which lead to gauged $`𝒩=2`$ supergravity in five dimensions as discussed in \[8,,9\]. The latter theory has a potential for the scalar fields $`X^I`$ that play the role of local coordinates on the moduli space of Kähler deformations of $`Y`$. Unfortunately, the scalar potential always exhibits a run-away behavior, so that compactification of M theory on $`Y`$ with non-zero $`G`$-flux does not lead to new vacua in the effective five-dimensional theory. It is worth mentioning that most of the material in this section is not new and has appeared in the literature in various form. In particular, we follow the steps of where analogous compactification on Calabi-Yau three-folds without $`G`$-fluxes was studied. In the context of Type II string theory compactifications with background fluxes were discussed in the work where similar results were found. Closer to the subject of our paper is the work where compactification of M theory with a $`G`$-flux was investigated and the induced superpotential was derived. In order to make the paper self-consistent, below we perform once again all the steps of the Kaluza-Klein reduction in the form that will be convenient later.
The Kaluza-Klein reduction on a Calabi-Yau three-fold $`Y`$ yields $`h^{1,1}`$ abelian gauge fields entering $`h^{1,1}1`$ vector multiplets and a gravity multiplet . The vector fields come from the light modes of the 3-form field $`C`$ in eleven dimensions. Namely, for the field strengths we have a decomposition:
$$GdA^I\omega _I$$
where $`\omega _IH^{(1,1)}(Y)`$ is a basis of $`(1,1)`$-forms. Each vector multiplet contains besides the gaugino a real scalar which comes from the reduction of the internal metric $`g_{a\overline{b}}=it^I(\omega _I)_{a\overline{b}}`$, where $`t^I`$ are the Kähler moduli. Identifying expectation values of the vector multiplet scalar fields with $`t^I`$, we can write the Kähler form as follows:
$$𝒦=t^I\omega _I$$
As we will see below, the scalar parameterizing the volume of the Calabi-Yau decouples from the vector multiplets and enters the universal hypermultiplet. This volume scalar is defined by:
$$𝒱=\sqrt{g_Y}=\frac{1}{3!}𝒦𝒦𝒦=\frac{1}{6}C_{IJK}t^It^Jt^K$$
and the scalars $`\varphi ^A`$ ($`A=1\mathrm{}h^{1,1}1`$) entering the vector multiplets are obtained from
$$1=\frac{1}{6}C_{IJK}X^IX^JX^K\mathrm{with}t^I=𝒱^{\frac{1}{3}}X^I$$
i.e. $`X^I=X^I(\varphi ^A)`$. In what follows we denote $``$ the manifold parameterized by the scalar fields $`\varphi ^A`$, see figure below.
In addition to the volume scalar the universal hypermultiplet contains a real scalar which is dual to the 4-form field in 5 dimensions:
$$GdC_3^{}da$$
and a complex scalar coming from:
$$Gdm\mathrm{\Omega }+d\overline{m}\overline{\mathrm{\Omega }}$$
In addition to these scalars further scalars are related to non-trivial elements of $`H^{(2,1)}(Y)`$, which build up the remaining hyper multiplets. These fields are not important for our analysis, so, we will ignore them.
Fig. 1: Scalar components $`\varphi ^A`$ of vector multiplets parametrize the space $``$ defined by the hypersurface equation (2.1). At extrema of $`W`$, the normal vector $`X_I`$ has to be parallel to the flux vector $`\alpha _I`$.
In order to obtain the canonical Einstein-Hlibert term in 5d, we have to perform a Weyl rescaling (Einstein versus string frame) which is related to the volume of the internal space. In five dimensions this rescaling is given by:
$$ds_E^2=𝒱^{\frac{2}{3}}ds_{str}^2,\sqrt{g_{str}}=\sqrt{g_E}𝒱^{\frac{5}{3}}.$$
with the eleven-dimensional metric written as $`ds_{11}^2=ds_{str}^2+ds_{CY}^2`$. Combining this rescaling with the rescaling of the scalars in (2.1), the reduction of the Ricci scalar yields :
$$S_5=\left[\frac{1}{2}R\frac{1}{2}G_{IJ}(X)X^IX^J\frac{1}{2}\frac{𝒱𝒱}{𝒱^2}\right]$$
where one has to use the relation $`G_{IJ}(X)X^IX^J=0`$ and $`G_{IJ}`$ as function of the $`t^I`$ coordinates is defined by:
$$G_{IJ}(t)=\frac{i}{2𝒱}\omega _I^{}\omega _J=\frac{1}{2}\left[\frac{C_{IJK}t^K}{𝒱}\frac{1}{4}\frac{(C_{IKL}t^Kt^L)(C_{JMN}t^Mt^N)}{𝒱^2}\right]$$
After rescaling into $`X`$ coordinates it takes the form:
$$G_{IJ}(t)=𝒱^{\frac{2}{3}}G_{IJ}(X)$$
with
$$G_{IJ}(X)=\frac{1}{2}\left[C_{IJK}X^K\frac{1}{4}(C_{IKL}X^KX^L)(C_{JMN}X^MX^N)\right].$$
Notice, that $`\frac{3}{2}G_{IJ}X^J=X_I`$ is the normal vector and $`_AX^I`$ are tangent vectors on the scalar manifold $``$, as shown on Fig.1.
In order to perform the reduction of the $`G^{}G`$ term in (2.1), consider the gauge field term (2.1) which in five dimensions becomes:
$$\sqrt{g_{str}}𝒱G_{IJ}(t)F_{\mu \nu }^IF_{\mu ^{}\nu ^{}}^Jg_{str}^{\nu \nu ^{}}g_{str}^{\mu \mu ^{}}=\sqrt{g_E}G_{IJ}(X)F_{\mu \nu }^IF_{\mu ^{}\nu ^{}}^Jg_E^{\nu \nu ^{}}g_E^{\mu \mu ^{}}$$
and (2.1) yields:
$$\sqrt{g_{str}}𝒱(dC_3)^2=\sqrt{g_E}𝒱^2(dC_3)^2.$$
We are looking for potentials that we can obtain from non-trivial $`G`$-fluxes. The flux quantization condition can be written in the following form<sup>1</sup> In general, the periods $`\alpha _I`$ are only required to be half-integer .:
$$_YG_{flux}\omega _I=\alpha _I=\mathrm{integer}$$
Since the internal space remains a Calabi-Yau, in particular Ricci-flat, the only source for a potential comes from the $`G^2`$ term. The topological term contains a derivative in the uncompactified space and therefore cannot give a potential. Let us consider the example discussed in \[8,,9\]:
$$G_{flux}=\frac{1}{𝒱}\alpha ^I{}_{}{}^{}\omega _{I}^{}$$
with $`\alpha ^I=G^{IJ}(t)\alpha _J`$ in agreement with (2.1). This yields:
$$_{M_{11}}G_{flux}^{}G_{flux}=2_{M_5}\sqrt{g_{str}}\frac{1}{𝒱}\alpha ^I\alpha ^JG_{IJ}(t)=2_{M_5}\sqrt{g_E}\frac{1}{𝒱^2}\alpha _I\alpha _JG^{IJ}(X)$$
Note the difference between $`G_{IJ}(t)`$ and $`G_{IJ}(X)`$, cf. (2.1). Writing the volume scalar as:
$$𝒱=e^{2\phi }$$
the potential becomes:
$$V(X,\phi )=e^{4\phi }\left(\alpha _I\alpha _JG^{IJ}(X)\right)$$
A potential of this form was originally found in .
If we include the $`G`$-flux in the topological term we obtain after compactification
$$_{M_{11}}GCG=_{M_5}GA^I_Y\omega _I^{}\omega _J\alpha ^J=_{M_5}GA^I\alpha _I$$
and after dualization the 4-form $`G`$, the corresponding scalar $`a`$ becomes charged under the gauge field $`A^I\alpha _I`$. This effect, as well as the generation of the potential (2.1) can also be understood from the reduction of the eleven-dimensional supersymmetry transformations:
$$\delta e_M^A=i\overline{\eta }\mathrm{\Gamma }^A\psi _M,\delta C_{MNP}=3i\overline{\eta }\mathrm{\Gamma }_{[MN}\psi _{P]}.$$
$$\delta \psi _M=_M\eta \frac{1}{288}(\mathrm{\Gamma }_{M}^{}{}_{}{}^{PQRS}8\delta _M^P\mathrm{\Gamma }^{QRS})G_{PQRS}\eta .$$
Here the supersymmetry parameter $`\eta `$ is an eleven-dimensional Majorana spinor. Under the split 11=5+6, we decompose it as $`\eta =ϵ\xi `$ where $`ϵ`$ is an anti-commuting spinor in five non-compact dimensions.
If there were no background $`G`$-fluxes, then the resulting supersymmetry transformations in five dimensions would correspond to the usual (not gauged) supergravity theory which does not allow a scalar potential. This supergravity would have $`SU(2)`$ $`R`$-symmetry group. Gauging a $`U(1)`$ subgroup of the $`R`$-symmetry group, one obtains a gauged supergravity first found by Gunaydin, Sierra and Townsend . They considered only vector multiplets which effectively means that the volume of the internal space is assumed to be fixed. In this case the supersymmetry transformations in the gauged theory differ only by the extra terms:
$$\delta \lambda ^{iA}=P^A\delta ^{ij}ϵ_j$$
in the gaugino variation, and:
$$\delta \psi _\mu ^i=\frac{i}{2\sqrt{6}}P_0\gamma _\mu \delta ^{ij}ϵ_j$$
in the variation of gravitino. In our notation $`P^A^AW`$ and $`P_0W`$. On the other hand, allowing for general $`G`$-fluxes also yields a dynamical volume scalar which is equivalent to a rescaling of $`W`$ combined with an additional term in the potential $`V`$.
It is easy to see, for example, how (2.1) comes from the supersymmetry transformations (2.1) with a $`G`$-flux. If the background field $`G`$ has non-zero components only in the internal space, then only the first term in brackets is relevant<sup>2</sup> In order to obtain (2.1), we also make a decomposition of gamma-matrices $`\mathrm{\Gamma }_\mu =\gamma _\mu \gamma _7`$ and $`\mathrm{\Gamma }_m=1\gamma _m`$ in the formula (2.1). The eleven-dimensional gamma-matrices $`\mathrm{\Gamma }^M`$ are hermitian for $`M=1,\mathrm{},10`$ and anti-hermitian for $`M=0`$.. The extra term (2.1) can be obtained in a similar way.
To summarize, in the large volume limit the effective five-dimensional gauged supergravity action reads :
$$\begin{array}{cc}\hfill S& \sqrt{G}\left[\frac{1}{2}Rg^2V\frac{1}{4}G_{IJ}F_{\mu \nu }^IF^{\mu \nu J}\frac{1}{2}g_{AB}_\mu \varphi ^A^\mu \varphi ^B\frac{1}{2}h_{rs}D_\mu q^rD^\mu q^s\right]\hfill \\ & +C_{IJK}F^IF^JA^K\hfill \end{array}$$
with
$$g_{AB}=_AX^I_BX^JG_{IJ}$$
subject to the constraint $`\frac{1}{6}C_{IJK}X^IX^KX^J=1`$. Using the convention of the metric of the universal hypermultiplet $`h_{rs}`$ is given by
$$\begin{array}{cc}\hfill h_{rs}dq^rdq^s=\frac{1}{4𝒱^2}d𝒱^2+\frac{1}{2𝒱^2}\left[da+i(md\overline{m}\overline{m}dm)\right]^2+\frac{1}{𝒱}dmd\overline{m}=u\overline{u}+v\overline{v}& \end{array}$$
where
$$u=\frac{dm}{\sqrt{𝒱}},v=\frac{1}{2𝒱}(d𝒱+ida+md\overline{m}\overline{m}dm)$$
and this metric parameterizes the coset $`SU(2,1)/U(2)`$. Recall, $`𝒱`$ is the volume scalar, the axionic scalar $`a`$ comes from the dualization of the five-dimensional 3-form field and the complex scalar $`m`$ was introduced in (2.1). Notice, due to the non-trivial flux only the axionic scalar $`a`$ becomes charged:
$$D_\mu q^r=\{_\mu 𝒱,_\mu a+A_\mu ^I\alpha _I,_\mu m,_\mu \overline{m}\}$$
In the supergravity theory this corresponds to a gauging of the axionic shift symmetry $`aa+const`$. In order to understand the structure of the potential, we have to understand the gauging on the supergravity side . Obviously, the $`G`$-fluxes correspond to a gauging along the Killing vector:
$$k=_a=\frac{i}{2𝒱}(_v_{\overline{v}})$$
The Killing prepotentials have the following form:
$$𝒫_I=\left(\begin{array}{cc}\frac{i}{4𝒱}\alpha _I& 0\\ 0& \frac{i}{4𝒱}\alpha _I\end{array}\right)$$
and obey the relations:
$$k_I^u𝒦_{uv}=_s𝒫_I=_v𝒫_I+[\omega _v,𝒫_I]$$
where $`\omega _v`$ is the $`v`$ component of the SU(2) connection and $`𝒦_{uv}`$ is the triplet of Kähler forms<sup>3</sup> Here $`J^x`$ denotes the triplet of complex structures. $`𝒦_{uv}^x=h_{uw}(J^x)_v^w`$:
$$\omega =\left(\begin{array}{cc}\frac{1}{4}(v\overline{v})& u\\ u& \frac{1}{4}(v\overline{v})\end{array}\right),𝒦=\left(\begin{array}{cc}\frac{1}{2}(u\overline{u}v\overline{v})& u\overline{v}\\ v\overline{u}& \frac{1}{2}(u\overline{u}v\overline{v})\end{array}\right)$$
The gauging fixes the potential to be of the following form:
$$\begin{array}{cc}\hfill V& =4tr(𝒫_I𝒫_J)\left[2X^IX^JG^{IJ}\right]+2X^IX^Jh_{uv}k_I^uk_J^v\hfill \\ & =4e^{4\phi }\left[g^{AB}_AW_BW\frac{4}{3}W^2\right]+2e^{4\phi }W^2|k|^2\hfill \end{array}$$
where
$$W\alpha _IX^I$$
Inserting this expression into (2.1), one finds that the last two terms in the first line cancel and the first term agrees with the potential in (2.1).
So far our discussion was purely classical. However, it is very easy to incorporate corrections due to the non-minimal terms in the action (2.1). It was found by Strominger that corrections due to the terms proportional to the fourth power of the Riemann curvature simply lead to the shift (redefinition of the dilaton field):
$$e^{2\phi }e^{2\phi }+\frac{\chi (Y)}{152^{10}\pi ^8}$$
where $`\chi (Y)`$ is the Euler number of the Calabi-Yau space $`Y`$. In string theory this would be a one-loop correction to the metric on the moduli space of the universal hypermultiplet. We note that as long as we consider only the universal hypermultiplet, we do not expect further corrections especially no instanton corrections. In addition, the shift (2.1) effectively puts a lower bound on the “quantum” volume of the Calabi-Yau space, which leads to some qualitative changes of the supergravity solutions. As we will see in the section 4, the domain wall describes a supergravity solution with monotonically decreasing volume and if it eventually reaches this lower bound, we can keep the volume constant and allowing afterwards only internal deformations as described by the scalars in the vector multiplets. As consequence, at the point where this “quantum” volume is reached, the volume scalar effectively decouples from our supergravity solution and the scalars in the vector multiplets settle down at the extremum of the potential. This configuration is described by an AdS vacuum. Of course, this interpretation makes sense only if $`\chi (Y)>0`$.
3. More Superpotentials From Calabi-Yau Calibrations
In this section we discuss a way to derive the effective superpotentials via identification of BPS domain walls with branes wrapped over supersymmetric cycles. Although in this paper we are mainly interested in M theory compactifications on Calabi-Yau three-folds, we will also consider string theory compactifications.
Let us start with a general compactification of string theory or M theory on a compact oriented manifold $`Y`$ to $`(d+1)`$ non-compact dimensions. In other words, the (real) dimension of space $`Y`$ is equal to $`(9d)`$ in string theory, or $`(10d)`$ in M theory. Trying to keep the discussion as general as possible, we make only a few minor assumptions about the geometry of the space-time. Namely, we assume that compactification on $`Y`$ preserves some supersymmetry, so that it makes sense to talk about BPS domain walls in $`(d+1)`$-dimensional effective theory. Non-compact space-time is assumed to be a maximally symmetric homogeneous space with zero or negative cosmological constant, i.e. Anti de Sitter space or a Minkowski space.
Assuming further the existence of a $`(d+k1)`$-brane in the theory we start with, we can construct a BPS domain wall in the effective field theory by wrapping this brane over a supersymmetric $`k`$-cycle $`\mathrm{\Sigma }Y`$, of course, if there is one. Indeed, a simple counting of dimensions shows that the resulting object should be codimension one in the non-compact space-time. Notice, supersymmetric branes that we consider represent a magnetic source for some field strength in string theory or M theory, depending on the model in question. Let us call this field strength $``$. Thus, as we move across the domain wall in $`(d+1)`$ dimensions, the field strength jumps, $`+\mathrm{\Delta }`$. The change of the flux, $`\mathrm{\Delta }`$, is determined by the geometry of the $`(d+k1)`$-brane that we used to construct the domain wall. Namely, we have:
$$\mathrm{\Delta }=\widehat{[\mathrm{\Sigma }]}$$
where the cohomology class $`\widehat{[\mathrm{\Sigma }]}H^{}(Y,\text{ZZ})`$ is Poincaré dual to the homology class $`[\mathrm{\Sigma }]`$.
Let us now return to the BPS property of the domain wall. Since BPS states have the least possible mass, and the domain wall in question is represented by a $`(d+k1)`$-brane wrapped over $`k`$-dimensional cycle $`\mathrm{\Sigma }`$, we conclude that $`\mathrm{\Sigma }`$ should have the minimal volume in its homology class. Due to this last property, calibrated geometries introduced by Harvey and Lawson turn out to be very useful in a study of supersymmetric brane configurations (see for a review and a list of references). Here we give only the definition of a calibrated submanifold and refer the reader to the original paper for further details. A closed $`k`$-form $`\mathrm{\Psi }`$ is called a calibration if its restriction to the tangent space $`T_x\mathrm{\Sigma }`$ is not greater than the volume form of $`\mathrm{\Sigma }`$ for every submanifold $`\mathrm{\Sigma }Y`$. By saying this we mean that $`\mathrm{\Psi }|_{T_x\mathrm{\Sigma }}\mathrm{vol}(T_x\mathrm{\Sigma })`$ is satisfied provided that $`\mathrm{\Psi }|_{T_xS}=c\mathrm{vol}(T_xS)`$ for some real coefficient $`c1`$. Furthermore, if $`\mathrm{\Psi }|_{T_x\mathrm{\Sigma }}=\mathrm{vol}(T_x\mathrm{\Sigma })`$ for every point $`x\mathrm{\Sigma }`$, the submanifold is called a calibrated submanifold with respect to the calibration $`\mathrm{\Psi }`$. It follows that calibrated submanifolds have the minimal volume in their homology class:
$$\mathrm{Vol}(\mathrm{\Sigma })=_\mathrm{\Sigma }\mathrm{\Psi }$$
The last assumption we are going to make is that the mass of our BPS domain wall in the $`(d+1)`$-dimensional effective theory is determined by the usual BPS formula:
$$M_{BPS}=|\mathrm{\Delta }W|$$
where $`W`$ is the effective superpotential. Then, combining the formulas (3.1), (3.1) and (3.1) together we obtain the following formula for the superpotential generated by a flux $`H^{}(Y)`$:
$$W=_Y\mathrm{\Psi }$$
The approach via calibrated geometries that we have outlined above can be applied to a computation of tree-level superpotentials induced by background fluxes in compactifications of string theory and M theory on Calabi-Yau manifolds \[20,,22,,28\], and to the derivation of membrane instanton superpotentials in M theory compactifications on $`G_2`$ manifolds . Although all the results agree with what one finds studying the supersymmetry conditions, it would be also interesting to derive the effective superpotentials directly from the Kaluza-Klein reduction of the Lagrangian, cf. . It is clear that in the case of Calabi-Yau four-folds, non-minimal terms like the anomaly term $`CI_8(R)`$ and, perhaps, their supersymmetric completion must play an important role .
Compactifications on $`Spin(7)`$ manifolds preserve only two real supercharges in the effective field theory. It was demonstrated in that the BPS mass condition (3.1) is modified in such theories by the one-loop quantum anomaly $`WW+\frac{W^{\prime \prime }}{4\pi }`$. Therefore, one might expect that the effective superpotential induced by a four-form flux in compactification on $`Spin(7)`$ manifold is given by the appropriate modifications of the formula (3.1) which takes into account one-loop quantum anomaly. It would be interesting to see this anomaly by a direct computation of the superpotential via Kaluza-Klein reduction of the ten-dimensional supergravity action or supersymmetry transformations.
In this paper we focus on the case where $`Y`$ is a Calabi-Yau three-fold. There are two types of calibrations on Calabi-Yau three-folds. The first type of calibrations — so-called Kähler calibrations — includes closed forms of even degree constructed from various powers of the Kähler form $`𝒦`$:
$$\mathrm{\Psi }=\frac{1}{p!}𝒦^p$$
Apart from the trivial examples corresponding to $`p=0`$ or $`p=3`$, the submanifolds calibrated by such $`\mathrm{\Psi }`$ are holomorphic curves and divisors in $`Y`$. The second type of calibrations — the special Lagrangian calibration:
$$\mathrm{\Psi }=\mathrm{Re}(\mathrm{\Omega })$$
corresponds to special Lagrangian submanifolds in $`Y`$. Here $`\mathrm{\Omega }H^{3,0}(Y)`$ is the unique holomorphic 3-form.
Clearly, the formula (3.1) can be used in compactifications of heterotic string theory on Calabi-Yau three-folds. In this case, four-dimensional effective field theory has $`𝒩=1`$ supersymmetry. The only way to construct a BPS domain wall in four non-compact dimensions is to consider a five-brane wrapped around a special Lagrangian cycle in $`Y`$. Since a five-brane is a source for the Neveu-Schwarz three-form field strength, eq. (3.1) yields:
$$W=_YH\mathrm{\Omega }$$
We believe that this formula can be derived by the direct arguments, similar to what we used in the previous section.
It turns out that the formula (3.1) can be also applied to theories with larger supersymmetry. Recently, Taylor and Vafa studied the effect of background fluxes in Type II string theory on (non-compact) Calabi-Yau three-folds. They found that it leads to partial supersymmetry breaking via generation of the effective superpotential (3.1) and reconciled it with the results of \[17,,18\]. In particular, in Type IIA string theory on a Calabi-Yau three-fold $`Y`$ the effective superpotential induced by the flux $``$ has the following form :
$$W=_Ye^𝒦$$
which is exactly what follows from (3.1) with the Kähler calibration (3.1).
One goal of the present paper is to demonstrate that effective superpotential of the form (3.1) is also generated in compactification of M theory on a Calabi-Yau space $`Y`$ with a four-form field flux $`G`$. The resulting field theory in five dimensions has $`𝒩=2`$ local supersymmetry. In $`𝒩=2`$ five-dimensional gauged supergravity theories all our assumptions, including the BPS formula (3.1), are justified by the relation between central charge of $`𝒩=2`$ supersymmetry algebra and the gravitino mass as discussed in \[5,,11\]. Since the four-form field strength $`G`$ is the only possible flux in M theory, the formula (3.1) predicts the following simple superpotential $`W_Y𝒦G_{flux}=\alpha _It^I`$. In the five-dimensional $`𝒩=2`$ gauged supergravity theory we expect the effective superpotential $`W`$ to be a function of the scalar fields $`X^I`$ from vector multiplets, rather than $`t^I`$ which also include a volume scalar $`𝒱`$ from the universal hypermultiplet. Since $`\alpha _I`$ are integer numbers, after the appropriate rescaling we obtain the following superpotential:
$$W=\alpha _IX^I$$
which is nothing but the effective superpotential (2.1) found in the previous section via direct Kaluza-Klein reduction.
Notice that variation of the potential (3.1) with respect to the fields $`X^I`$ leads to the condition:
$$G=0$$
which means that there are no supersymmetric vacua in compactification of M theory on Calabi-Yau three-folds with non-trivial fluxes. In other words, the space of supersymmetric vacua has only one component corresponding to $`=0`$, cf. (1.1).
4. Domain Wall Solutions
Motivated by , one may hope to understand non-perturbative effects in realistic models a la Randall-Sundrum via embedding the corresponding domain wall solutions in M theory or string theory. Since $`𝒩=2`$ five-dimensional supergravity can be obtained from compactification of M theory on a Calabi-Yau three-fold $`Y`$, it is natural to assume that the domain wall is constructed out of M5-brane wrapped over a holomorphic curve $`\mathrm{\Sigma }Y`$, see also . Then, topology of $`Y`$ and $`\mathrm{\Sigma }`$ determine the spectrum of the low-energy theory on the five-brane, and the appropriate embedding $`\mathrm{\Sigma }Y`$ may give us a theory close to the Standard Model. Note, because the curve $`\mathrm{\Sigma }`$ is holomorphic in $`Y`$, the effective four-dimensional theory has $`𝒩=1`$ supersymmetry.
Interested in domain wall solutions in the five-dimensional supergravity we write the metric as:
$$ds^2=e^{2U(y)}\left[dt^2+dx_1^2+dx_2^2+dx_3^2\right]+e^{2\gamma U(y)}dy^2$$
where the constant $`\gamma `$ fixes the coordinate system and will be chosen later. This ansatz contains no restrictions as long as we regard the four-dimensional domain wall as a flat Minkowski space, but this parameterization will enable us to find an analytic solution below. Keeping the flat Minkowski space means also that the solution cannot carry electric and/or magnetic charges, but can carry a topological charge given by the difference of the cosmological constants. It is thus consistent to set all the gauge fields to zero. Moreover, investigating the equations of motion coming from the Lagrangian, we find that the complex scalar $`m`$ and the axion $`a`$ can be neglected because they do not show up in the potential. We will keep all scalars $`t^I`$, i.e. the scalars of in the vector multiplets $`\varphi ^A`$ and the volume scalar $`𝒱=e^{2\phi }`$.
4.1. Solution of the $`5d`$ Killing spinor equations
To ensure supersymmetry we have to solve the Killing spinor equations. Since the gauge fields are trivial for our domain wall the relevant variations are:
$$\begin{array}{cc}\hfill \delta \psi _\mu =& \left(_\mu +\frac{1}{4}\omega _\mu ^{ab}\mathrm{\Gamma }_{ab}+\frac{1}{2}g\mathrm{\Gamma }_\mu e^{2\phi }W\right)ϵ,\hfill \\ \hfill \delta \lambda _A=& \left(\frac{i}{2}g_{AB}\mathrm{\Gamma }^\mu _\mu \varphi ^B+i\frac{3}{2}ge^{2\phi }_AW\right)ϵ,\hfill \\ \hfill \delta \zeta =& e^{2\phi }\left(\frac{i}{2}\mathrm{\Gamma }^\mu _\mu e^{2\phi }i\mathrm{\hspace{0.17em}3}gX^Ik_I\right)ϵ\hfill \end{array}$$
with $`W=\alpha _IX^I`$. The scalar fields $`\varphi ^A`$ parameterize the manifold $``$ defined by (2.1), and the only non-trivial hypermultiplet field $`e^{2\phi }=𝒱`$ gives the Calabi-Yau volume.
Let us start with the gravitino variation $`\delta \psi `$. For our ansatz of the metric, the only non-zero components of the vielbeine and spin connection are:
$$e^m=e^Udx^m,e^y=e^{\gamma U}dy,\omega ^{my}=e^{(\gamma +1)U}U^{}dx^m$$
where $`m=0,1,2,3`$ and the corresponding gravitino variation becomes:
$$\delta \psi _m=\left(\frac{1}{2}e^{(\gamma +1)U}U^{}\mathrm{\Gamma }_m\mathrm{\Gamma }_y+\frac{1}{2}\mathrm{\Gamma }_mge^Ue^{2\phi }W\right)ϵ$$
Using the projector $`(1+\mathrm{\Gamma }_y)ϵ=0`$ we find:
$$ge^{2\phi }W=e^{\gamma U}U^{}.$$
From the $`\delta \psi _y`$ component:
$$0=\delta \psi _y=\left(_y+\frac{1}{2}e^{\gamma U}\mathrm{\Gamma }_yge^{2\phi }W\right)ϵ$$
we obtain the Killing spinor after using (4.1):
$$ϵ=e^{\frac{U}{2}}\left(1\mathrm{\Gamma }_y\right)ϵ_0$$
where $`ϵ_0`$ is any constant spinor.
Moreover, using (4.1) we can also solve the hyperino variation:
$$\begin{array}{cc}\hfill 0& =\left(\frac{1}{2}\mathrm{\Gamma }^\mu _\mu e^{2\phi (y)}3gW\right)ϵ\hfill \\ & =\left(\frac{1}{2}e^{\gamma U}(e^{2\phi })^{}3e^{\gamma U}U^{}e^{2\phi }\right)ϵ\hfill \\ & =\frac{1}{2}e^{\gamma U}e^{2\phi }\left(2\phi ^{}6U^{}\right)ϵ\hfill \end{array}$$
and therefore
$$e^{6U}=e^{2(\phi \phi _0)}=𝒱/\mathrm{}$$
where $`\mathrm{}=e^{2\phi _0}`$ is the integration constant. Finally, we come to the gaugino variation $`\delta \lambda _A`$ which gives:
$$\begin{array}{cc}\hfill 0& =\frac{i}{2}\left(g_{AB}\mathrm{\Gamma }^\mu _\mu \varphi ^B3ge^{2\phi }_AW\right)ϵ\hfill \\ & =\frac{i}{2}\left(\mathrm{\Gamma }^\mu _AX^I_BX^JG_{IJ}_\mu \varphi ^B3ge^{2\phi }_A(\alpha _IX^I)\right)ϵ\hfill \\ & =\frac{i}{2}_AX^I\left(e^{\gamma U}\frac{3}{2}_yX_I3ge^{2\phi }\alpha _I\right)ϵ.\hfill \end{array}$$
Because $`_AX^I`$ defines tangent vectors, the expression in brackets has to be proportional to the normal vector $`X_I`$:
$$\frac{3}{2}e^{\gamma U}(X_I)^{}3ge^{2\phi }\alpha _I=3e^{\gamma U}U^{}X_I=\frac{3}{\gamma }(e^{\gamma U})^{}X_I$$
the coefficient on the rhs can be verified by contracting the equation with $`X^I`$ and using (4.1). Next, replacing $`e^{2\phi }`$ by employing (4.1) and taking $`\gamma =4`$ we get
$$\frac{1}{2}_y\left(e^{2U}X_I\right)=g\alpha _I/\mathrm{}$$
and thus the solution is
$$X_I\frac{1}{6}C_{IJK}X^JX^K=e^{2U}\frac{1}{3}H_I=e^{2U}\frac{1}{3\mathrm{}}\left(q_I+6g\alpha _Iy\right)$$
where $`q_I`$ are arbitrary constants. This solution agrees with the one derived in , but notice also the close relationship to the attractor equations which extremize the supersymmetry central charge, or in our case, the superpotential $`W`$ and which state that at extrema of $`W`$ the normal vector $`X_I`$ becomes parallel to the flux vector $`\alpha _I`$. These extrema are reached at $`y\pm \mathrm{}`$ where the scalars $`\varphi ^A`$ becomes constant and $`W`$ extremal, due to (4.1). Remember, because of the run-away behavior of the volume, extrema of $`W`$ are not extrema of the supergravity potential $`V`$.
As we discussed at the end of section 2, quantum corrections yield a lower bound for the volume, which is mainly given by the Euler number of the Calabi-Yau space. So, if we assume that in this “quantum” region the universal hypermultiplets effectively decouples and if we approximate the volume by the lower bound $`𝒱=e^{2\phi _0}=\mathrm{}`$, we find the same solution for eq. (4.1), but with $`\gamma =+2`$. In this case the spacetime metric becomes asymptotically anti de Sitter, which is expected because for a fixed volume, the potential has extrema.
4.2. Domain walls from five-branes in $`T^6`$
In the last section we showed that the Killing spinor equations are solved if the supergravity fields satisfy the eqs. (4.1) and (4.1) with the metric ansatz given by (4.1). Following , let us consider a simple example $`Y=T^6`$, where the intersection form is given by:
$$\frac{1}{6}C_{IJK}X^IX^JX^K=X^1X^2X^3$$
For this example the equations (4.1) become:
$$\begin{array}{cc}& H_1=e^{2U}(X^2)(X^3)\hfill \\ & H_2=e^{2U}(X^1)(X^3)\hfill \\ & H_3=e^{2U}(X^1)(X^2)\hfill \end{array}$$
and thus:
$$X^1=\frac{e^{2U}}{H_1},X^2=\frac{e^{2U}}{H_2},X^3=\frac{e^{2U}}{H_3},e^{6U}=H_1H_2H_3$$
For the generic case of a “running volume” ($`\gamma =4`$) the domain wall metric reads:
$$ds^2=(H_1H_2H_3)^{1/3}\left[dt^2+dx_1^2+dx_2^2+dx_3^2\right]+(H_1H_2H_3)^{4/3}dy^2$$
and the volume is $`𝒱=e^{6U}=H_1H_2H_3`$ (setting $`\mathrm{}=1`$). In order to understand the domain wall from the M theory perspective, we can rescale the solution and obtain for the string frame and for the Kähler class moduli:
$$\begin{array}{cc}& ds_{str}^2=(H_1H_2H_3)^{1/3}\left[dt^2+dx_1^2+dx_2^2+dx_3^2\right]+(H_1H_2H_3)^{2/3}dy^2\hfill \\ & t^I=𝒱^{1/3}X^I=e^{2U}X^I=\frac{(H_1H_2H_3)^{2/3}}{H_I}\hfill \end{array}$$
In an infinite volume limit we can decompactify this solution and the 11-d metric becomes
$$ds_{11}^2=\frac{1}{(H_1H_2H_3)^{1/3}}[dt^2+dx_1^2+dx_2^2+dx_3^2+(H_2H_3d\omega _1+cycl.)+H_1H_2H_3dy^2]$$
where $`d\omega _{1,2,3}`$ are 2-d line elements and this configuration is an intersection $`M5\times M5\times M5`$ over a common 3-brane.
On the other hand in the fixed volume case, the $`X^I`$ field are the same, but since $`\gamma =2`$ the metric differs
$$ds^2=(H_1H_2H_3)^{1/3}\left[dt^2+dx_1^2+dx_2^2+dx_3^2\right]+(H_1H_2H_3)^{2/3}dy^2$$
which yields $`AdS_5`$ for large $`y`$.
4.3. Discussion of some global aspects
Solving the local supergravity equations is not enough to describe domain walls, which are typically gravitational kink solutions that interpolate between vacua at $`y=\pm \mathrm{}`$. Interesting cases are interpolating solutions between vacua with different cosmological constants on both sides, i.e. the scalar fields flow between extrema of the potential. But there are also dilatonic domain walls, where the potential typically does not allow for isolated extrema and at least one scalar field “runs away”. This resembles the linear dilaton vacua appearing in certain string backgrounds. In the compactified theory, this run-away behavior signals a strong or weak coupling region, where the internal volume either diverges or shrinks. This is exactly the case for the solution that we described, where the volume of the internal space as described by the volume scalar $`\phi `$ diverges for $`y+\mathrm{}`$. On the other hand, the scalars in the vector multiplets are fixed by the attractor equations (4.1) and become constant asymptotically, fixed only by the flux vector $`\alpha _I`$. Therefore, the Killing spinor equations imply that asymptotically $`_AW=0`$ and we reach an extremum of the superpotential $`W(\varphi ^A)`$.
Note, the supergravity solution for the scalars $`X^I(y)`$ describes a trajectory on the (curved) moduli space and since it solves the equations of motion this trajectory is geodesic with radial coordinate $`y`$ as affine parameter. This geodesic is fixed if we fix the two endpoints, i.e. two vacua. Let us stress, that it is not enough to fix only one endpoint, say, at $`y=+\mathrm{}`$, we have also to choose where the solution should flow at $`y=\mathrm{}`$, i.e. on the other side of the wall.
But what happens if we pass the point $`y=0`$? Our solution is valid for all values of $`y`$ and the point $`y=0`$ is generically not singular. So, we have to discuss the continuation to negative values of $`y`$. Because the domain wall is an intersection of five-branes, the flux-vector $`\alpha _I`$ should change at least the sign while passing the five-branes. As consequence, the product $`\alpha _Iy`$ remains positive and we avoid singularities due to zeros of harmonic functions<sup>4</sup> Note, a vanishing harmonic function $`H_I`$ means a vanishing cycle $`X_I`$. at finite values of $`y`$. A trivial possibility is to treat both sides symmetrically and therewith identifying both asymptotic vacua. More interesting, especially from the RG-flow point of view, is to patch together different vacua. An interesting case would be a solution interpolating between different vacua of a given superpotential $`W`$, but these domain walls are expected to be singular, because due to the global convexity of the moduli space the attractor equations (4.1) have only one solution for a given Kaehler cone and different extrema of $`W`$ have to lie on disconnected branches of the moduli space. One can also consider a domain wall describing the flow towards vanishing volume. In this case the metric develops a singularity where the 4-dimensional world volume is squeezed to zero size, for examples see \[35,,11,,36\]. Let us comment on them in more detail.
The first thing to notice is that we can always approach this singular point by a proper choice of the vector $`q_I`$, which fixes the point $`X^I(y=0)`$. A vanishing volume of the internal space yields always a singularity in supergravity solutions, but as we discussed earlier, quantum corrections or better higher curvature corrections provide a cut-off for the volume. This lower bound (2.1) was basically given by the Euler number of the internal manifold and therefore the regular supergravity solution can allow for at most 4 unbroken supercharges. By a simple shift in $`y`$ we can always arrange that we reach this “quantum volume” at $`y=0`$ and it is natural to describe the other side of the wall by the solution with a fixed volume, i.e. $`\gamma =2`$ in the solution described in section 4.1. Therefore in this regularized supergravity solution, the volume scalar flows from infinity (infinite volume) towards a lower bound and the scalars in the vector multiplets extremize on both sides the superpotential, i.e. they flow between fixpoints. Notice, the superpotential does not need to be the same on both sides, e.g. we may change the flux vector $`\alpha _I`$ on both sides and/or the intersection form but, due to the attractor equation combined with the convexity of $``$, a given superpotential $`W(\varphi ^A)`$ has a unique extremum, where the flux vector $`\alpha _I`$ is parallel to the normal vector $`X_I`$ (see Fig. 1). Moreover this type of domain wall provides an interesting example from AdS/CFT perspective, because the gauge theory couplings which are dual to the Kaehler class moduli $`t^I`$ are UV-free, related on the sugra side to the infinite volume region, and flow in the IR to a non-trivial conformal fixpoints, where the supergravity solution becomes $`AdS_5`$ with fixed scalars.
Acknowledgements
We would like to thank K. & M. Becker, J. Gomis, C. Vafa and E. Witten for useful discussions. The work of K.B. was supported by a Heisenberg Fellowship of the DFG. The work of S.G. was supported in part by the Caltech Discovery Fund, grant RFBR No 98-02-16575 and Russian President’s grant No 96-15-96939.
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warning/0001/nucl-th0001042.html | ar5iv | text | # Proton decay of high-lying states in odd nuclei
## I Introduction
Nucleon transfer reactions induced by hadronic probes at intermediate energies favour the excitation of high angular momentum states lying above particle emission threshold . The observed structures originate from the coupling of the initial single-particle mode with more complex states. The coupling of single-particle states with surface vibrations is mainly responsible for the damping of the single-particle mode . The particle decay of highly excited states gives the opportunity to study in detail the damping process. For example, the relative contributions of the direct and statistical components to the damping of single-particle mode can be found. Up to now, experimental data and theoretical calculations are available mainly for the neutron decay of high-lying single-particle modes. Very recently, the proton decays of high-lying states in <sup>41</sup>Sc, <sup>59</sup>Cu and <sup>91</sup>Nb have been measured .
In odd nuclei, the simplest excited states can be described as admixtures of single-particle states or quasiparticle states coupled to collective excitations (or phonons) of the even-even core. This weak coupling picture has been successfully applied to obtain the strength functions of a variety of odd nuclei . A method for calculating particle escape widths in the framework of the quasiparticle-phonon model (QPM) has been suggested in Ref. where the inclusive semi-direct neutron decays of high-lying states in <sup>209</sup>Pb have been studied using the general procedure of Refs. which was established for the case of nucleon emission from giant resonances. In a previous work we have extended the method of Ref. in order to calculate non-statistical particle decays of excited states in odd nuclei leading to exclusive channels which correspond to the ground and low-lying excited states of the even-even core. In the present work we apply our method to study the proton decay. Using the QPM we calculate the partial cross sections and branching ratios for the proton decay of the high angular momentum states in <sup>41</sup>Sc, <sup>59</sup>Cu and <sup>91</sup>Nb and compare them with experimental data.
This paper is organized as follows: in Sec. II we describe briefly our theoretical approach to treat the direct nucleon decay of high angular momentum states of single-particle type. In Sec. III a comparison of the calculated and measured branching ratios for non-statistical proton decay of high-lying states in <sup>41</sup>Sc, <sup>59</sup>Cu and <sup>91</sup>Nb is presented. Finally, in Sec. IV conclusions are drawn.
## II Theory
### A The projection operator method
The projection operator method is a convenient approach to treat problems involving single-particle continua in a many-body context. The general formalism was introduced by Yoshida and Adachi and it has been applied to studies of high-lying states in odd nuclei. The detailed expressions can be found in Ref.. Here, we simply recall the main features.
We write the hamiltonian of the A+1 system in the form:
$`H`$ $`=`$ $`h+H_{core}+H_{coupl}.`$ (1)
The first term describes the motion of a particle in an average potential $`U`$ created by the particles in the core:
$`h`$ $`=`$ $`{\displaystyle \frac{1}{2m}}^2+U.`$ (2)
The core hamiltonian is a sum of single-particle hamiltonians $`h_i`$ and two-body residual interactions $`V_{i,j}`$:
$`H_{core}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}h_i+{\displaystyle \underset{i<j}{\overset{A}{}}}V_{i,j}.`$ (3)
The last term of $`H`$ is a sum of interactions between the odd particle and the core particles:
$`H_{coupl}`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{A}{}}}V_{0,i}.`$ (4)
The physical spectrum of $`h`$ consists of a small number of bound states $`\{\phi _i,e_i\}`$ and a continuum of scattering states $`\{\phi _e,e\}`$ which form altogether a complete, orthogonal basis. The projection operator method consists in introducing another complete set of orthogonal basis states which is a direct sum of two complementary subsets, a first subset of discrete, orthonormal states $`\{\varphi _\alpha ,ϵ_\alpha \}`$ which span the single-particle $`q`$ space and a second subset of $`\{\varphi _ϵ,ϵ\}`$ continuum states spanning the complementary $`p1q`$ space. We will denote by $`a_\alpha ^{},a_\alpha `$ ($`a_ϵ^{},a_ϵ`$) the creation and annihilation operators of state $`\phi _\alpha `$ ($`\phi _ϵ`$).
In the spirit of the QPM the hamiltonian $`H_{core}`$ is treated in the random-phase approximation (RPA) in a discrete space, i.e., the particle-hole configurations of RPA are built only with q-space states. We denote by $`E_\nu `$ and $`O_\nu ^{}`$ the energies and creation operators of these RPA states which describe core excitations. If $`0>`$ represents the RPA ground state of the core, the properties of the (A+1)-nucleus except for its nucleon decay properties can be described in terms of the one-particle states $`a_\alpha ^{}0>`$ and one-particle-plus-phonon states $`[a_\beta ^{}O_\nu ^{}]0>`$. We can write:
$`d_i>`$ $``$ $`d_i^{}0>`$ (5)
$`=`$ $`\left({\displaystyle \underset{\alpha }{}}C_\alpha ^{(i)}a_\alpha ^{}+{\displaystyle \underset{\beta ,\nu }{}}D_{\beta ,\nu }^{(i)}[a_\beta ^{}O_\nu ^{}]\right)0>.`$ (6)
We call Q space the space spanned by the (real) state vectors $`d_i>`$ and $`Q`$ the corresponding projection operator. The amplitudes $`C_\alpha ^{(i)}`$ and $`D_{\beta ,\nu }^{(i)}`$ , and the energies $`\omega _i`$ of $`d_i>`$ are determined by diagonalizing $`H`$ in the RPA, i.e., one solves:
$`[QHQ,d_i^{}]=\omega _id_i^{},`$ (7)
within the approximation of commutator linearization. The distribution of $`C^{(i)}^2`$ represents the strength function from which one can deduce the spectroscopic factors.
To allow for nucleons to decay it is necessary to introduce state vectors where the odd particle has a non-zero probability of being at infinity. This is achieved by constructing the P space complementary to Q space and consisting of all states which are linear combinations of the following one-particle and one-particle-plus-phonon configurations:
$`ϵ>a_ϵ^{}0>,ϵ,\nu >[a_ϵ^{}O_\nu ^{}]0>.`$ (8)
The present definition of P space neglects continuum effects on the phonons $`O_\nu ^{}`$ which can also in principle couple to the continuum and emit nucleons by themselves. Actually, these effects should have a small influence on the non-statistical particle decay of the (A+1)-nucleus since the most important phonons contributing to the particle-phonon coupling are the low-lying collective states of the core.
The direct sum of Q space and P space is by construction the complete particle-plus-phonon space in which the hamiltonian $`H`$ should be solved. It is completely equivalent to solve in the simpler Q space the more complicated effective hamiltonian:
$`(E)`$ $``$ $`QHQ+QHP{\displaystyle \frac{1}{E^{(+)}PHP}}PHQ`$ (9)
$``$ $`H_{QQ}+W(E),`$ (10)
where $`P`$ is the projection operator onto P space ($`P+Q=1`$) and $`E`$ is the energy of the system. The hamiltonian $``$ is complex and energy dependent. For each value of $`E`$ one has to find the set of complex states and eigenenergies:
$`𝒟_i`$ $``$ $`𝒟_i^{}0,`$ (11)
$`\mathrm{\Omega }_i`$ $``$ $`\overline{\omega }_i{\displaystyle \frac{i}{2}}\mathrm{\Gamma }_i^{},`$ (13)
satisfying:
$`[(E),𝒟_i^{}]=\mathrm{\Omega }_i𝒟_i^{}.`$ (14)
### B Escape widths
We consider a direct transfer reaction $`a+Ab+(A+1)^{}`$ followed by a sequential decay $`(A+1)^{}p+A^{}`$ where the (A+1)-nucleus in a highly excited state decays by a semi-direct proton emission. In this process, the $`𝒟_i>`$ states will act as doorway states. If we describe the reaction mechanism in a simple approach, e.g., a distorted wave Born approximation (DWBA), we can write down the scattering amplitude from an initial channel $`i`$ where the target $`A`$ is in its ground state $`0>`$ to a final channel $`f`$ where the residual nucleus is left in some excited state $`\nu >=O_\nu ^{}0>`$ with excitation energy $`E_\nu `$ while the escaping proton has an energy $`EE_\nu `$. Using the complex bi-orthogonal basis $`\{𝒟_i>,<\overline{𝒟}_i\}`$, we have:
$`T_{fi}`$ $`=`$ $`{\displaystyle \underset{lj}{}}{\displaystyle \underset{i}{}}{\displaystyle \frac{<\varphi _{lj}^{()}(EE_\nu ),\nu H𝒟_i><\overline{𝒟}_i,bV0,a>}{E\overline{\omega }_i+i\mathrm{\Gamma }_i^{}/2}},`$ (15)
where $`V`$ is the interaction inducing the particle transfer from $`a`$ to $`A`$, $`\varphi ^{()}`$ is an incoming wave of $`p`$ space at energy $`EE_\nu `$, and the sum over $`(l,j)`$ is restricted by the angular momenta and parities of states $`𝒟_i>`$ and $`\nu `$. The case where the final channel is the ground state of the residual nucleus corresponds to the above expression with $`\nu =0,E_\nu =0`$.
The nucleon-transfer matrix element $`<\overline{𝒟}_i,bV0,a>`$ is proportional to the one-quasiparticle amplitude of the state $`𝒟_i>`$. In analogy with the amplitude $`C^{(i)}`$ of Eq.(6), we denote it by $`𝒞^{(i)}`$. We also introduce the partial escape amplitudes of state $`𝒟_i>`$ to channel $`\nu `$:
$`\gamma _{i,\nu }(lj)`$ $``$ $`\sqrt{2\pi }<\varphi _{lj}^{()}(EE_\nu ),\nu H𝒟_i>.`$ (16)
The partial widths are:
$`\mathrm{\Gamma }_{i,\nu }^{}{\displaystyle \underset{lj}{}}\gamma _{i,\nu }(lj)^2.`$ (17)
To obtain an expression for the cross section simple enough to lend itself to a discussion in terms of escape widths, let us assume furthermore that interference terms between different doorway states can be neglected. The density of $`|𝒟_i`$ states is large, and for each interval centered around $`E`$ and containing $`N`$ states we define locally averaged quantities:
$`C^2(E)`$ $``$ $`{\displaystyle \underset{iI}{}}|C^{(i)}|^2/N,`$ (18)
$`\mathrm{\Gamma }_\nu ^{}(E)`$ $``$ $`{\displaystyle \underset{iI}{}}\mathrm{\Gamma }_{i,\nu }^{}/N,`$ (19)
$`\mathrm{\Gamma }^{}(E)`$ $``$ $`{\displaystyle \underset{\nu }{}}\mathrm{\Gamma }_\nu ^{}(E).`$ (20)
Then, one can rewrite the cross sections in the following form:
$`\sigma _\nu (E)C^2(E)\mathrm{\Gamma }_\nu ^{}(E)/\mathrm{\Gamma }^{}(E).`$ (21)
The branching ratios can be calculated by the following formula:
$`B_\nu ={\displaystyle \frac{\sigma _\nu (E)}{_\mu \sigma _\mu (E)}}={\displaystyle \frac{\mathrm{\Gamma }_\nu ^{}(E)}{\mathrm{\Gamma }^{}(E)}}.`$ (22)
### C Inputs of the model
The above formalism is applied to study the semi-direct proton decay of the nuclei <sup>41</sup>Sc, <sup>59</sup>Cu and <sup>91</sup>Nb with the aim of comparing the predictions with existing data from exclusive measurements.
Although the QPM model Soloviev et al. is not fully consistent since the residual interaction between quasiparticles is not derived from the quasiparticle mean field, it has the advantage that its two-body residual interaction is chosen of a multipole-multipole separable form and therefore, it enables one to work with configuration spaces of large dimensions without facing the cumbersome problem of diagonalizing large matrices. This is very important for the results to be meaningful as it was already shown in a previous work on neutron emission.
The mean potentials entering the single-particle hamiltonians $`h_i`$ are of Woods-Saxon form. Their parameters are chosen according to Refs. with some readjustments such that, once the coupling to phonons is included, the ground states are located at their respective experimental positions with respect to the proton separation threshold. The QPM hamiltonian includes the monopole pairing which must be taken into account for nuclei with open shells. This is the case when we calculate the core properties of <sup>59</sup>Cu and <sup>91</sup>Nb.
The residual particle-hole interaction of Eq. (3) is taken of a separable form in coordinate space with effective interaction strengths considered as adjustable parameters. For the radial interaction form factor we have used $`f(r)=dU/dr`$ where $`U(r)`$ is the central part of the Woods-Saxon potential. For each excitation mode of the even-even core corresponding to a given angular momentum, parity and isospin, the interaction strength is found by requiring that the lowest collective state calculated in RPA be at the experimental energy. These interaction strengths are also used for the quasiparticle-phonon coupling. We have modified the single-particle proton spectrum for <sup>91</sup>Nb in comparison with Ref. to get a better description of high-lying $`1g_{7/2}`$, $`1h_{11/2}`$ subshells, but so far as energies of levels near the Fermi surface are practically the same there are no changes in the properties of low-lying vibrational states in comparison with those of Ref.. For example, the calculated B(E$`\lambda `$) values of the first low-lying collective states $`3_1^{}`$, $`5_1^{}`$ in <sup>40</sup>Ca are B(E3)=1.36$`e^2fm^6`$, B(E5)=2.67$`e^2fm^{10}`$ that can be compared with the corresponding experimental values 1.24$`e^2fm^6`$ and 2.97$`e^2fm^{10}`$. For all nuclei under consideration the calculated values of B(E$`\lambda `$) are in general agreement with experiment. In the present model, these low-lying phonons of the core are the physical channels where the initial state in the excited odd-A nucleus can decay by semi-direct proton emission. In the actual calculation, a very large number of RPA phonons are included in the quasiparticle-phonon basis, but for <sup>40</sup>Ca the first quadrupole state cannot be treated as a one-phonon or particle-hole state and it is outside the present model. On the other hand, the experimental data show that the proton decay of <sup>41</sup>Sc to the $`2_1^+`$ state in <sup>40</sup>Ca is very weak.
At high excitation energies, when the level density becomes large, it is often more convenient to calculate the strength function instead of solving the secular equations (14). One defines the strength function $`\overline{C}_\alpha ^2(E)`$ as the strength distribution $`|C_\alpha ^{(i)}|^2`$ folded with an averaging function $`\rho `$ :
$`\overline{C}_\alpha ^2(E)`$ $`=`$ $`{\displaystyle \underset{i}{}}|C_\alpha ^{(i)}|^2\rho (E\mathrm{\Omega }_i),`$ (23)
where $`\rho (z)`$ is usually chosen as a Breit-Wigner function, $`\rho (z)=\frac{\mathrm{\Delta }}{2\pi }\frac{1}{z^2+\mathrm{\Delta }^2/4}`$. It is shown in Ref. that the calculation of the strength function $`\overline{C}_\alpha ^2(E)`$ can be done without the detailed knowledge of the amplitudes $`|C_\alpha ^{(i)}|^2`$. The same averaging procedure can be applied to calculate the partial escape widths and the branching ratios. In all calculations performed in this work we have adopted the value $`\mathrm{\Delta }`$ = 0.1 MeV.
## III Results and discussion
Using the methods presented above we have calculated the cross sections and branching ratios for the non-statistical proton decay of high-lying states with high angular momenta in <sup>41</sup>Sc, <sup>59</sup>Cu and <sup>91</sup>Nb going to the ground and low-lying collective states of <sup>40</sup>Ca, <sup>58</sup>Ni and <sup>90</sup>Zr, respectively.
An example of cross section calculations by the strength function method is shown in Fig. 1 for the case of proton decay of the $`g_{9/2}`$ states in <sup>41</sup>Sc. Since the calculated cross section is proportional to the square of the one quasiparticle amplitude of the wave function this figure enables one to see the $`1g_{9/2}`$ strength distribution in <sup>41</sup>Sc, too. The $`1g_{9/2}`$ strength is distributed over a broad energy interval due to the coupling with the collective vibrations. Results of our calculations for the $`1g_{9/2}`$ proton particle strength and experimental data obtained from the <sup>40</sup>Ca(<sup>3</sup>He,dp) reaction at $`E_{{}_{}{}^{3}He}`$=240 MeV are presented in Table 1. It is seen from this table that the calculations reproduce reasonably well the integral characteristics of the $`1g_{9/2}`$ proton strength distribution in <sup>41</sup>Sc.
The partial contributions of the non-statistical proton decay into the ground state and $`3_1^{}`$, $`5_1^{}`$ states of <sup>40</sup>Ca are presented in Fig. 2. It is seen from Fig. 2 that the partial cross sections are very energy dependent and as a result different channels can dominate at some excitation energies. The solid curve in Fig. 1 presents the sum of all 3 partial channels. As one can see from Fig. 2 at excitation energies below 8.4 MeV the ground state channel dominates, but there are strong transitions to the $`5_1^{}`$ state at 7.5 and 8.0 MeV. The $`3_1^{}`$ and $`5_1^{}`$ channels begin to contribute importantly starting from 8.4 MeV and they become the main contributors after 9.5 MeV.
The sum of partial widths for different channels in some energy intervals are given in Table 2. As one can see from Table 2 for the energy interval 2.0 - 12.4 MeV the proton decay to the $`5_1^{}`$ state gives a contribution of about 54% and the $`3_1^{}`$ channel contributes about 37% in the total sum. There is a predominance of the $`5_1^{}`$ channel in comparison with the $`3_1^{}`$ one in the energy interval 7.0 - 12.4 MeV. Such a behaviour can be understood from the structure of the decaying states and the angular momentum and energy carried away by the emitted proton. It is seen from Eq. (16) that the partial escape amplitude for fixed values of $`lj`$ is proportional to the contribution of the quasiparticle-plus-phonon configuration made of a quasiparticle with the same $`lj`$ and a phonon of the final state. The cross section depends on the square of these amplitudes summed over all $`lj`$ allowed by the angular momentum coupling rules. In the case of the proton decay of the $`g_{9/2}`$ states at excitation energies 7–12 MeV into the $`3_1^{}`$ channel the outgoing protons carry $`(l,j)`$=(1,3/2) mainly whereas it is $`(l,j)`$=(1,3/2) and $`(l,j)`$=(1,1/2) for the $`5^{}`$ channel. Thus, penetration factor arguments favour the latter channel. It is well known that the strongest coupling between the single-particle states and phonons takes place for the collective phonons. That is why the coupling with the low-lying vibrations is mainly responsible of the damping of the high-lying single-particle modes and giant resonances in nuclei. This is the case for the $`3_1^{}`$ and $`5_1^{}`$ states in <sup>40</sup>Ca. By this reason the strengths of the $`3_1^{}`$ and $`5_1^{}`$ channels are redistributed in a broad energy interval. The decrease of the spectroscopic strength with an increase of the excitation energy leads to a decreasing of cross sections of the channels discussed above. It is worth to mention that a very similar behaviour for the $`3_1^{}`$ and $`5_1^{}`$ channels for the neutron decay of the high angular momentum states in <sup>209</sup>Pb also takes place . Our conclusions about the role of different channels for the proton decay in <sup>41</sup>Sc agree with observations made in Ref., but no separation of the direct part from the statistical one for the proton decay was done in that work.
As a second example we consider the proton decay of the $`g_{9/2}`$ states in <sup>59</sup>Cu. Results of our calculations for the $`g_{9/2}`$ strength distribution in <sup>59</sup>Cu and experimental data are presented in Table 3. The calculations reproduce correctly the strength near 7 MeV and this part is mainly responsible for the proton decay. For the excitation energy interval (2 - 8.9) MeV in <sup>59</sup>Cu the summed ground state width is equal to 10.3 eV and this is much less than in the <sup>41</sup>Sc case. An additional contribution from the 2<sup>+</sup> channel has a summed width that is six times less than that of the ground state channel. There are no transitions to other vibrational states. This is easily understood if one looks at the structure of the states under discussion. For the $`2^+`$ channel the outgoing protons carry $`(l,j)`$=(4,9/2) and there is a suppression of such transitions due to penetration factor effects. The configurations including $`3_1^{}`$ and $`5_1^{}`$ states coupled with continuum are located at somewhat higher energies. This is why no proton decay can proceed to the $`3_1^{}`$ and $`5_1^{}`$ states in <sup>59</sup>Ni.
Very recently, detailed experimental information about proton decay of isobaric analog states (IAS) in <sup>91</sup>Nb has been shown in Ref. . The IAS are strongly excited by means of the <sup>90</sup>Zr$`(\alpha ,t)`$<sup>91</sup>Nb reaction. A sharp peak is seen at 12 MeV excitation energy. Around this excitation energy three IAS have been observed earlier at 11.93 MeV, 12.07 MeV and 12.15 MeV. It is shown in Ref. that at 12 MeV by means of $`(\alpha ,t)`$ reaction predominantly the $`h_{11/2}`$ state is excited. The proton decays to the ground state and some low-lying states are observed. The partial differential cross sections for each final state are given. The data reveal proton decays predominantly to $`5_1^{}`$ and $`3_1^{}`$ excited states of <sup>90</sup>Zr.
Experimental data for energies and spectroscopic factors of IAS and results of our calculations are given in Table 4. Besides, this table also contains calculated partial widths for the proton decay of three IAS to the ground, 2$`{}_{1}{}^{}{}_{}{}^{+}`$, 3$`{}_{1}{}^{}{}_{}{}^{}`$, 4$`{}_{1}{}^{}{}_{}{}^{+}`$ and 5$`{}_{1}{}^{}{}_{}{}^{}`$ states.
As one can see from Table 4 the calculated excitation energies for the $`7/2^+`$ IAS are higher than the experimental ones but the spectroscopic factors extracted from experimental data within 20% accuracy are reproduced reasonably well.
Let us discuss the proton partial widths for these IAS. It is seen from Table 4 that for the first $`7/2^+`$ state the main decay channel is the $`4_1^+`$ channel. For this state the $`(l,j)`$=(0,1/2) protons are most important for the $`4_1^+`$ channel and there is a rather weak transition for the $`2_1^+`$ channel due to the $`(l,j)`$=(2,3/2) protons. The ground state width is proportional to the one quasiparticle strength (the spectroscopic factor). Since the contribution of the one quasiparticle component in the norm of the wave function of this IAS is much less than contributions of the quasiparticle-plus-phonon components the transition to the ground state is weaker than to the $`4_1^+`$. The particularities of the partial proton widths for the second $`7/2^+`$ state can be understood from the structure of this state, too. For the $`4_1^+`$ channel the $`(l,j)`$=(2,3/2),(2,5/2) and (4,7/2) protons can contribute besides the $`(l,j)`$=(0,1/2) protons. As a result the decay width to the $`4_1^+`$ state increases strongly for the second $`7/2^+`$ IAS in comparison with that of the first $`7/2^+`$ state. An increase of the contribution of the configuration constructed from the first quadrupole phonon and $`2d_{3/2}`$ single-particle state in comparison with their contribution to the structure of the first $`7/2^+`$ state results in an essential growth of the decay width to the $`2_1^+`$ channel. The ground state width becomes larger mainly because of the increased spectroscopic factor. The transitions to the $`5_1^{}`$ state can take place due to the outgoing protons with $`(l,j)`$=(5,11/2) but they are supressed because of penetration factors and a small contribution of relevant components in the wave function structure.
In the case of the $`11/2^{}`$ IAS the $`5^{}`$ channel dominates and this is due to the outgoing protons carrying $`(l,j)`$=(0,1/2),(2,3/2) and (2,5/2). The components containing the $`3_1^{}`$ phonon have a small contribution in the wave function structure and as a result the width for the $`3_1^{}`$ channel is small, too. In spite of unfavourable penetration factor the $`2_1^+`$ width is 4 times larger than that of the $`3_1^{}`$ channel. This is due to the configuration consisting of the $`2_1^{}`$ phonon coupled with $`h_{11/2}`$ which contributes about 15% in the norm of the $`11/2^{}`$ IAS.
The comparison of the calculated partial widths with the data of Ref. shows that there is a qualitative agreement. The main channels of the proton decay of (11/2)<sup>-</sup> IAS are reproduced in the calculated structure. The dominance of the $`5_1^{}`$ channel is well established and the calculated partial width of 3.3 keV (Table 4) is in agreement with the measured one (2.9 keV). The $`3_1^{}`$ channel is more pronounced than ground state channel but the calculated partial widths for both channels are much less than the measured ones. The calculations indicate a large decay to the $`2_1^+`$ channel. Such decay is discussed in Ref. but quantitative estimations for the contribution of the $`2_1^+`$ channel in the cross section have not been evaluated in that work.
The summed partial widths for the proton decay of the three IAS to the low-lying vibrational states are also presented in Table 4. As one can see from the last row of Table 4 the $`4_1^+`$ channel exhausts almost 60% of the total cross section, the $`5_1^{}`$ channel contributes about 26% and the $`2_1^+`$ channel gives about 10% of the total strength. The contribution of the ground state channel for the proton decay of the three IAS in <sup>91</sup>Nb is 4% only.
## IV Conclusions
A microscopic approach based on the QPM has been applied to calculate the non-statistical proton decays of high angular momentum states excited in one-nucleon transfer reactions. Partial cross sections and branching ratios for proton emission from high-lying states in <sup>41</sup>Sc, <sup>59</sup>Ni and <sup>91</sup>Nb have been evaluated. The calculated branching ratios demonstrate the existence of strong direct proton decays to the low-lying vibrational states in <sup>41</sup>Sc and <sup>91</sup>Nb and enables one to understand particularities of the decay to different channels. One can conclude from an analysis of calculated partial cross sections that, for high angular momentum states the non-statistical proton decay is more favourable into the higher angular momentum and lower excitation energy final states when the two following conditions are fulfilled: a strong particle-phonon coupling in that channel and a penetration factor which is not hindered by angular momentum or energy. A similar conclusion was also reached for the neutron decay case. A general agreement with existing experimental data is obtained and the predicted branching ratios for different channels can be used to analyze future experimental data.
###### Acknowledgements.
We would like to thank G. Crawley and S. Fortier for fruitful discussions and correspondence. Ch.S. and V.V.V. thank the hospitality of IPN-Orsay where the main part of this work was done. This work is partly supported by IN2P3-JINR agreement and by the Bulgarian Science Foundation (contract No Ph. 801). |
warning/0001/astro-ph0001267.html | ar5iv | text | # Sunyaev-Zel’dovich Surveys: Analytic treatment of cluster detection
## 1 Introduction
Cosmologists have long appreciated the value of the Universe’s biggest objects, galaxy clusters. Besides being a collection of galaxies well suited for studies of galaxy formation, studies focussed on the global properties of clusters provide information on the nature of dark matter; the relative proportions of hot gas, dark matter and stars; and on scenarios of structure formation, including constraints on the universal density parameter, $`\mathrm{\Omega }_\mathrm{o}`$. One example of the latter that comes to mind in anticipation of observations with the new generation of X–ray satellites, Chandra and XMM, is the use of the redshift evolution of the cluster abundance to constrain $`\mathrm{\Omega }_\mathrm{o}`$ (Oukbir & Blanchard 1992, 1997; Bartlett 1997; Henry 1997; Bahcall & Fan 1998; Borgani et al. 1999; Eke et al. 1998; Viana & Liddle 1999a); another is the now classic cluster baryon fraction test (White et al. 1993).
While clusters have been extensively studied in the optical and X–ray bands, observations based on weak gravitation lensing and the Sunyaev–Zel’dovich (SZ) effect (Sunyaev & Zel’dovich 1972) are just coming to fruition. In the case of SZ observations, important samples consisting of several tens of clusters pre–selected in the X–ray are beginning to permit cosmologists to capitalize on the potential of combined SZ/X–ray observations (Carlstrom et al. 1996, 1999). Full maturity of the field will be heralded by the realization of purely SZ–based sky surveys. In what we might refer to as the “SZ–band”, one can then imagine performing cluster science analogous to what is now done in the X–ray, e.g., the construction of cluster counts, redshift distributions, luminosity functions, etc., all viewed via the unique characteristics of the the SZ effect. For example, several authors have emphasized the advantages of the SZ effect, over similar X–ray based efforts, to constrain $`\mathrm{\Omega }_\mathrm{o}`$ via the cluster redshift distribution, as well as to study cluster physics out to very large redshifts (provided the clusters are out there, the very question of $`\mathrm{\Omega }_\mathrm{o}`$ itself) (Korolyov et al. 1986; Bond & Meyers 1991; Bartlett & Silk 1994; Markevitch et al. 1994; Barbosa et al. 1996; Eke et al. 1996; Colafrancesco et al. 1997; Holder et al. 1999). Such pure SZ surveys will be performed: the Planck Surveyor will supply an almost full–sky catalog of several thousand clusters detected uniquely by their SZ signal; and advances in both detector technology and observing techniques now offer the exciting prospect of performing purely SZ–based surveys from the ground, with both large format bolometer arrays and dedicated interferometers.
I discuss in this paper some aspects of the science accessible to pure SZ surveys by examining the nature of their cluster selection. Because of the close analogy with X–ray studies, it is useful for this purpose to compare and contrast SZ–based cluster searches to those based on X–ray observations. The redshift independence of the surface brightness of a cluster (of given properties) means that SZ cluster detection is inherently more efficient than X–ray detection at finding high redshift objects. Equally important is that although the SZ effect and X–rays both “see” the hot intracluster medium (ICM), they do so in significantly different ways. In particular, the well–known fact that the SZ effect scales as the gas pressure implies that the flux density, $`S_\nu `$, is simply proportional to the total thermal energy of the gas. This makes modeling especially simple, for this quantity depends only on the total gas mass and the effectiveness of gas heating during collapse, in stark contrast to the X–ray emission that depends also on the density and temperature distribution of the gas. This simplicity is an advantage because any theoretical interpretation of survey results requires an adequately modeled relation between the observable and the theoretically relevant quantity of cluster mass.
These remarks concern essentially the physics of the ‘emission’ mechanism itself. Of equal relavence is the nature of object selection imposed by the eventual detection algorithm used to extract sources from a set of observations (a map); and this in turn depends crucially, as for any survey, on the particular combination of sensitivity and angular resolution of the observations. The objects detected by Planck will not be the same as those selected by ground–based surveys, and the final catalogs should be viewed as complementary. Planck will produce a shallow ($``$ tens of mJy) large–area survey, while the ground–based instruments will perform deeper surveys ($`<1`$ mJy) over smaller sky areas (several square degrees). Most clusters remain unresolved at the Planck resolution of $`510`$ arcmins, and this characterizes the kinds of objects accessible to this survey, e.g., the counts and the redshift distributions. The higher angular resolution of future ground–based instruments (on the order of an arcmin) will resolve many clusters and impose different selection criteria that will define the counts and redshift distributions of the final catalog.
This is a central issue of the present study were, motivated by the possibility of ground–based surveys, I examine the detection of resolved clusters. While the detection of unresolved sources is principally dependent on observational sensitivity, and the final selection is more or less one of apparent flux – $`S_\nu \theta _\mathrm{c}^2i_\nu `$ – the detection of resolved sources is a more complicated cuisine involving individually the characteristic source size, $`\theta _\mathrm{c}`$, and surface brightness, $`i_\nu `$. The specific goal of the present work is to quantify in terms of observing parameters the abundance, masses and redshifts of clusters detectable by ground–based surveys, with the particular aims of understanding optimal object extraction and the accessible science. For example, one of the key questions facing any survey is one of observing strategy: given a fixed, total amount of observing time, should one “go deep”, with long integrations on a few fields, or instead “go wide”, covering more fields to higher sensitivity. If one is out to maximize the number of detected objects, the answer depends on the slope of the counts. One gains by going deeper if the integrated counts are steeper than $`S_\nu ^2`$, assuming that noise diminishes as $`1/\sqrt{t}`$; otherwise, a larger area yields more objects.
The cluster selection criteria of a survey may be compactly summarized by a minimum detectable mass as a function of redshift – $`M_{\mathrm{det}}(z)`$. Together with a suitable mass function (we shall use the formalism of Press and Schechter 1974), this quantity determines both the source counts and redshift distributions of the final source catalog. Thus, in very concrete terms, we must examine $`M_{\mathrm{det}}(z)`$ and understand the influence of the observationally imposed restrictions on $`\theta _\mathrm{c}`$ and $`i_\nu `$. Given a set of observations, i.e., a map, one could imagine many different algorithms to extract astrophysical sources, and $`M_{\mathrm{det}}(z)`$ will depend upon this choice. There is in principle an optimal method, one which preserves signal–to–noise over the entire range of source surface brightness and size. It is characterized by a decreasing surface brightness limit with object size – the greater number of object pixels permits lower surface brightness detections. This algorithm is difficult to apply in practice, and more standard approaches search instead for a minimum number of connected pixels above a preset threshold, thereby establishing a fixed cut on source surface brightness. Detection signal–to–noise is no longer constant, rather increasing with $`\theta _\mathrm{c}`$, and these methods loose large, and in–principle detectable, low surface brightness objects. All of this will be reflected in the resulting functions $`M_{\mathrm{det}}(z)`$.
Throughout the discussion, we will be guided by the characteristics of two potential types of ground–based instruments: large format bolometer arrays, epitomized by BOLOCAM (Glenn et al. 1998<sup>1</sup><sup>1</sup>1http://phobos.caltech.edu/ lgg/bolocam/bolocam.html ), and interferometer arrays optimized for SZ observations, as suggested by Carlstrom et al. (1999)<sup>2</sup><sup>2</sup>2While writing, I became aware of another project – the Arcminute MicroKelvin Imager. See Kneissl R. 2000, astro-ph/0001106. BOLOCAM is a 151–element bolometer array under construction at Caltech for operation in three bands – 2.1 mm, 1.38 mm (the null of the thermal SZ effect) and 850 $`\mu `$m. At the Caltech Submillimeter Observatory, it is expected that the array will be diffraction limited to $``$ arcminute resolution, or better, and limited in sensitivity by atmospheric emission (rather than detector noise). With its $`9`$–arcmin field–of–view, one could imagine surveying a square degree to sub–mJy sensitivity in these bands. Carlstrom et al. (1999) have recently expounded the virtues of interferometric techniques using telescope arrays specifically designed for SZ observations. They have proposed the construction of such an array, operating at a wavelength of 1 cm, and estimated that it would be capable, in the course of one year of dedicated observations, of covering $`10`$ square degrees to a limiting sensitivity of $`0.3`$ mJy at arcminute resolution. In summary, then, we are interested in considering SZ observations at arcminute angular resolution and to sub–mJy sensitivity at both centimeter and millimeter wavelengths.
The paper is organized as follows: a rapid review of the SZ effect is given in the next section, followed by a discussion of the unique aspects of SZ cluster detection. Section 3 details the cluster population model employed, based on the Press–Schechter (Press & Schechter, 1974) mass function and the isothermal $`\beta `$–model. Since we shall focus on issues of cluster selection as imposed by survey parameters, the cluster model will be restricted to the simple example of a self–similar population. The next section (Section 4) introduces the principal figures (Figures 1,2 and 3) of the present work by consideration of unresolved cluster detection; this case will also be used as a benchmark against which to examine the effects of resolved detection. Section 5 then develops the principle themes of resolved SZ cluster detection, starting with consideration of the optimal, constant signal–to–noise method, and followed by detailed study of cluster detection based on the standard algorithm. A final discussion (Section 6) then more closely examines the number of detections to be expected from ground–based surveys and gives a non–exhaustive list of some important issues still to be treated. Section 7 concludes.
Key results will be the $`M_{\mathrm{det}}(z)`$ curves presented in Figure 1, quantifying the nature of SZ detected clusters, and the conclusion that resolved source counts are lower and steeper than expectations based on simple unresolved source count calculations, Figure 2. To the point, the latter implies that surveys at arcminute resolution gain objects with an observing strategy of “going deep”. The cosmological density parameter is denoted by $`\mathrm{\Omega }_\mathrm{o}8\pi G\rho /3H_\mathrm{o}^2`$, the vacuum density parameter by $`\lambda _\mathrm{o}\mathrm{\Lambda }/3H_\mathrm{o}^2`$ and the Hubble constant by $`H_\mathrm{o}h100`$ km/s/Mpc; unless otherwise indicated, $`h=1/2`$ and $`\lambda _\mathrm{o}=0`$.
## 2 The Particular Value of the SZ Effect
We begin by establishing our notation in recalling the basic formulas of the SZ effect. The change in surface brightness relative to the unperturbed cosmic microwave background (CMB), caused by inverse Compton scattering in the hot ICM, is expressed as
$$i_\nu (\theta )=y(\theta )j_\nu (x)$$
(1)
where $`xh_p\nu /kT_\mathrm{o}`$ is a dimensionless frequency expressed in terms of the energy of the unperturbed CMB Planck spectrum at $`T_\mathrm{o}=2.725`$ K (Mather et al. cmb:temp2 (1999)). The spectral shape is embodied in the function $`j_\nu `$,
$`j_\nu (x)`$ $`=`$ $`2{\displaystyle \frac{(kT_\mathrm{o})^3}{(h_pc)^2}}{\displaystyle \frac{x^4\text{e}^x}{(\text{e}^x1)^2}}\left[{\displaystyle \frac{x}{\mathrm{tanh}(x/2)}}4\right]`$
$``$ $`2{\displaystyle \frac{(kT_\mathrm{o})^3}{(h_pc)^2}}f_\nu `$
$`=`$ $`(2.28\times 10^4\mathrm{mJy}/\mathrm{arcmin}^2)f_\nu `$
while the amplitude is given by the Compton $`y`$–parameter
$$y\text{d}l\frac{kT}{m_ec^2}n_e\sigma _T$$
(3)
an integral of the pressure along the line–of–sight at position $`\theta `$ relative to the cluster center. Here, $`T`$ is the temperature of the ICM (really, the electrons), $`m_e`$ is the electron rest mass, $`n_e`$ the ICM electron density, and $`\sigma _T`$ is the Thompson cross section. Planck’s constant is written in these expressions as $`h_p`$, the speed of light in vacuum as $`c`$, and Boltzmann’s constant as $`k`$. These formulae apply in the non–relativistic limit of low electron (and photon) energies; relativistic extensions have recently been made by several authors (e.g., Rephaeli 1995; Stebbins 1997; Challinor & Lasenby 1998; Itoh et al 1998; Pointecouteau et al. 1998; Sazonov & Sunyaev 1998). The spectral shape of the distortion is unique, becoming negative at wavelengths larger than $`1.4`$ mm (relative to “blank” sky) and positive a shorter wavelengths. This offers a way of clearly separating the effect from other astrophysical emissions.
All of the physics is in the Compton $`y`$–parameter, an apparently innocuous–looking expression. In fact, it holds the key to all of the pleasing aspects of the SZ mechanism. First of all, the conspicuous absence of an explicit redshift dependence is the well–known result that the SZ surface brightness is redshift–independent, determined only by cluster properties. This should be contrasted to other emission mechanisms which all experience “cosmic dimming” \[$`\iota (1+z)^4`$\] due to the expansion of the Universe. This is countered in the SZ effect by the increasing energy density towards higher $`z`$ of the CMB, the source of photons for the effect.
Another very important aspect of the SZ mechanism resides in the fact that its amplitude is proportional to the pressure, or thermal energy, of the ICM. This appears most clearly when we consider the total flux density from a cluster, found by integrating the surface brightness over the cluster face:
$`S_\nu (x,M,z)`$ $`=`$ $`j_\nu (x)D_a^2(z){\displaystyle \text{d}V\frac{kT(M,z)}{m_ec^2}n_e(M,z)\sigma _T}`$ (4)
$``$ $`M_{\mathrm{gas}}<T>`$
The integral is over the entire virial volume of the cluster. In this expression, $`D_a(z)`$ is the angular–size distance in a Friedmann–Robertson–Walker metric –
$`D_{\mathrm{ang}}(z)=2cH_\mathrm{o}^1\left[{\displaystyle \frac{\mathrm{\Omega }_\mathrm{o}z+(\mathrm{\Omega }_\mathrm{o}2)(\sqrt{1+\mathrm{\Omega }_\mathrm{o}z}1)}{\mathrm{\Omega }_\mathrm{o}^2(1+z)^2}}\right]`$
$`=cH_\mathrm{o}^1D(z)`$ (5)
where I introduce the dimensionless quantity $`D(z)`$. We see clearly that the final result is simply proportional to the total thermal energy of the ICM, $`𝑑VnT`$. This is extremely important, because it means that the SZ flux density is insensitive (strictly speaking, completely so for the total flux density and for fixed thermal energy) to either the spatial distribution of the ICM or its temperature structure, making modeling much simpler than in the case of X–ray emission. Consider that in X–ray modeling one prefers the X–ray temperature over luminosity as a more robust indicator of cluster mass, but even the temperature has some sensitivity to the gas distribution, because it is all the same an emission weighted temperature that is actually observed. We would expect the temperature appearing in the second line of Eq. (4), which is the true mean electron energy, to demonstrate an even better correlation with virial mass than the observed X–ray temperature. Simple scaling arguments lead one to believe that this correlation should be $`TT_{\mathrm{virial}}M^{2/3}(1+z)`$, from which we deduce
$$S_\nu f_{\mathrm{gas}}(M,z)M^{5/3}(1+z)D^2(z)$$
(6)
where $`f_{\mathrm{gas}}`$ is the gas mass fraction contributed by the ICM to the total cluster mass.
The SZ mechanism therefore conveniently reduces all the potential complexity of the ICM to just its total thermal energy, $`f_{\mathrm{gas}}<T>`$. This quantity may nevertheless be influenced by several factors. For example, the gas mass fraction in Eq. (6) has carefully been written as a general function of both mass and redshift. In simulations this quantity is most often constant, the majority of gas being primordial and simply falling into the cluster at formation. One could imagine other possibilities (e.g., Bartlett & Silk 1994; Colafrancesco & Vittorio 1994) that would lead to a more important dependence on either mass or redshift, although metallicity arguments seem to require that most of the gas be primordial, at least in the more massive systems (Metzler & Evrard 1994; Elbaz et al. 1995). While it appears from numerical studies that shocking during cluster formation efficiently heats the ICM to $`80`$%– $`100`$% of the virial temperature (Metzler & Evrard 1994; Bryan & Norman 1998), additional sources of heating could in principle change the temperature of the gas relative to that of the potential, i.e., $`TT_{\mathrm{virial}}`$. Such heating may not always produce the most obvious effects – remember that it is the total thermal energy of the gas that counts, and understanding the change of this quantity with heating in a gravitational potential requires careful modeling. Although models studied so far do not lead to a strong effect (Metzler & Evrard 1994), we shall at times be discussing rather low mass systems, for which these effects are poorly understood theoretically and observationally. Finally, the exact form of the virial temperature–mass relation depends in part on the dark matter profile of the collapsing proto–cluster; once again, numerical experiments seem to indicate that this does not change too much, i.e., one finds a good T–M relation with rather small scatter (Evrard et al. 1996; Bryan & Norman 1998). Putting all of this together, a relation of the form (6) between the observable, $`S_\nu `$, and cluster mass appears quite reasonable and rather robust; and in any case, the modeling uncertainties are always easier to understand than in the case of X–rays, due to the all important insensitivity of the SZ flux density to spatial/temperature structure of the ICM.
The conclusion is that the SZ flux density should be a very good halo mass detector, in principle sensitive to all halos with significant amounts of hot gas and over a large range of redshifts. All of these remarks concern to a large extent the total SZ flux density of Eq. (4), and therefore apply primarily to situations where the clusters are unresolved. It is still true that, even when a cluster is resolved, the SZ signal is proportional to the total thermal energy of the gas, but now only of that portion contained within the column defined by the beam. After first outlining the cluster population model employed, we shall tackle in detail the additional complexities introduced by resolved cluster observations.
## 3 Modeling the Source Population
The central ingredient of a model for the cluster population and its evolution is the mass function, $`n(M,z)`$, which gives the number density of collapsed, virialized objects as a function of mass and redshift. The exact form of this function depends on the statistical properties of the primordial density fluctuations. For Inflationary–type scenarios, in which these fluctuations are Gaussian, a reasonable expression for the mass function appears to be the Press–Schechter formula (Press & Schechter 1974)
$$n(M,z)dM=\sqrt{\frac{2}{\pi }}\frac{<\rho >}{M}\nu (M,z)\left|\frac{\text{d}\mathrm{ln}\sigma (M)}{\text{d}\mathrm{ln}M}\right|\text{e}^{\nu ^2/2}\frac{dM}{M}$$
(7)
The quantity $`\rho `$ represents the comoving cosmic mass density and $`\nu (M,z)\delta _c(z)/\sigma (M,z)`$, with $`\delta _c`$ equal to the critical linear over–density required for collapse and $`\sigma (M,z)`$ the amplitude of the density perturbations on a mass scale $`M`$ at redshift $`z`$. Numerical studies ascribe rather remarkable accuracy to the simple expression of Eq. (7) (Lacey & Cole 1994; Eke et al. 1996; Borgani et al. 1999), and we shall adopt it in the following. More explicitly, $`\delta _c(z,\mathrm{\Omega }_\mathrm{o},\lambda _\mathrm{o})`$ and $`\sigma (M,z)=\sigma _\mathrm{o}(M)\times (D_\mathrm{g}(z)/D_\mathrm{g}(0))`$, with $`D_\mathrm{g}(z,\mathrm{\Omega }_\mathrm{o},\lambda _\mathrm{o})`$ being the linear growth factor. It is essentially through $`D_\mathrm{g}`$ that the dependence on cosmology ($`\mathrm{\Omega }_\mathrm{o}`$, $`\lambda _\mathrm{o}`$) enters the mass function, with $`\mathrm{\Omega }_\mathrm{o}`$ being the more important of the two as the dependence on $`\lambda _\mathrm{o}`$ is relatively weak (see, e.g., Bartlett casa (1997) for a detailed discussion). This dependence on $`\mathrm{\Omega }_\mathrm{o}`$ in the exponent means that the cluster abundance as a function of redshift is a very sensitive probe of the density parameter (e.g., Oukbir & Blanchard 1992, 1997), and is the motivation for many efforts in all wavebands to find clusters at high redshifts. As emphasized by several authors (Barbosa et al. 1996; Eke et al. 1996; Colafrancesco et al. 1997; Bartlett et al. 1998; Holder & Carlstrom 1999; Mohr et al. 1999), the SZ effect is particularly well positioned in this arena (see also below).
It is clear that the important theoretical variables are cluster mass and redshift. Although redshift is directly measurable, the mass appearing in Eq. (7) must be translated into an observational quantity suitable for the type of observations under consideration. As mentioned above, one of the pleasant features of the SZ effect is the simplicity of this relation. Using the simulations of Evrard et al. (1996) to normalize the $`TM`$ relation, we can quantitatively express the total SZ flux density of a cluster (e.g., Eqs. 4 & 6) as
$`S_\nu =(34\text{mJy}h^{8/3})f_\nu (x)f_{\mathrm{gas}}\mathrm{\Omega }_o^{1/3}\left[{\displaystyle \frac{\mathrm{\Delta }_{\mathrm{NL}}(z)}{178}}\right]^{1/3}M_{15}^{5/3}`$
$`(1+z)D^2(z)`$ (8)
where the mass $`M_{15}M/10^{15}\mathrm{M}_{}`$ refers to the cluster virial mass and $`f_{\mathrm{gas}}`$ is possibly a function of both mass and redshift (see also Barbosa et al. 1996, but note that the definition there of $`D(z)`$ differs by a factor of $`2`$). Evrard (1997) finds $`f_{\mathrm{gas}}=0.06h^{1.5}`$, while Mohr et al. (1998) find marginal evidence for a decrease in lower mass systems (see also Carlstrom et al. 1999 for recent work based on SZ images); there is little information on any possible evolution with redshift at present. Other quantities appearing in this equation are the mean density contrast for virialization, $`\mathrm{\Delta }_{\mathrm{NL}}(z,\mathrm{\Omega }_\mathrm{o},\lambda _\mathrm{o})`$ ($`=178`$ for $`\mathrm{\Omega }_\mathrm{o}=1`$, $`\lambda _\mathrm{o}=0`$), and the dimensionless functions $`f_\nu `$ and $`D(z)`$ introduced in Eqs. (2) and (2).
Observations for which clusters are unresolved measure this total flux density, and therefore this is all that is needed in order to calculate the unresolved source counts, as we will do in the next section. For resolved sources, on the other hand, the detection criteria are more complicated. Contrary to the point source limit, the details of the cluster SZ profile now play an important role. I will employ a simple isothermal $`\beta `$–model to describe this profile:
$$i_\nu (\theta )=\frac{y_\mathrm{o}j_\nu (x)}{(1+\theta ^2/\theta _\mathrm{c}^2)^\alpha }$$
(9)
The exponent $`\alpha =0.5(3\beta 1)`$, where $`\beta `$ is the exponent of the three–dimensional ICM density profile: $`n(1+r^2/r_\mathrm{c}^2)^{3\beta /2}`$, $`r_\mathrm{c}`$ being the physical core radius. Local X–ray observations indicate that $`\beta 2/3`$, a value I adopt throughout for the calculations. In this case, $`\alpha =1/2`$, a rather significant value, as will be discussed shortly. This profile will be assumed to hold out to the virial radius, $`R_\mathrm{v}`$, of the cluster.
The $`\beta `$–profile of Eq.(9) is empirically described by $`y_\mathrm{o}`$, a sort of central surface brightness (actually, it is $`y_\mathrm{o}j_\nu `$ that has units of surface brightness, but it is simpler to work with $`y_\mathrm{o}`$), and $`\theta _\mathrm{c}`$. In these terms, there is nothing specific to the SZ effect. The physics of the SZ effect appears only when we make the connection between these empirical parameters and the theoretically interesting ones, namely, mass and redshift, via relations of the kind $`y_\mathrm{o}(M,z)`$ and $`\theta _\mathrm{c}(M,z)`$. As our principle goal in this work is to understand the selection effects of resolved SZ cluster detection, the model for cluster evolution will be kept simple: a constant gas mass fraction, $`f_{\mathrm{gas}}=0.06h^{1.5}`$ (Evrard 1997), over cluster mass and redshift, and a core radius scaling with the virial radius $`R_\mathrm{v}`$, i.e., $`x_\mathrm{v}R_\mathrm{v}/r_\mathrm{c}=const`$. Unless otherwise specified, this constant will be given a value of 10. One deduces from simple scaling arguments that
$$R_\mathrm{v}=(1.69h^{2/3}\mathrm{Mpc})M_{15}^{1/3}(1+z)^1\mathrm{\Omega }_o^{1/3}\left(\frac{178}{\mathrm{\Delta }_{\mathrm{NL}}(z)}\right)^{1/3}$$
where the normalization is taken from the spherical collapse model. This scaling relation is about as robust as the relation for cluster temperature; in fact, the two are essentially the same, since $`TM/R_\mathrm{v}`$. Some dependence of the normalization on mass and redshift could appear if the density profile around a peak forming a cluster changed significantly with these two quantities. In the following, we shall ignore this possibility, which numerical simulations seem to indicate is a small effect in any case. This then fixes the relation
$$r_\mathrm{c}(M,z)=R_\mathrm{v}(M,z)/x_\mathrm{v}$$
(10)
For the axially symmetric surface brightness of Eq. (9), the integral defining the total SZ flux density may be written
$`S_\nu (M,z)`$ $`=`$ $`j_\nu 2\pi {\displaystyle 𝑑\theta \theta y(\theta )}`$
$`=`$ $`2\pi j_\nu y_\mathrm{o}(M,z)\theta _\mathrm{c}^2(M,z)\left(\sqrt{1+x_\mathrm{v}^2}1\right)`$
Using Eq. (3) for $`S_\nu (M,z)`$ in this expression, we deduce
$`y_\mathrm{o}(M,z)`$ $`=`$ $`(6.40\times 10^5h^2)f_{\mathrm{gas}}\mathrm{\Omega }_\mathrm{o}\left({\displaystyle \frac{\mathrm{\Delta }_{\mathrm{NL}}(z)}{178}}\right)`$ (11)
$`M_{15}(1+z)^3\left({\displaystyle \frac{x_\mathrm{v}^2}{\sqrt{1+x_\mathrm{v}^2}1}}\right)`$
Together with the $`\beta `$–profile (Eq. 9), Eqs. (10) and (11) define our cluster evolution model. As mentioned, it is self–similar, and we see the expected scaling $`r_\mathrm{c}M^{1/3}/(1+z)^1`$ and $`y_\mathrm{o}M(1+z)^3`$. This is most probably an oversimplified description of the actual cluster population, but it nevertheless provides a ‘standard’ with which we may understand the nature of the selection effects imposed by resolved cluster detection, and a benchmark for comparing more detailed models. It is important in the following that one does not forget the model dependence of our results, which can be retraced to this point of the discussion.
## 4 Unresolved Detections
This section is dedicated to the simple case of unresolved SZ detection, which will be used as a reference in the following discussion of resolved detection. It also offers an introduction to the main figures, Figures 1, 2 and 3, summarizing the essential results of the present work. They are constructed for two representative cosmologies: a critical model $`\mathrm{\Omega }_\mathrm{o}=1`$, and an open model ($`\lambda _\mathrm{o}=0`$) with $`\mathrm{\Omega }_\mathrm{o}=0.3`$. For the counts and redshift distributions of Figures 2 and 3, I have used a CDM–like power spectrum with “shape parameter” fixed at $`\mathrm{\Gamma }=0.25`$; both models are normalized to the present day abundance of X–ray clusters – $`\sigma _8=0.6`$ and $`\sigma _8=1.0`$ for the critical and open models, respectively (e.g., Blanchard et al. 1999; Borgani et al. 1999; Viana & Liddle 1999b).
Observations for which most clusters are unresolved measure the total SZ flux density. One can then simply invert Eq. (3) to find the corresponding limiting detection mass as a function of redshift, $`M_{\mathrm{det}}^{\mathrm{ur}}(z,S_\nu )`$:
$`M_{\mathrm{det}}^{\mathrm{ur}}(z,S_\nu )`$ $`=`$ $`(0.12\times 10^{15}h^{8/5}\mathrm{M}_{})\left({\displaystyle \frac{\mathrm{S}_\nu }{\mathrm{mJy}}}\right)^{3/5}`$ (12)
$`(f_\nu f_{\mathrm{gas}})^{3/5}\mathrm{\Omega }_\mathrm{o}^{1/5}\left({\displaystyle \frac{178}{\mathrm{\Delta }_{\mathrm{NL}}}}\right)^{1/5}`$
$`D^{6/5}(z)(1+z)^{3/5}`$
Integrating the mass function over redshift and over masses greater than this limit directly yields the source counts:
$$\frac{dN}{d\mathrm{\Omega }}(>S_\nu )=_0^{\mathrm{}}𝑑z\frac{dV}{dzd\mathrm{\Omega }}_{M_{\mathrm{det}}(z,S_\nu )}^{\mathrm{}}𝑑M\frac{dn}{dM}(M,z)$$
(13)
The corresponding redshift distribution is simply obtained as the integrand of the $`z`$–integral.
Figure 1 compares the various detection masses as a function of redshift for observations at $`2`$ mm, e.g., a bolometer array, and at $`1`$ cm, representative of an interferometer; in each case the upper (red) curve corresponds to the open model. For the moment, concentrate only on the the dashed lines, which give the result for unresolved detection, Eq. (12), at a flux density of $`S_\nu =1.5`$ mJy. These curves remain unchanged from Figure 1b to 1c, both at $`1`$ cm but differing in angular resolution, because resolution is irrelevant for point sources (ignoring source confusion issues). Observe that in all cases the detection mass decreases with redshift beyond $`z1`$. This remarkable behavior is directly attributable to the fact that the SZ surface brightness is independent of distance. As already emphasized, the distance appearing in Eq. (12) is the angular distance and not the luminosity distance, a factor of $`(1+z)^2`$ larger. At high $`z`$ the redshift dependence therefore scales as $`z^{9/5}`$, one power coming from the assumed redshift scaling of the virial temperature and the rest from the decrease in angular distance (focusing) as $`1/z`$. A self–similar cluster model, implicitly assumed in this context by the constancy of $`f_{\mathrm{gas}}`$, thus predicts that SZ observations are more sensitive to objects at large, rather than intermediate, redshifts. This overall behavior would not change even if we broke the self–similarity with a declining gas mass fraction with mass; such a dependence could only modify the rate of decrease with $`z`$. On the other hand, an explicit decrease in $`f_{\mathrm{gas}}`$ with redshift stronger than $`(1+z)^3`$ would cause $`M_{\mathrm{det}}^{\mathrm{ur}}`$ to actually increase with redshift. It is perhaps not so surprising that at close range, small $`z`$, the detection mass also drops; this is simply due to the increasing angular size of the object creating an increase in total flux density (the source is assumed to always remain unresolved in this discussion).
From the difference between Figures 1a and 1b,c, we see that, at a given sensitivity, the $`2`$ mm observations probe farther down in mass. This is nothing more than the spectral shape of the SZ effect, described by the function $`j_\nu `$: the biggest decrement occurs precisely near $`2`$ mm (the maximum emission of the effect is around $`750\mu `$m). The resolved detection mass limits, to be shortly discussed, depend also on the angular resolution.
Source counts for the two cosmological scenarios are given in Figure 2. These have been calculated using Eq. (13) and the appropriate detection mass. In order to shed some light on the importance of low mass objects to these results, the counts are presented in pairs, one curve for a low mass cut–off of $`10^{13}\mathrm{M}_{}`$ and one for a cut–off of $`10^{14}\mathrm{M}_{}`$. Note that the $`x`$–axis denotes the pixel noise, $`\sigma _{\mathrm{pix}}`$, and not a limiting source flux density; in the present situation of unresolved detections, this just means that the corresponding limiting flux density is $`q_{\mathrm{det}}\times \sigma _{\mathrm{pix}}`$.
The first thing to remark from Figure 2 is the large difference between the two cosmological models. The presence of clusters at high redshift in a low–density model shows up in the integrated counts, as confirmed by the corresponding redshift distributions shown in Figure 3, where the huge difference in cluster abundance at large redshift is evident. It is for this reason that the redshift distribution of SZ sources is a potentially powerful tool for constraining $`\mathrm{\Omega }_\mathrm{o}`$ (Barbosa et al. 1996, Bartlett et al. 1998). This is of foremost importance and represents one of the primary motivating factors behind this type of survey.
This situation of unresolved sources applies in practice to missions such as the Planck Surveyor, as discussed, for example, by Barbosa et al. (1996) and Aghanim et al. (1997). The higher angular resolution of possible ground–based surveys calls for examination of resolved source detection.
## 5 Resolved Detections
In this, the principle section of this paper, we treat in detail the issue of resolved SZ cluster detection. The context will be one of arcminute resolution (pixel size) and sub–mJy sensitivity, as targeted by the up–coming ground–based instruments. It is worth being very explicit about the nature of the observations: the simplest case to imagine corresponds to that of an image produced by a bolometer array, such as BOLOCAM. In this case each point on the image, a ‘pixel’, represents a sample point of the sky brightness, as transformed by the optics of the observing system. The optical response may be divided into that of the telescope–plus–atmosphere (defining the projection of the sky onto the focal plane) and the optics proper to the detector (which act on the focal–plane image). There is a difference between bolometer arrays and the familiar example of a CCD camera working in the visible. For the latter, atmospheric seeing and telescope optics project the sky onto the focal plane by convolving with a Gaussian, and the camera itself then convolves this focal–plane image with a square top–hat, one centered on each pixel. The difference with a bolometer array lies in the fact that the CCD camera defines sharp, well–defined pixel boundaries, while a bolometer array, with its set of cones, convolves the focal–plane image with something closer to a Gaussian. This means that, unlike CCDs, the pixels of a bolometer array ‘overlap’ in the focal plane. This has little consequence for the ensuing discussion, but it is all the same worth keeping in mind.
This picture is not completely accurate when it comes to interferometers. Such instruments actually directly sample the Fourier transform of the sky. The result may often be modeled by a real sky image convolved with an effective, synthesized beam, but this beam lacks sensitivity on large scales, i.e., large spatial wavelengths on the sky (short baselines). Thus, the effective beam cannot not be precisely a Gaussian, and it is especially important to correctly model the loss of response on large scales for extended objects such as clusters. For the ensuing discussion, I adopt the bolometer picture, applying it at times rather indiscriminately to characterize ground–based observations; a future work will consider the details specific to interferometric observations (see also the recent work of Holder et al. 1999).
For a bolometer array, the response of the entire optical chain (atmosphere-telescope-detector) is often adequately modeled as a bi–dimensional Gaussian (if one is lucky, a symmetric one!), and for proper sampling, respecting Shannon, the sample period must be 2 – 3 times smaller than the beam FWHM. We will characterize a survey by the pixel size and sensitivity per pixel of its images – $`\mathrm{\Omega }_{\mathrm{pix}}`$, a solid angle, and $`\sigma _{\mathrm{pix}}`$, a flux density. Note that because the pixels ‘overlap’ in the focal plane, what precisely is meant by $`\mathrm{\Omega }_{\mathrm{pix}}`$ is the square of the separation between sample points, $`\theta _{\mathrm{pix}}`$; the concept is a bit more ambiguous than in the case of a CCD camera. Thus, proper sampling means that the pixel scale $`\mathrm{\Omega }_{\mathrm{pix}}\theta _{\mathrm{pix}}^2\theta _{\mathrm{fwhm}}^2/4`$. It is also worth explicitly remarking that, in the following, I assume that the noise is uncorrelated (from pixel to pixel) and uniform over the image.
Given, then, a map of the SZ sky, we would like to understand how to extract clusters and the nature of the selection imposed by our technique. In addition to the observational parameters $`\mathrm{\Omega }_{\mathrm{pix}}`$ and $`\sigma _{\mathrm{pix}}`$, this will depend on the form of the extended emission of the sources, a complication avoided in the case of unresolved cluster detection; this represents an important difference between the two situations. Employing the $`\beta `$–model introduced previously, Eq. (9), we see that a cluster SZ profile may be described by a characteristic central surface brightness, $`y_\mathrm{o}`$, and an angular size, $`\theta _\mathrm{c}`$ (the core radius). When couched in terms of the purely empirical parameters of $`\mathrm{\Omega }_{\mathrm{pix}}`$, $`\sigma _{\mathrm{pix}}`$, $`y_\mathrm{o}`$ and $`\theta _\mathrm{c}`$, we have before us a rather classic and well–known problem of Astronomy. The only difference with galaxies in the optical is the form of the source profile. All physics specific to the SZ effect itself appears only in the relation of the empirical source descriptors – $`(y_\mathrm{o},\theta _\mathrm{c})`$ – to the theoretically meaningful ones of cluster mass, $`M`$, and redshift, $`z`$.
The procedure in the following is then always the same: quantify the detection algorithm in terms of $`\mathrm{\Omega }_{\mathrm{pix}}`$ and $`\sigma _{\mathrm{pix}}`$, and then translate this, via the isothermal $`\beta `$–model, into a $`M_{\mathrm{det}}(z;\mathrm{\Omega }_{\mathrm{pix}},\sigma _{\mathrm{pix}})`$. I employ a notation where the imminently interesting independent variables of a function appear before the “;”, and parameterizing ones afterward. Thus, as written, the detection mass is primarily a function of redshift, parameterized by the survey properties $`\mathrm{\Omega }_{\mathrm{pix}}`$ and $`\sigma _{\mathrm{pix}}`$. This function teaches us about the kinds of objects we detect, and leads directly to the survey counts and the redshift distribution of our clusters, via Eq. (13). These latter quantities are the key indicators of the science content of the survey.
This procedure will be applied to two source extraction methods in the following, and the results compared to those for an unresolved SZ survey. We will refer to the first as “optimal detection”, because it extracts sources in such a way as to preserve the signal–to–noise across the entire range of detectable surface brightness and source size. This is achieved by lowering the surface brightness limit for large sources, possible due to the greater number of covered pixels. The second method, routinely used by such packages as SExtractor (Bertin & Arnouts 1996), searches for a minimum number of connected pixels above a preset threshold. The important difference with the first technique is the imposition of a fixed surface brightness limit, independent of source size. The signal–to–noise of the detections is no longer constant, but increases with source size. This technique may be considered sub–optimal in the sense that it loses in–principle detectable low surface brightness sources, a fact well appreciated in the case of optical galaxy surveys.
### 5.1 Optimal case
Optimal detection selects all sources with a flux density
$$S_\nu q_{\mathrm{opt}}N^{1/2}\sigma _{\mathrm{pix}}$$
(assuming spatially uncorrelated and uniform noise) where $`N`$ is the number of pixels covered by the cluster and $`q_{\mathrm{opt}}`$ represents a threshold, say $`q_{\mathrm{opt}}35`$; in fact, $`q_{\mathrm{opt}}=S/N`$, the signal–to–noise of the detection. Notice also that, as advertised, the limiting surface brightness decreases with object size: $`<i_\nu >S_\nu /Nq_{\mathrm{opt}}\sigma _{\mathrm{pix}}/\sqrt{N}`$. One extracts in this way all objects detectable at a given $`S/N`$, and for this reason we may refer to the method as optimal. The number of object pixels is simply found as $`N=\pi \theta _{\mathrm{vir}}^2/\mathrm{\Omega }_{\mathrm{pix}}`$, where $`\theta _{\mathrm{vir}}=R_\mathrm{v}/D_{\mathrm{ang}}`$ is angular virial radius. This permits us to express the detection mass as
$`M_{\mathrm{det}}^{\mathrm{opt}}(z,S_\nu )`$ $`=`$ $`(0.19\times 10^{15}h^2\mathrm{M}_{})\left({\displaystyle \frac{\mathrm{q}_{\mathrm{opt}}\sigma _{\mathrm{pix}}}{\mathrm{mJy}}}\right)^{3/4}`$ (14)
$`\left({\displaystyle \frac{\mathrm{arcmin}^2}{\mathrm{\Omega }_{\mathrm{pix}}}}\right)^{3/8}(f_\nu f_{\mathrm{gas}})^{3/4}\mathrm{\Omega }_\mathrm{o}^{1/2}`$
$`\left({\displaystyle \frac{178}{\mathrm{\Delta }(z)}}\right)^{1/2}D^{3/4}(z)(1+z)^{3/2}`$
As written, this criteria uses an aperture corresponding to the full angular size of the object – $`S_\nu `$ is understood to be the total SZ flux density in Eq. (3). For resolved sources, one would like to chose an aperture which optimizes the signal–to–noise ratio of the detection. Interestingly, a 3D gas profile close to $`r^2`$, corresponding to a SZ surface brightness $`y\theta ^1`$, results in a constant signal–to–noise with aperture radius. A $`\beta `$–model with $`n(1+r^2/r_\mathrm{c}^2)^{3\beta /2}`$ and $`\beta 2/3`$ exhibits this behavior at large radii, for example: $`y(\theta )(1+\theta ^2/\theta _c^2)^{1/2}`$. In this case, the signal–to–noise of a SZ detection increases from the center of the cluster image out to the core radius, $`r_\mathrm{c}`$, beyond which it turns over to a constant out to the virial radius. The situation is different for X–ray images, where the surface brightness falls off more rapidly, diving under the background at large radii. From this we conclude that the simple criteria given above provides in fact an optimum SZ detection (at least as long as $`\beta `$ remains close to 2/3, as appears to be the case locally).
The detection mass Eq. (14) is displayed in Figure 1 as the solid lines. Compared to the hypothetical point source results, observations resolving clusters are less efficient at detecting clusters, especially at intermediate redshifts. This is easy to understand as the effect of distributing a given flux density over $`N`$ pixels, each adding a noise with variance of $`\sigma _{\mathrm{pix}}`$, resulting in a total noise level over the object image of $`\sqrt{N}\sigma _{\mathrm{pix}}`$. A point source, in contrast, is only subject to the noise of one pixel, $`\sigma _{\mathrm{pix}}`$. Hence, high resolution at fixed sensitivity “resolves out” a certain fraction of objects. The consequences for the source counts are clear and will be discussed shortly. These curves retain the same asymptotic behavior as before, namely a greater sensitivity to low masses at high redshift. Despite the fact that the object covers a larger number of noisy pixels as $`z`$ decreases, the optimal method is able to take proper advantage of the greater total flux density to detect low mass objects locally, just as in the case of unresolved point sources. We shall see that this does not follow for the standard detection routine (the dot–dashed lines), due to its additional surface brightness constraint (discussed below). By comparing Figures 1b and 1c, which differ only in their angular resolution, we note that for a given sensitivity, lower resolution observations are the more effective. This is traceable to the fact that the flux density of a source is dispersed over fewer (noisy) pixels than would be the case at a higher angular resolution. This indicates that low resolution observations at a given wavelength and sensitivity are to be prefered, at least for detection purposes. There is, however, a limit set by eventual source confusion.
We have just seen from Figure 1 that low surface brightness clusters are “resolved out” at high resolution. This leads to overall lower counts that are also much steeper than the equivalent for unresolved point sources. Generally speaking, the unresolved counts do not deviate too much from a Euclidean law, $`S_\nu ^{3/2}`$; on the other hand, the resolved counts can be much steeper. The examples shown in Figure 2 are in fact steeper than $`S_\nu ^2`$, indicated by the dotted lines, down to essentially the faintest flux levels attainable in immediately foreseeable observations. This is critical for optimizing an observing strategy with a fixed amount telescope time, $`T`$. Consider the common situation in which the final map noise decreases with integration time as $`1/\sqrt{t}`$; then, the solid angle covered in time $`T`$, with individual field integrations of duration $`t`$, scales with sensitivity as $`T/t\sigma _{\mathrm{pix}}^2`$. Hence, if the integrated source counts are steeper than $`\sigma _{\mathrm{pix}}^2`$, one gains objects by “going deep”, integrating longer on each individual field, rather than “going wide”, with shorter integrations covering a larger total solid angle. The important conclusion to draw from Figure 2 is then that the way to optimize the number of detected objects in a survey with arcminute resolution is by “going deep”, down to the point were the counts begin to flatten out. In our examples, this does not occur until the very lowest flux levels deemed at present reasonable. It should be emphasized that this conclusion rests on the results calculated here in the context of a self–similar cluster population. It is all the same suggestive and important in the fact that it is contrary to the conclusion one would draw based on unresolved source count calculations. The opposite holds for surveys with low angular resolution where the majority of sources remain unresolved, such as the Planck Surveyor observations.
Finally, as to be expected from the “loss” of objects at intermediate redshifts, the redshift distribution for optimally selected objects lies under the corresponding point source examples, and is somewhat flatter. All the same, the two cosmological models are easily distinguished with an enormous “leverage” at high $`z`$.
### 5.2 Standard algorithms
Standard detection routines typically identify sources as a minimum number of contiguous pixels all above a preset threshold, usually $`q_{\mathrm{st}}`$ times the pixel noise $`\sigma _{\mathrm{pix}}`$. This is not the same criteria as above, in the optimal case, because we have now established a fixed surface brightness threshold – $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}/\mathrm{\Omega }_{\mathrm{pix}}`$ – independent of object size (or luminosity). Previously, we allowed ourselves to lower this threshold for larger sources, in order to pick–up low surface brightness objects while maintaining a constant signal–to–noise; for this reason, it was an optimum detection algorithm. Here, the surface brightness is instead a fixed, while the signal–to–noise increases with object size as $`S/N=q_{\mathrm{st}}\sqrt{N}`$. A further difference is that the surface brightness cut imposes a minimum detectable mass at $`z=0`$. We obviously expect this method to detect fewer objects than the optimal approach.
Consider application of the standard algorithm to a SZ profile, empirically described in the $`\beta `$–model by $`y_\mathrm{o}`$ and $`\theta _\mathrm{c}`$. Although our final goal is to understand the selection on mass and redshift imposed by the detection criteria, it is quite useful, firstly, to gain insight into the workings of detection in terms of $`y_\mathrm{o}`$ and $`\theta _\mathrm{c}`$. As mentioned, what is actually recorded at each sample point (pixel), say by a bolometer camera, is the sky signal integrated over the beam $`B`$, which we will take to be axially symmetric. Thus, for a pixel at position $`\widehat{n}`$ (a unit vector on the sphere):
$$S_\nu ^{\mathrm{obs}}(\widehat{n})=𝑑\mathrm{\Omega }^{}i_\nu (\widehat{n}^{})B(\widehat{n}\widehat{n}^{})$$
For our calculations, we shall furthermore adopt a Gaussian beam, so that a cluster appears as a $`\beta `$–profile smeared by a Gaussian of dispersion $`\sigma _\mathrm{b}=\theta _{\mathrm{fwhm}}/\sqrt{8\mathrm{ln}2}`$,
$$B=e^{\theta ^2/2\sigma _\mathrm{b}^2}$$
where $`\theta `$ is the angle from the beam axis and $`\theta _{\mathrm{fwhm}}`$ is the beam full–width at half–maximum. Notice that we take the image to be in flux density units. By placing the coordinate origin at the cluster center, so that now $`\widehat{n}`$ is simply marked by the angular distance $`\theta `$ from the origin (small angle approximation), the beam–smeared profile of a cluster may be written as
$$S_\nu ^{\mathrm{obs}}(\theta )=y_\mathrm{o}\theta _\mathrm{c}^2j_\nu 𝒢[\theta /\theta _\mathrm{c};\sigma _\mathrm{b}/\theta _\mathrm{c}]$$
explicitly separating out a dimensionless function
$`𝒢(r;p)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑xxe^{\frac{1}{2}\left(\frac{x}{p}\right)^2}`$
$`{\displaystyle _0^{2\pi }}𝑑\varphi {\displaystyle \frac{\mathrm{\Theta }[x_\mathrm{v}^2r^2x^22xr\mathrm{cos}\varphi ]}{\left(1+r^2+x^2+2xr\mathrm{cos}\varphi \right)^\alpha }}`$
parameterized only by the ratio $`p=\sigma _\mathrm{b}/\theta _\mathrm{c}`$. The Heavyside function, $`\mathrm{\Theta }`$, cuts off the integral beyond the virial radius.
It is this smeared profile of a cluster that is subject to the detection criteria that a minimum number of connected pixels, $`N_{\mathrm{min}}`$, must lie above the threshold $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}`$. This amounts to demanding that the object image above a flux density of $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}`$ cover a minimum solid angle of $`N_{\mathrm{min}}\mathrm{\Omega }_{\mathrm{pix}}`$. Let $`\theta _{\mathrm{det}}`$ be the angular size of a cluster above the detection threshold, which may be calculated as the root of the following equation:
$$S_\nu ^{\mathrm{obs}}(\theta _{\mathrm{det}})=q_{\mathrm{st}}\sigma _{\mathrm{pix}}$$
(15)
We will say that a cluster is detected if $`\theta _{\mathrm{det}}`$ is large enough to cover $`N_{\mathrm{min}}`$ pixels. In the present analytic treatment, we will simply impose a lower limit to $`\theta _{\mathrm{det}}`$ and ignore any complications arising from the discreteness of the image. Using the function $`𝒢`$, Eq. (15) may be written in compact form as
$$𝒢[\theta _{\mathrm{det}}/\theta _\mathrm{c};\sigma _\mathrm{b}/\theta _\mathrm{c}]=\frac{q_{\mathrm{st}}\sigma _{\mathrm{pix}}}{y_\mathrm{o}\theta _\mathrm{c}^2j_\nu }\frac{\widehat{y}}{y_\mathrm{o}}\left(\frac{\sigma _\mathrm{b}}{\theta _\mathrm{c}}\right)^2$$
(16)
introducing the parameter $`\widehat{y}(q_{\mathrm{st}}\sigma _{\mathrm{pix}})/(\sigma _\mathrm{b}^2j_\nu )`$ characterizing the experimental set–up. It is clear that the solution will be given as $`\theta _{\mathrm{det}}/\theta _\mathrm{c}`$, that it will be a function of $`\sigma _\mathrm{b}/\theta _\mathrm{c}`$ and $`y_\mathrm{o}`$, and that it will be parameterized by $`\widehat{y}`$, i.e.,
$$[\theta _{\mathrm{det}}/\theta _\mathrm{c}](\sigma _\mathrm{b}/\theta _\mathrm{c},y_\mathrm{o};\widehat{y})$$
To understand the role of $`\widehat{y}`$, study the result in the limit as $`\theta _\mathrm{c}\mathrm{}`$; this will be particularly important below, when we consider the non–zero detection mass at zero redshift imposed by the surface brightness cut. In this large–object limit, $`𝒢(r;p)𝒢(r0;p0)2\pi p^2`$. We thus find that in order to be detected an object must have a central surface brightness
$$y_\mathrm{o}>\frac{\widehat{y}}{2\pi }$$
(17)
In other words, $`\widehat{y}`$ indeed embodies the surface brightness cut. This will be used shortly.
Figure 4 shows the solution over the $`(\theta _\mathrm{c}/\sigma _\mathrm{b},y_\mathrm{o})`$–plane for two reasonable values of $`|\widehat{y}|`$. To understand this figure, separate the plane into a region occupied by resolved sources – $`\theta _\mathrm{c}/\sigma _\mathrm{b}>>1`$ – and the region of point sources, $`\theta _\mathrm{c}/\sigma _\mathrm{b}<<1`$:
* Resolved sources: It is clear that by increasing $`y_\mathrm{o}`$ at fixed $`\theta _\mathrm{c}/\sigma _\mathrm{b}(>>1)`$, we see an ever increasing portion of the ICM. The cluster ‘lights-up’ until we see all of it, out to the virial radius (beyond which we assume that the gas has not been heated), and the solution flattens out at this point to the adopted value $`x_\mathrm{v}=R_\mathrm{v}/r_\mathrm{c}=10`$; beyond this, there is no more cluster to be seen. Of course, in the other direction, the object is eventually lost as we decrease $`y_\mathrm{o}`$ to the point where even the central parts of the cluster do not rise above the detection threshold.
* Point–source limit ($`\theta _\mathrm{c}/\sigma _\mathrm{b}0`$): In this extreme, the source profile becomes that of the beam, normalized to the total source flux density. This latter quantity scales as $`y_\mathrm{o}j_\nu \theta _\mathrm{c}^2`$, so that as $`\theta _\mathrm{c}/\sigma _\mathrm{b}`$ continues to decrease at fixed $`\widehat{y}`$ (i.e., holding $`\sigma _\mathrm{b}`$ constant), the imprint of the object gradually sinks below the detection threshold and $`\theta _{\mathrm{det}}0`$; this explains the cut-off at low $`\theta _\mathrm{c}/\sigma _\mathrm{b}`$ for a given central surface brightness.
So far, nothing extraordinary, but rather the standard issues of extended object detection given a particular intensity profile. Modeling more specific to the SZ effect enters only when we apply the SZ–based relations between the central surface brightness and angular size – $`y_\mathrm{o}`$ and $`\theta _\mathrm{c}`$ – and the theoretically meaningful quantities of cluster mass, $`M`$, and redshift, $`z`$. These relations allow us to translate a surface like that of Figure 4 into an equivalent surface over the $`(z,M)`$–plane, as shown for three different cases in Figure 5 (all using our adopted self–similar cluster model). Note that these latter surfaces (Figure 5) are not uniquely parameterized by $`\widehat{y}`$, because the translation from the axes in Figure 4 to mass and redshift explicitly involves $`\sigma _\mathrm{b}`$. Thus, there are now two governing parameters, which we can take to be $`\widehat{y}`$ and $`\sigma _\mathrm{b}`$:
$$\theta _{\mathrm{det}}[z,M;\widehat{y},\sigma _\mathrm{b}]$$
To study Figure 5 in detail, recall the simple scaling relations $`y_\mathrm{o}nMM\mathrm{\Delta }_{\mathrm{NL}}(z)(1+z)^3`$ and $`\theta _\mathrm{c}M^{1/3}\mathrm{\Delta }_{\mathrm{NL}}^{1/3}(z)(1+z)^1/D_{\mathrm{ang}}(z)`$, valid if the core radius scales with virial radius (it may not, but this has been assumed in the construction of the figure). Notice that mass and redshift are mixed in a nontrivial way in the expressions for surface brightness and angular extent. In particular, a cluster of given mass becomes more centrally bright towards higher redshift, due to a higher gas density (scaling $``$ the background), while its angular extent at first decreases rapidly, as $`1/z`$ at low redshift, and then approaches an asymptote, since $`(1+z)D_{\mathrm{ang}}(z)2c(H_\mathrm{o}\mathrm{\Omega }_\mathrm{o})^1`$ towards large $`z`$.
Consider firstly the low–redshift region of Figure 5b (the various effects are most clearly displayed in panel b), where the $`z`$-dependence is dominated by $`D_{\mathrm{ang}}`$; here, mass uniquely parameterizes the surface brightness, i.e., $`y_\mathrm{o}`$, while $`z`$ changes only $`\theta _\mathrm{c}`$:
* At constant $`M`$, decreasing $`z`$ increases the angular size of a cluster, so that, for resolved objects that are not completely below the surface brightness detection threshold ($`q_{\mathrm{st}}\sigma _{\mathrm{pix}}/\mathrm{\Omega }_{\mathrm{pix}}`$), $`\theta _{\mathrm{det}}`$ rises as $`1/z`$; this corresponds to the $`\theta _{\mathrm{det}}/\theta _\mathrm{c}const`$ behavior in Figure 4. For smaller objects, of low mass, the beam profile determines the scaling of $`\theta _{\mathrm{det}}`$ with $`z`$, and this corresponds to the point–source limit of Figure 4.
* At fixed (low) redshift, low mass objects eventually fall below the surface brightness limit, and $`\theta _{\mathrm{det}}`$ reaches zero; on the other hand, massive clusters are ‘lit–up’ out to their virial radius, at which point $`\theta _{\mathrm{det}}`$ attains $`\theta _{\mathrm{vir}}=x_\mathrm{v}\theta _\mathrm{c}`$, which grows as $`M^{1/3}`$.
An important aspect of detection procedure, mentioned at the beginning of this section and now evident from Figure 5, is the existence of a minimum detectable mass in the limit of zero redshift (on the $`M`$–axis). This characteristic mass is established by the surface brightness threshold – $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}/\mathrm{\Omega }_{\mathrm{pix}}`$, and its existence represents a fundamental difference with the previous case of optimal detection, where arbitrarily low mass (low surface brightness) clusters where picked up if they were large enough, i.e., very close at $`z=0`$. We already saw in Eq. (17) how $`\widehat{y}`$ summarizes the surface brightness constraint. The low redshift detection mass limits seen in Figure 1 and here in Figure 5 are indeed reproduced numerically from Eq. (17), once Eq. (11) is used to convert $`y_\mathrm{o}`$ into a mass at $`z=0`$.
At redshifts approaching unity and beyond, $`M`$ and $`z`$ are fully mixed in the expressions for $`y_\mathrm{o}`$ and $`\theta _\mathrm{c}`$:
* Massive clusters well above the surface brightness threshold increase in surface brightness with redshift to the point where they are completely seen, all the way out to $`\theta _{\mathrm{vir}}`$; at even larger $`z`$, $`\theta _{\mathrm{det}}`$ reflects the gradual fall–off to the asymptote set by the angular–size distance. This explains the ridge running down the surface in Figure 5 around $`z=1`$. Less massive clusters, on the other hand, only reach the point of full illumination at higher $`z`$, well into the asymptotic behavior of $`\theta _{\mathrm{vir}}`$, and hence the ridge tends to be washed out at the low mass end. Finally, the central surface brightness of very low mass clusters falls below the detection threshold at ever larger redshifts, i.e., the boundary $`\theta _{\mathrm{det}}=0`$ moves outward in $`z`$ as $`M`$ is decreased.
Compare now the three panels of Figure 5. We observe the greater sensitivity at $`2`$ mm, due to the spectrum of the SZ effect, by the fact that a given cluster of mass $`M`$ and $`z`$ produces a smaller $`\theta _{\mathrm{det}}`$ at $`1`$ cm wavelength in panel b (notice the change in scale along the $`M`$–axis between panels a and b). The same remarque applies to the greater sensitivity, at a given noise level, of the lower resolution observations exemplified in panel c. These characteristics will be inherited by the detection mass curves, our next topic.
#### 5.2.1 Detection mass as a function of redshift
Since $`\theta _{\mathrm{det}}`$ increases monotonically with $`M`$, the contours displayed on the top and bottom faces of Figure 5 represent curves of minimum detectable mass, $`M_{\mathrm{det}}^{\mathrm{st}}(z;\theta _{\mathrm{det}})`$, each one for a different detection threshold defined by different values of $`\theta _{\mathrm{det}}`$ (indicated in arcminutes on the contours in the figure). All of these contours, however, are defined for the same value of $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}`$ set by the governing parameters $`\widehat{y}`$ and $`\sigma _\mathrm{b}`$ (or $`\theta _{\mathrm{fwhm}}`$). In contrast to the optimal routine, a detection in the standard case is specified by not one, but two parameters – the pair $`(q_{\mathrm{st}},\theta _{\mathrm{det}})`$. This embodies the fact that a detection must satisfy two criteria: a minimum flux density, $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}\theta _{\mathrm{det}}^2/\mathrm{\Omega }_{\mathrm{pix}}`$, and a minimum surface brightness, $`q_{\mathrm{st}}\sigma _{\mathrm{pix}}/\mathrm{\Omega }_{\mathrm{pix}}`$. In practice, the choice of values for $`q_{\mathrm{st}}`$ and $`\theta _{\mathrm{det}}`$ may be somewhat of a black art, but once made it specifies the survey’s characteristic $`M_{\mathrm{det}}^{\mathrm{st}}(z)`$. For the ensuing examples, I make the choice motivated by the following considerations: Note that the signal–to–noise of a detection $`S/N=q_{\mathrm{st}}\sqrt{N_{\mathrm{min}}}`$. The detection criteria imposed as $`\theta _{\mathrm{det}}=\sqrt{N_{\mathrm{min}}}\theta _{\mathrm{pix}}`$ may then be expressed in terms of the $`S/N`$:
$$\theta _{\mathrm{det}}=\frac{1}{q_{\mathrm{st}}}\left(\frac{S}{N}\right)\left(\frac{\theta _{\mathrm{pix}}}{\theta _{\mathrm{fwhm}}}\right)\theta _{\mathrm{fwhm}}$$
(18)
For reference, recall that in the optimal approach the parameter $`q_{\mathrm{opt}}`$ was exactly the $`S/N`$. Now, at fixed signal–to–noise, a larger $`q_{\mathrm{st}}`$ leads to a smaller $`N_{\mathrm{min}}`$ (i.e., $`\theta _{\mathrm{det}}`$), which facilitates the detection of fainter point sources, because their flux density is buried in less noise (fewer pixels). On the other hand, a large value of $`q_{\mathrm{st}}`$ disfavors finding low surface brightness objects; thus, a compromise is called for. One reasonable choice would be $`\theta _{\mathrm{det}}=(1/2)\theta _{\mathrm{fwhm}}`$, corresponding to a minimum detection $`S/N3`$, with $`q_{\mathrm{st}}3`$ and $`\theta _{\mathrm{pix}}/\theta _{\mathrm{fwhm}}1/2`$. I henceforth adopt these values for the following examples, which now completely specifies our detection routine.
The dot–dashed lines in Figure 1 show the resulting standard mass detection curves. They all lie above the optimal detection curves, implying a lower overall sensitivity, as expected; and as in the previous cases, they fall with $`z`$. The low resolution examples in Figure 1c show a slight turn–down at low redshift, but under no circumstances will they ever reach the origin at $`z=0`$, as do the unresolved and optimal resolved detection curves: as already mentioned, there always remains a non–zero detection mass at low $`z`$ in the standard case, due to the surface brightness cut. This constraint may be neatly summarized, using our earlier result, as $`y_\mathrm{o}[M_{\mathrm{det}}^{\mathrm{st}}(z=0),0]=\widehat{y}/2\pi `$ (but it must be noted that this relies on our use of the self–similar cluster model). The loss of close–by, low–mass halos is particularly noteworthy for the study of low–mass halos; to find them in an SZ survey, one of the important potentials of such efforts, will require special “tuning” of detection criteria, to more closely approach the optimal routine.
#### 5.2.2 Counts
Our next goal is to use the detection mass to calculate the cluster counts and redshift distributions. It is worth noting in passing that one can envision several different kinds of source counts: as a function of the value of $`\theta _{\mathrm{det}}`$, as a function of some aperture flux density (fixed or isophotal) or as a function of survey sensitivity, $`\sigma _{\mathrm{pix}}`$. Only the last one, however, is useful for optimizing an observing strategy, i.e., to answer the question of whether it is better to ‘go deep’, or to ’go wide’ when performing a survey. We are thus brought to consider in detail the detection mass as a function of detector sensitivity – $`M_{\mathrm{det}}^{\mathrm{st}}(z;\sigma _{\mathrm{pix}})`$. The most direct efficient way to do this for a large number of different sensitivities is by returning to Eq.(16), fixing $`\theta _{\mathrm{det}}(=1/2\theta _{\mathrm{fwhm}})`$ and then finding the detection mass as the root of the equation for each $`z`$. This avoids having to calculate the entire $`\theta _{\mathrm{det}}`$ surface for each $`\sigma _{\mathrm{pix}}`$ ($`\widehat{y}`$) just to extract a single contour.
Using the result of this operation in Eq.(13), we find the counts and redshift distributions displayed in Figures 2 and 3 as the dot–dashed lines. Not surprisingly, the counts are lower at a given sensitivity than the corresponding optimal counts, and they are also slightly steeper. Similar remarks apply to the results of Figure 3. All of this is easily understood from the loss of low–surface brightness objects relative to the optimal routine. The essential conclusion concerning observation strategy remains the same: down to low flux densities, deeper integrations should yield more sources.
## 6 Discussion
The effects of resolving clusters must be properly modeled to understand the capabilities of possible ground–based surveys, as is clear from, for example, Figure 2: predicted counts are lower and steeper for resolved clusters relative to hypothetical SZ point sources. Besides lowering the expectations for the number of detectable sources, these results also suggest that deep integrations are more efficient than wide and shallow ones. The actual number of clusters expected for realistic ground–based performance are model dependent. For a self–similar cluster population, one could reasonably expect between 10–100 clusters/sq. deg. down to $`0.1`$ mJy at $`2`$ mm and with $`\theta _{\mathrm{fwhm}}=1`$ arcmin, as shown in Figure 2a. This number depends in addition on the source detection method employed: the standard routine counts may perhaps be considered realistic, while the optimal method counts indicate instead the best one could hope to achieve. At $`1`$ cm, for equivalent sensitivity and at the lower resolution of $`\theta _{\mathrm{fwhm}}=2`$ arcmin, one expects an order of magnitude lower surface density (Figure 1c). In sum, a square degree survey at $`2`$ mm could yield $`10100`$ detections depending on the exact cluster model and the detection algorithm; a 10 square degree survey at $`1`$ cm to the same sensitivity ($`0.1`$ mJy) could produce similar numbers. Both types of survey may soon be achievable, with instruments similar to BOLOCAM (Glenn et al. 1998) or a detected interferometer array (Carlstrom et al. 1999)
One of the primary interests of opening this new window onto the Universe is to the search for high redshift clusters. The details of resolved cluster detection do not change the important and tell–tail difference between the redshift distributions in different cosmological models: the expected number of high redshift clusters is a sensitive function of $`\mathrm{\Omega }_\mathrm{o}`$, as demonstrated by the redshift distributions given in Figure 3. Observations of such redshift distributions should prove a valuable tool for constraining $`\mathrm{\Omega }_\mathrm{o}`$ and for understanding evolution of the cluster environment.
There are several important issues that have not been dealt with in the present work. One concerns eventual source confusion, an effect that depends on the beam size and the exact value of the counts. This effect may very well be important even on arcminute scales, as noted by Aghanim et al. (1997). As these authors also point out, the issue is complicated by the fact that, due to the extended nature of clusters, one must also contend with source blending. Detailed modeling of these effects really requires simulations.
Another important issue not addressed in the present work concerns the question of radio source contamination. With sufficient frequency coverage, on can always identify SZ sources by their unique spectrum. Most often, though, spectral coverage is limited and contamination may become problematic. Its importance depends on the observation frequency, and the counts at millimeter wavelengths are in fact a subject of current fundamental research; thus, the nature of contamination at in the millimeter is much more model dependent than in the centimeter.
Finally, I note once again that the present work is based on a simple cluster model, because the principal motivation has been to understand the nature of resolved cluster detection by comparison to the more classic unresolved case. Any attempt at a more exact examination of the number counts and redshift distributions requires more detailed cluster modeling. Such work would, in addition, permit an interesting comparison of the relative efficiencies of SZ and X–ray observations to finding high redshift clusters, in practice. The SZ effect is clearly inherently more efficient, but to really address this question, one should consider the actual achievable sensitivities of the two approaches.
## 7 Conclusions
There are clear and important differences in the conclusions one draws concerning SZ surveys depending on whether clusters are considered as point sources or as extended. For low resolution surveys, such as expected from the Planck Surveyor, most clusters will remain unresolved; however, when discussing the arcminute resolution more applicable to possible future ground–based surveys, we have seen that it is important to model the clusters as resolved sources in order to properly understand the nature of detectable objects. For a given sensitivity, high angular resolution “resolves out” some clusters, lowering and steepening the final source counts. Relative to optimal resolved detection, standard algorithms tend to in addition loose low mass, low redshift clusters due to their imposed surface brightness cut, further steepening and lowering the counts. With a fixed total observation time and a given frequency and angular resolution, we have seen that our results imply that deep integrations yield more objects than shallow ones covering a large area.
Some important issues still to be explored concern the questions of source confusion and blending, and radio source contamination. A detailed comparison of SZ and X–ray surveys would also be of interest, which implies more detailed cluster modeling than employed here.
All the same, the numbers from the self–similar cluster model should be, within all the present uncertainties of these predictions, illustrative of what may be soon achieved from the ground. It appears that both in the millimeter and in the cm, ground–based SZ surveys could be capable of detecting up to $`100`$ clusters in total, a respectable statistical catalog.
###### Acknowledgements.
I am very pleased to thank K. Romer and J. Mohr for their SZ workshop at the centennial AAS meeting in Chicago, which was the starting point for this work. I am also grateful for the hospitality of P. Rosati at the European Southern Observatory and of J. Willick at Stanford University where some of this work was carried out. Thanks to B. Keating for helpful discussions. |
warning/0001/quant-ph0001094.html | ar5iv | text | # Dark-State Polaritons in Electromagnetically Induced Transparency
## Abstract
We identify form-stable coupled excitations of light and matter (“dark-state polaritons”) associated with the propagation of quantum fields in Electromagnetically Induced Transparency. The properties of the dark-state polaritons such as the group velocity are determined by the mixing angle between light and matter components and can be controlled by an external coherent field as the pulse propagates. In particular, light pulses can be decelerated and “trapped” in which case their shape and quantum state are mapped onto metastable collective states of matter. Possible applications of this reversible coherent-control technique are discussed.
Dark resonances and electromagnetically induced transparency (EIT) can be used to make a resonant, optically opaque medium transparent by means of quantum interference. Associated with the induced transparency is a dramatic modification of the refractive properties of the media. These can result, for instance, in very slow group velocities . In the present contribution we study the propagation of quantum fields in EIT media. We demonstrate the existence of formstable quantum excitations associated with such propagation, which we term “dark-state polaritons”. The polaritons are coherent superpositions of photonic and Raman-like matter branches. We show that their group velocity is directly related to the ratio of the two contributions. This ratio can be externally controlled by adiabatically changing a coherent control field as the pulse propagates. In particular, dark-state polaritons can be stopped and re-accelerated in such a way that their shape and quantum state are preserved. In this process the quantum state of light is ideally transfered to collective atomic excitations and vise versa.
The possibility to coherently control the propagation of quantum light pulses via dark-state polaritons opens up interesting applications involving the generation of non-classical states of atomic ensembles (in squeezed or entangled states), reversible quantum memories for light waves , and high resolution spectroscopy . Furthermore, the combination of the present technique with studies on few-photon nonlinear optics can be used, in principle, for processing of quantum information stored in collective excitations of matter. Finally, the present technique may provide an interesting tool to study quantum scattering phenomena in systems involving coherent cold collisions. In this regard the present work opens a link between nonlinear optics for light waves and nonlinear atom optics. E.g. an interaction (or entanglement) between light waves can be induced by a collisional interaction of atoms (e.g. s-wave scattering); alternatively an interaction between atoms can be induced via optical nonlinearities.
We consider a medium consisting of $`\mathrm{\Lambda }`$-type 3-level atoms with two meta-stable lower states as shown in Fig. 1. A quantum field described by the slowly-varying dimensionless operator
$`\widehat{E}(z,t)={\displaystyle \underset{k}{}}a_k(t)\mathrm{e}^{ikz}\mathrm{e}^{i\frac{\nu }{c}(zct)}`$ (1)
couples resonantly the transition between the ground state $`|b`$ and the excited state $`|a`$. $`\nu =\omega _{ab}`$ is the carrier frequency of the optical field. The upper level $`|a`$ is furthermore coupled to the stable state $`|c`$ via a coherent control field with the slowly-varying, real Rabi-frequency $`\mathrm{\Omega }(t)`$. For the purposes of the present discussion the external field can be treated classically. We assume that initially (i.e before the quantum pulse arrives) all atoms are in their ground states $`|b_j`$. To describe the quantum properties of the medium, we use collective, slowly varying atomic operators, appropriately averaged over small but macroscopic volumes containing $`N_z1`$ particles at position $`z`$,
$$\widehat{\sigma }_{\alpha \beta }(z,t)=\frac{1}{N_z}\underset{j=1}{\overset{N_z}{}}|\alpha _j\beta _j|\mathrm{e}^{i\omega _{\alpha \beta }t}.$$
(2)
The interaction between light and atoms is governed by the Hamiltonian
$$\widehat{V}=N\frac{\mathrm{d}z}{L}\left(\mathrm{}g\underset{k}{}a_k\mathrm{e}^{ikz}\widehat{\sigma }_{ab}(z)+\mathrm{}\mathrm{\Omega }\widehat{\sigma }_{ac}(z)\right)+h.c.$$
(3)
Here $`g=\mathrm{}\sqrt{\frac{\nu }{2\mathrm{}ϵ_0V}}`$ is the atom-field coupling constant with $`\mathrm{}`$ being the dipole moment of the $`ab`$ transition and $`V`$ the quantization volume. $`N`$ is the number of atoms in this volume and $`L`$ its length in $`z`$ direction.
The evolution of the Heisenberg operator corresponding to the optical field can be described in slowly varying amplitude approximation by the propagation equation
$`\left({\displaystyle \frac{}{t}}+c{\displaystyle \frac{}{z}}\right)\widehat{E}(z,t)=igN\widehat{\sigma }_{ba}(z,t).`$ (4)
The atomic evolution is governed by a set of Heisenberg-Langevin equations
$`{\displaystyle \frac{}{t}}\widehat{\sigma }_{\mu \nu }=\gamma _{\mu \nu }\sigma _{\mu \nu }+{\displaystyle \frac{i}{\mathrm{}}}[\widehat{V},\widehat{\sigma }_{\mu \nu }]+F_{\mu \nu },`$ (5)
where $`\gamma _{\mu \nu }`$ are the transversal decay rates and $`\widehat{F}_{\mu \nu }`$ are $`\delta `$-correlated Langevin noise operators.
We now assume that the Rabi-frequency of the quantum field is initially much smaller than $`\mathrm{\Omega }`$ and that the number of photons in the input pulse is much less than the number of atoms. We will show that the Rabi-frequency of the quantum field will then be much smaller than $`\mathrm{\Omega }`$ at all times. In such a case the atomic equations can be treated perturbatively in $`\widehat{E}`$. In zeroth order only $`\widehat{\sigma }_{bb}=\mathrm{𝟏}`$ is different from zero and in first order one finds
$`\widehat{\sigma }_{ba}`$ $`=`$ $`{\displaystyle \frac{i}{\mathrm{\Omega }(t)}}{\displaystyle \frac{}{t}}\widehat{\sigma }_{bc},`$ (6)
$`\widehat{\sigma }_{bc}`$ $`=`$ $`{\displaystyle \frac{g\widehat{E}}{\mathrm{\Omega }}}{\displaystyle \frac{i}{\mathrm{\Omega }}}\left[\left({\displaystyle \frac{}{t}}+\gamma _{ba}\right)\left({\displaystyle \frac{i}{\mathrm{\Omega }}}{\displaystyle \frac{}{t}}\widehat{\sigma }_{bc}\right)+\widehat{F}_{ba}\right].`$ (7)
In the above equations we disregarded a (small) decay of the Raman coherence ($`\gamma _{bc}`$).
The propagation equations simplify considerably if we assume a sufficiently slow change of $`\mathrm{\Omega }`$, i.e. adiabatic conditions . Introducing a normalized time $`\stackrel{~}{t}=t/T`$ where $`T`$ is a characteristic time scale and expanding the r.h.s. of (7) in powers of $`1/T`$ we find in lowest non-vanishing order
$$\widehat{\sigma }_{bc}(z,t)=g\frac{\widehat{E}(z,t)}{\mathrm{\Omega }(t)}.$$
(8)
Note that $`\widehat{F}_x(t)\widehat{F}_y(t^{})\delta (tt^{})=\delta (\stackrel{~}{t}\stackrel{~}{t}^{})/T`$. Thus in the perturbative and adiabatic limit the propagation of the quantum light pulse is governed by the equation
$`\left({\displaystyle \frac{}{t}}+c{\displaystyle \frac{}{z}}\right)\widehat{E}(z,t)={\displaystyle \frac{g^2N}{\mathrm{\Omega }(t)}}{\displaystyle \frac{}{t}}{\displaystyle \frac{\widehat{E}(z,t)}{\mathrm{\Omega }(t)}}.`$ (9)
If $`\mathrm{\Omega }`$ is constant, the term on the r.h.s. simply leads to a modification of the group velocity of the quantum field according to $`v_g=c/(1+\frac{g^2N}{\mathrm{\Omega }^2})`$. In the general case the field equation of motion will acquire an additional term proportional to $`(\dot{\mathrm{\Omega }}/\mathrm{\Omega })\widehat{E}`$ which describes reversible changes in quantum amplitudes due to stimulated Raman scattering.
One can obtain a very simple solution of eq.(9) by introducing a new quantum field $`\widehat{\mathrm{\Psi }}(z,t)`$ via the canonical transformation
$`\widehat{\mathrm{\Psi }}(z,t)=\mathrm{cos}\theta (t)\widehat{E}(z,t)\mathrm{sin}\theta (t)\sqrt{N}\widehat{\sigma }_{bc}(z,t),`$ (10)
$`\mathrm{cos}\theta (t)={\displaystyle \frac{\mathrm{\Omega }(t)}{\sqrt{\mathrm{\Omega }^2(t)+g^2N}}},\mathrm{sin}\theta (t)={\displaystyle \frac{g\sqrt{N}}{\sqrt{\mathrm{\Omega }^2(t)+g^2N}}}.`$ (11)
$`\widehat{\mathrm{\Psi }}`$ obeys the following equation of motion
$`\left[{\displaystyle \frac{}{t}}+c\mathrm{cos}^2\theta (t){\displaystyle \frac{}{z}}\right]\widehat{\mathrm{\Psi }}(z,t)=0,`$ (12)
which describes a shape-preserving propagation with velocity $`v=v_g(t)=c\mathrm{cos}^2\theta (t)`$:
$$\widehat{\mathrm{\Psi }}(z,t)=\widehat{\mathrm{\Psi }}\left[zc_0^td\tau \mathrm{cos}^2\theta (\tau ),t=0\right].$$
(13)
Several interesting properties of the new field should be noted. First of all, by introducing a plain-wave decomposition $`\widehat{\mathrm{\Psi }}(z,t)=_k\widehat{\mathrm{\Psi }}_k(t)\mathrm{e}^{ikz}`$ one finds that the mode operators $`\widehat{\mathrm{\Psi }}_k`$ and $`\widehat{\mathrm{\Psi }}_k^{}`$ obey the commutation relations
$`[\widehat{\mathrm{\Psi }}_k,\widehat{\mathrm{\Psi }}_k^{}^+]=\delta _{k,k^{}}\left[\mathrm{cos}^2\theta +\mathrm{sin}^2\theta {\displaystyle \frac{1}{N}}{\displaystyle \underset{j}{}}(\widehat{\sigma }_{bb}^j\widehat{\sigma }_{cc}^j)\right].`$ (14)
In the linear limit considered here, where the number density of photons is much smaller than the density of atoms, $`\widehat{\sigma }_{bb}^j1,\widehat{\sigma }_{cc}^j0`$. Thus the new field possesses bosonic commutation relations and we can associate with it bosonic quasi-particles (polaritons). Furthermore one immediately verifies that all number states created by $`\widehat{\mathrm{\Psi }}_k^{}`$ are dark-states :
$`|D_n^k={\displaystyle \frac{1}{\sqrt{n!}}}\left(\widehat{\mathrm{\Psi }}_k^{}\right)^n|0|b_1\mathrm{}b_N,`$ (15)
where $`|0`$ denotes the field vacuum. In particular, the states $`|D_n^k`$ do not contain the excited atomic state and are thus immune to spontaneous emission. Furthermore, they are eigenstates of the interaction Hamiltonian with eigenvalue zero, $`\widehat{V}|D_n^k=0`$. For these reasons we call the quasi-particles “dark-state polaritons”.
To summarize, we have found a shape-preserving, polariton-like superposition $`\widehat{\mathrm{\Psi }}`$ of an electromagnetic field and collective Raman coherences. This excitation is not of soliton type since no special pulse-shape or pulse area is required. It is related to the classical adiabaton solutions of pulse-pair propagation in $`\mathrm{\Lambda }`$-type media in the limit of one strong and one weak field. We emphasize however that the field can here be in any quantum state. In particular it does not need to have a coherent component with a well defined phase.
One of the most interesting aspects of dark-state polaritons is the possibility to coherently control their properties by changing $`\mathrm{\Omega }(t)`$. For example, by adiabatically rotating $`\theta (t)`$ from $`0`$ to $`\pi /2`$ one can decelerate and stop an input light pulse. It is remarkable that in this process pulse shape and quantum state of the initial light pulse are mapped onto collective, metastable states of matter in which they are stored. Likewise the dark-state polariton can be re-accelerated to the vacuum speed of light; in this process the stored quantum states is transferred back to the field. This is illustrated in Fig. 2, where we have shown the coherent amplitude of a dark-state polariton which results from an initial light pulse as well as the corresponding field and matter components. One recognizes that the pulse shape is preserved and that the stopping corresponds to a transfer from field to atomic Raman excitations. Explicitly, the mapping of the quantum states corresponds to the following unitary transformation:
$`\left({\displaystyle \underset{k,l,m\mathrm{}}{}}\xi _{k,l,m\mathrm{}}a_k^{}a_l^{}a_m^{}\mathrm{}|0\right)|b_1\mathrm{}b_N`$ (16)
$`\left({\displaystyle \underset{k,l,m\mathrm{}}{}}\xi _{k,l,m\mathrm{}}\sqrt{N}\sigma _{cb}^k\sqrt{N}\sigma _{cb}^l\sqrt{N}\sigma _{cb}^m\mathrm{}|b_1\mathrm{}b_N\right)|0,`$ (17)
as can be verified using expression (15) for the polariton state vectors.
The coherent transfer of quantum states between light and matter opens interesting prospectives for the generation of non-classical atomic ensembles in squeezed and entangled states, high-precision spectroscopy with resolution beyond the standard quantum limit as well as reversible quantum memories. Furthermore, by trapping correlated photons in separate media entangled states of separated atomic ensembles can be created. With respect to these applications the present paper is complementary to our earlier studies in which we showed that quantum states of light can be mapped onto Dicke-like collective states of an EIT medium in an optical resonator . The quantum states of matter generated in the case of the present paper are more complicated; however trapping the light in a traveling-wave geometry does not require special shaping of the classical driving pulses (quantum impedance matching), which is necessary in a cavity configuration.
We also note related studies on quantum memories for light involving mapping the quantum state of the field onto atoms by dissipative absorption . In contrast to these approaches the adiabatic passage technique used here allows for a complete and reversible excitation transfer of arbitrary quantum wavepackets.
Finally, our approach is also different from the mechanism suggested recently in , in which “freezing” of the light pulse in a laboratory frame was proposed using moving atoms.
The above analysis involves a perturbation expansion, an adiabatic approximation and disregards the decay of Raman coherence. In what follows the validity of these approximations is discussed. First of all, we note that making use of (8) one finds: $`g^2\widehat{E}^+\widehat{E}/|\mathrm{\Omega }|^2=\widehat{\sigma }_{cb}\widehat{\sigma }_{bc}`$. I.e., the ratio of the average intensities of quantum and control field is proportional to that of the matter field $`\widehat{\sigma }_{cc}`$. If the initial number of photons in the quantum field is much less than the number of atoms, $`\widehat{\sigma }_{cc}`$ is always much smaller than unity. Therefore the mean intensity of the quantum field remains small compared to that of the control field even when the latter is turned to zero.
In order to check the validity of the adiabatic approximation we consider the first correction to $`\widehat{\sigma }_{bc}`$:
$`\widehat{\sigma }_{bc}{\displaystyle \frac{g\widehat{E}}{\mathrm{\Omega }}}+{\displaystyle \frac{1}{\mathrm{\Omega }}}\left({\displaystyle \frac{}{t}}+\gamma _{ba}\right){\displaystyle \frac{1}{\mathrm{\Omega }}}{\displaystyle \frac{}{t}}{\displaystyle \frac{g\widehat{E}}{\mathrm{\Omega }}}+\mathrm{}`$ (18)
The non-adiabatic correction in (18) leads to a spectral narrowing (pulse spreading) of the quantum field due to the finite bandwidth of the transparency window , which results in a “pulse”-matching of quantum and classical control field . Using the adiabatic solution (13), one can verify that these corrections are small for propagation distances:
$`zz_{max}={\displaystyle \frac{g^2N}{\gamma _{ab}}}\times {\displaystyle \frac{L_p^2}{c}},`$ (19)
where $`L_p`$ is the length of the input pulse. Hence, in order to trap a pulse with negligible losses, it is required that
$`{\displaystyle \frac{g^2NL_p}{c\gamma _{ab}}}1.`$ (20)
This condition contains the number of atoms which is a signature of collective interactions. It should be contrasted to the strong-coupling condition corresponding to a quantum state transfer in cavity QED . We note, in particular, that in the optically dense medium the adiabatic condition (20) is much easier to implement.
The effect of the Raman coherence decay can be easily estimated using the explicit expression for the generated matter states (16). It is clear that the collective states containing $`n_e`$ atomic excitations will dephase at a rate $`\gamma _{bc}n_e`$. Hence, the time of the storage should be limited to $`t_s(\gamma _{bc}n_e)^1`$ to avoid decoherence .
In the discussion above we have considered the case where the control field only depends on time. This is valid, for instance, when the control field propagates in a direction perpendicular to that of the quantum field. In experiments involving hot atomic vapors co-propagation is required, however, in order to cancel Doppler broadening of the two-photon transition. In this case propagation effects of the control field need to be considered. If the quantum field is weak, the control field propagates as in free space and thus $`\mathrm{\Omega }(z,t)=\mathrm{\Omega }(tz/c)`$. In this case one finds:
$`\left({\displaystyle \frac{}{t}}+c\mathrm{cos}^2\theta (z,t){\displaystyle \frac{}{z}}\right){\displaystyle \frac{\widehat{E}(z,t)}{\mathrm{\Omega }(z,t)}}=0.`$ (21)
Since the group velocity is now also $`z`$-dependent, trapping of the pulse does not preserve the shape exactly. Nevertheless it is evident that trapping and a reversible transfer of the quantum state from light to atoms are still possible. In experiments, however, a more practical approach can be taken in which a light pulse enters the medium already with $`v_g^0c`$. In such a case retardation of the control field can be ignored and one has $`\mathrm{\Omega }(tz/c)\mathrm{\Omega }(t)`$. Since the index of refraction is close to unity there will be no reflection losses at the entrance plane. However the polariton pulse becomes spatially compressed according to $`L_p/L_p^0=v_g^0/c`$, and its amplitude grows according to the boundary condition $`\widehat{\mathrm{\Psi }}(0,t)=\sqrt{c/v_g^0}\widehat{E}(0,t)`$. In this way, the total energy of the polariton field inside the medium is equval to the energy of the light field outside. After entering the medium the polaritons can be manipulated as discussed above.
In conclusion we have shown that it is possible to control the propagation of quantum pulses in optically thick $`\mathrm{\Lambda }`$-type media. This coherent control mechanism is based on dark-state polaritons associated with EIT. In particular, a quantum light pulse can be “trapped”, in which case its shape and quantum state are preserved in stationary atomic excitations. The matter-like polariton can then be re-accelerated and converted back into a photon pulse. These properties of dark-state polaritons can be used for squeezing and entanglement transfer from light to atoms. Furthermore, we anticipate interesting applications involving nonlinear interactions between such polaritons.
We thank M.O. Scully for many stimulating discussions. This work was supported by the National Science Foundation. |
warning/0001/cond-mat0001276.html | ar5iv | text | # CA Models for Traffic Flow: Comparison with Empirical Single-Vehicle Data
## 1 Introduction
For a long time the modelling of traffic flow phenomena was dominated by two theoretical concepts (for a review, see e.g. ): Microscopic car-following models and macroscopic models based on the analogy between traffic flow and the dynamics of compressible viscous fluids. Both approaches are still used widely by traffic engineers but for practical purposes they are often not suitable, e.g. an efficient implementation for computer simulations of large networks is not possible. Macroscopic models use a large number of parameters which have partly no counterpart within empirical investigations. Moreover, the information which can be obtained using macroscopic models is incomplete in the sense that quantities concerning individual cars cannot be introduced or derived directly.
In order to fill this gap cellular automata (CA) models have been invented . CA’s are microscopic models which are by design well suited for large-scale computer simulations. A comparison of the simulations with empirical data shows that already very simple approaches give meaningful results. In particular they can be used to simulate dense networks like cities which are controlled by the dynamics at the intersections. However, for highway traffic a more detailed description of the dynamics seems to be necessary.
Recent empirical results show the existence of metastable states in traffic dynamics and the occurrence of synchronised flow , which can be identified by vanishing cross-correlations of the local density and the local flow . Moreover, a detailed analysis of single-vehicle data revealed important facts for the microscopic modelling of traffic. The time-headway distribution shows two characteristic peaks. Small time-headways ($`0.8sec`$) are a result of cars or clusters of cars moving with small headway but large velocity, a time-headway of $`2sec`$ can be identified with the drivers efforts for safety: It is recommended to drive with a distance of $`2`$ seconds. Additionally, the distance-headway gives the most important information for the adjustment of the car’s speed for the correct description of the car-car interaction. This is introduced in several models by the so-called optimal velocity (OV) curve. It has been shown that one universal optimal velocity curve for all density regimes does not exist, but individual curves for different densities can be calculated (see the article of Neubert et al. in these proceedings for more details). In fact, some model extensions of the CA model proposed by Nagel and Schreckenberg (NaSch) exist which are capable to reproduce metastable states or small time-headways , but up to now it it not possible to generate synchronised traffic and the correct microscopic properties mentioned above.
Here we propose a new CA model generalising the NaSch model and some earlier extensions. We compared our simulations with the corresponding data used in . The simulation data are evaluated by an artificial counting loop, i.e. we measured the speed and the time-headway of the vehicles at a given link of the lattice. This data set is analysed using the methods suggested in . In particular the density is calculated via the relation $`\rho =J/v`$ where $`J`$ and $`v`$ are the mean flow and the mean velocity of cars passing the detector in a time interval of $`1`$ minute. This dynamical estimate of the density gives correct results only if the velocity of the cars between two measurements is constant, but for accelerating or braking cars, e.g. in stop-and-go traffic, the results do not coincide with the real occupation. In addition to the aggregated data also the single-vehicle data of each passing car are analysed. Although the empirical data have been obtained on a two-lane highway, the simulations are performed on a single-lane road with one type of cars, because the empirical results show no systematic lane dependence which is a consequence of the applied speed limit.
In Sect. 2 we give a brief description of the new model definition consisting of the NaSch-rules and some extensions. Section 3 compares the simulation results of the new model with the corresponding empirical data. Finally, Sect. 4 concludes with a short summary and discussion.
## 2 New Approach
Traffic networks can be classified as complex systems of a multitude of individual interacting agents. In contrast to urban traffic where the flow is dominated by intersections, signals etc., car characteristics like different maximal velocities, acceleration capabilities and car lengths become important on highways. In order to allow for a more realistic modelling of these characteristics we reduce the cell length of the standard NaSch model (see for a detailed description of the model) to a length of $`l=1.5m`$. The acceleration and randomisation remains unaltered with one site per time step of $`1sec`$ which leads to a velocity discretisation of $`5.4km/h`$ which is slightly above ”comfortable” acceleration of about $`1m/sec^2`$ .
The update rules of the new model combine the original NaSch model and some recent extensions, namely a slow-to-start rule and an anticipation term . The slow-to-start rule allows to tune the velocity of the upstream front of a traffic jam directly. It turns out that for a realistic choice of the parameters the outflow of a jam does not achieve the capacity of the road. This empirical fact is known to lead to the existence of metastable states.
The next step towards a more realistic description of especially highway traffic is to introduce anticipation effects, i.e. the adjustment of speed also takes into account the expected behaviour of the leading vehicle. Anticipation leads to a much more efficient lane-usage in multi-lane traffic. Although both modifications significantly increase the realism of the simulations a complete description of traffic highway traffic is not yet possible. The main problem with the existing discrete models is that they fail to reproduce platoons of slow moving vehicles. These patterns are not as stable as in real traffic, i.e. the models overestimate the probability to form large compact jams.
This deficiency motivated us to prolong the range of interactions if a braking maneuver of the leading vehicle happens or, more figurative, we equipped the vehicles with brake lights. The event driven interaction leads to a timely adjustment of speeds and therefore to a more coherent movement of the vehicles in dense traffic. We implemented the reaction to a brake light simply by an increased randomisation parameter $`p_b>p`$.
## 3 Validation of the Model
Obviously, the fundamental diagram of the new model coincides very well with the empirical data (Fig. 1). In comparison we observe a more narrow distribution of densities. This further narrowing is simply an artifact of the discretisation of the velocities which determines the upper limit of detectable densities.
The slow-to-start rule has been introduced in order to reduce the outflow of a jam. This rule is responsible for the formation of large jams at high densities. To measure the outflow we used a megajam initialisation. Obviously, the outflow is reduced considerably (inset of Fig. 1). Using the autocorrelation function of the measured local density of a system initialised with a megajam it is possible to determine the jam velocity. A detailed analysis leads to a value of about $`12.75km/h`$.
As a next step for the validation of the model we compared single-vehicle data of the simulation with the corresponding empirical data. In order to give a correct comparison of our simulation data with the data of Neubert et al. we tried to identify the three traffic states found in . We therefore analysed the local data by means of the average velocity. A contiguous time series of minute averages above $`25m/sec`$ was classified as free flow, otherwise as congested flow. In Fig. 1 the cross-covariance $`cc(J,\rho )`$ of the flow and the local measured density for different traffic states is also shown. In the free-flow regime the flow is strongly coupled to the density indicating that the average velocity is nearly constant. Also for large densities, in the stop-and-go regime, the flow is mainly controlled by density fluctuations. In the mean density region there is a transition between these two regimes. At cross-covariances in the vicinity of zero the fundamental diagram shows a plateau which supports the interpretation of that synchronised flow leads to $`cc(J,\rho )0`$. In the further comparison of our simulation with the corresponding empirical data we used these traffic states for synchronised flow data and congested states with $`cc(J,\rho )>0.7`$ for stop-and-go data.
For the correct description of the car-car interaction the distance-headway (OV-curve) gives the most important information for the adjustment of the velocities. For densities in the free-flow regime it is obvious that the OV-curve (Figure 2) deviates from the linear velocity-headway curve of the NaSch model. Due to anticipation effects smaller distances occur, so that driving with $`v_{max}`$ is possible even within very small headways. This strong anticipation becomes weaker with increasing density and cars tend to have smaller velocities than the headway allows so that the OV-curve saturates for large distances. At headways of about $`50m`$ the simulation data are in good agreement with the empirical ones, but for large distances the acceleration behaviour of the NaSch model cannot be suppressed so that the velocity increases with the headway. In Fig. 2 the time-headway distributions for different density regimes are shown. The time-headways are calculated via the relation $`\mathrm{\Delta }t=\mathrm{\Delta }x/v`$ with a resolution of $`0.1sec`$. Due to the discrete nature of the model large fluctuations occur. In the free-flow state the anticipation rule is responsible for time-headways smaller than $`1sec`$. The ability to anticipate the predecessors behaviour is getting weaker with increasing density so that small time-headways are nearly vanished in the synchronised and stop-and-go state. Two peaks arise in these states: The peak at a time of $`2sec`$ can be identified with the driver’s efforts for safety: It is recommended to drive with a distance of $`2`$ sec. Nevertheless, with increasing density the NaSch peak at a time of $`1sec`$ (in the NaSch model the minimal time-headway is restricted to $`1sec`$) becomes dominant. The higher the density the stronger this peak structure is pronounced.
## 4 Summary and Discussion
Based on empirical data we tried to find a simple extension of the original NaSch model which is able to reproduce metastable states and synchronised flow as well as microscopic features like density-dependent OV-curves and characteristic time-headways.
First of all, the original NaSch cell length had to be reduced for a more realistic acceleration behaviour which is especially important on highways. For a more realistic car-car interaction anticipation terms seem to play a crucial role. On the one hand, anticipation of the predecessors movement in the next time step allows small time-headways and therefore high flows. On the other hand, braking anticipation by means of brake lights enables a driver to anticipate an imminent velocity reduction due to a jam. It is this braking anticipation which leads to synchronised states and to increased time-headways which results in plateaus in the local fundamental diagram.
Unfortunately, in this short contribution it is not possible to describe all the features of the new model . For example, a finer discretisation of the NaSch model or the corresponding limit $`v_{max}\mathrm{}`$ leads to metastable states analogous to the VDR-model which can be characterised by an order parameter.
For a further validation of this approach it is necessary to extend the model to multi-lane traffic. The implementation of the new model in the Online-Simulation of the Autobahn network in Nordrhein-Westfalen should show its suitability for realistic traffic simulations. |
warning/0001/math0001157.html | ar5iv | text | # Applications of another characterization of 𝛽ℕ∖ℕ
## Introduction
Topological problems that involve the behaviour of families of subsets of the set of natural numbers tend to have (moderately) easy solutions if the Continuum Hypothesis ($`\mathrm{𝖢𝖧}`$) is assumed. The reason for this is that one’s inductions and recursions last only $`\mathrm{}_1`$ steps and that at each intermediate step only countably many previous objects have to be dealt with.
An archetypal example is Parovičenko’s characterization, see , of the space $`\mathrm{}^{}`$ as the only compact zero-dimensional $`F`$-space of weight $`𝔠`$ without isolated points in which non-empty $`G_\delta `$-sets have non-empty interiors. The proof actually shows that $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is the unique atomless Boolean algebra of size $`𝔠`$ with a certain property $`R_\omega `$ and then applies Stone duality to establish uniqueness of $`\mathrm{}^{}`$. It runs as follows: consider two Boolean algebras $`A`$ and $`B`$ with the properties of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ just mentioned and well-order both in type $`\omega _1`$. Assume we have isomorphism $`h`$ between subalgebras $`A_\alpha `$ and $`B_\alpha `$ that contain $`\{a_\beta :\beta <\alpha \}`$ and $`\{b_\beta :\beta <\alpha \}`$ respectively. We need to define $`h(a_\alpha )`$ (if $`a_\alpha A_\alpha `$); consider $`S=\{aA_\alpha :a<a_\alpha \}`$ and $`T=\{aA_\alpha :a<a_\alpha ^{}\}`$. We must find $`bB`$ such that
1. $`h(a)<b`$ if $`aS`$;
2. $`h(a)<b^{}`$ if $`aT`$;
3. $`h(a)b0`$ and $`h(a)b^{}0`$ if $`aA_\alpha `$ but $`aST`$.
Property $`R_\omega `$ says exactly that this is possible — its proper formulation can be found after Definition 2.9.
This paper grew out of observations that in the Cohen model the Boolean algebra $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ retains much of the properties that were used above. In a sense to be made precise later, in Definition 2.3, $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ contains many subalgebras that are like $`A_\alpha `$ and $`B_\alpha `$ above ($`\mathrm{}_0`$-ideal subalgebras); even though these will not be countable the important sets $`S`$ and $`T`$ will be. We also define a cardinal invariant, $`𝔪_c`$, that captures just enough of $`R_\omega `$ to allow a Parovičenko-like characterization of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ in the $`\mathrm{}_2`$-Cohen model — this is Steprāns’ result alluded to in the abstract (Theorem 2.13). During the preparation of this paper we became aware of recent work on the weak Freese-Nation property in . Although the weak Freese-Nation property is stronger than our properties the proofs of the consequences are very similar; therefore we restrict, with few exceptions, ourselves to more topological (and new) applications. Perhaps the difference in approach (weak Freese-Nation versus $`(\mathrm{}_1,\mathrm{}_0)`$-ideal) is mostly a matter of taste but ours arose directly out of Steprāns original results and the essentially folklore facts about the effects of adding Cohen reals.
In Section 2 we shall formulate the properties alluded to above and prove that in the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ does indeed satisfy them. In Section 3 we select some results about $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ (or $`\mathrm{}^{}`$) that are known to hold in the Cohen model and derive them directly from the new properties — whenever we credit a result to some author(s) we mean to credit them with establishing that it holds in the Cohen model. In Sections 4 and 5 we investigate the properties themselves and their behaviour with respect to subalgebras and quotients. Finally, in Section 6 we investigate how much of an important phenomenon regarding $`\mathrm{}^{}`$ persists; we are referring to the fact that under $`\mathrm{𝖢𝖧}`$ for every compact zero-dimensional space $`X`$ of weight $`𝔠`$ or less the Čech-Stone remainder $`(\omega \times X)^{}`$ is homeomorphic to $`\mathrm{}^{}`$.
We would like to take this opportunity to thank the referee for a very insightful remark concerning our version of Bell’s example (see Definition 3.7), which enabled us to simplify the presentation considerably.
## 1. Preliminaries
### Boolean algebras
Our notation is fairly standard: $`b^{}`$ invariably denotes the complement of $`b`$.
For a subset $`S`$ of a Boolean algebra $`B`$, let $`S^{}`$ denote the ideal of members of $`B`$ that are disjoint from every element of $`S`$, i.e., $`S^{}=\{bB:(sS)(bs=0)\}`$. For convenience, we use $`b^{}`$ in place of $`\{b\}^{}`$. Also, let $`b^{}`$ be the principal ideal generated by $`b`$, namely $`\{aB:ab\}`$. Clearly $`b^{}`$ is equal to $`(b^{})^{}`$. Also, for subsets $`S`$ and $`T`$ we let $`ST`$ abbreviate $`(sS)(tT)(st=0)`$; in fact we shall often abbreviate $`st=0`$ by $`st`$.
### Cohen reals
‘The Cohen model’ is any model obtained from a model of the $`\mathrm{𝖦𝖢𝖧}`$ by adding a substantial quantity of Cohen reals — more than $`\mathrm{}_1`$. In particular ‘the $`\mathrm{}_2`$-Cohen model’ is obtained by adding $`\mathrm{}_2`$ many Cohen reals. Actually, since we are intent on proving our results using the *properties* of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ only, many readers may elect to take Lemma 2.2, Theorem 2.7 and the remark made after Proposition 2.12 on faith or else consult for the necessary background on Cohen forcing.
### The weak Freese-Nation property
A partially ordered set $`P`$ is said to have the *weak Freese-Nation property* if there is a function $`F:P[P]^\mathrm{}_0`$ such that whenever $`pq`$ there is $`rF(p)F(q)`$ with $`prq`$.
### Elementary substructures
Consider two structures $`M`$ and $`N`$ (groups, fields, Boolean algebras, models of set theory …), where $`M`$ is a substructure of $`N`$. We say that $`M`$ is an *elementary* substructure of $`N`$, and we write $`MN`$, if every equation, involving the relations and operations of the structures and constants from $`M`$, that has a solution in $`N`$ has a solution in $`M`$ as well.
The Löwenheim-Skolem theorem says that every subset $`A`$ of a structure $`N`$ can be enlarged to an elementary substructure $`M`$ of whose cardinality is the maximum of $`|A|`$ and $`\mathrm{}_0`$. The construction proceeds in the obvious way: in a recursion of length $`\omega `$ one keeps adding solutions to equations that involve ever more constants.
We prefer to think of an argument that uses elementary substructures as the lazy man’s closing off argument; rather than setting up an impressive recursive construction we say “let $`\theta `$ be a suitably large cardinal and let $`M`$ be an elementary substructure of $`H(\theta )`$” and add some words that specify what $`M`$ should certainly contain.
The point is that the impressive recursion is carried out inside $`H(\theta )`$, where $`\theta `$ is ‘suitably large’ (most of the time $`\theta =𝔠^+`$ is a good choice as everything under consideration has cardinality at most $`𝔠`$), and that it (or a nonessential variation) is automatically subsumed when one constructs an elementary substructure of $`H(\theta )`$.
In this paper we shall be working mostly with $`\mathrm{}_1`$-sized elementary substructures, most of which will be *$`\mathrm{}_0`$-covering*. The latter means that every countable subset $`A`$ of $`M`$ is a subset of a countable element $`B`$ of $`M`$. This is not an unreasonable property, considering that the ordinal $`\omega _1`$ has it: every countable subset of $`\omega _1`$ is a subset of a countable ordinal.
An $`\mathrm{}_0`$-covering structure can be constructed in a straightforward way. One recursively constructs a chain $`M_\alpha :\alpha <\omega _1`$ of countable elementary substructures of $`H(\theta )`$ with the property that $`M_\beta :\beta \alpha M_{\alpha +1}`$ for all $`\alpha `$. In the end $`M=_{\alpha <\omega _1}M_\alpha `$ is as required: if $`AM`$ is countable then $`AM_\alpha `$ for some $`\alpha `$ and $`M_\alpha M`$.
For just a few of the results we indicate two proofs: a direct one and one via elementarity — we invite the reader to compare the two approaches and to reflect on their efficacy.
## 2. Two new properties of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$
In this section we introduce two properties that Boolean algebras may have. We shall prove that in the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ has both and that in the $`\mathrm{}_2`$-Cohen model their conjunction actually characterizes $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$.
### $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebras
We begin by defining the $`\mathrm{}_0`$-ideal subalgebras alluded to in the introduction.
###### Definition 2.1.
For a Boolean algebra $`B`$, we will say that a subalgebra $`A`$ of $`B`$ is $`\mathrm{}_0`$-ideal if for each $`bBA`$ the ideal $`\{aA:a<b\}=Ab^{}`$ has a countable cofinal subset.
Of course, by duality, the ideal $`b^{}A`$ is countably generated as well; thus in $`\mathrm{}_0`$-ideal subalgebras the phrase “$`S`$ and $`T`$ are countable” from the introduction is replaced by “$`S`$ and $`T`$ have countable cofinal subsets”.
The main impetus for this definition comes from following result.
###### Lemma 2.2 (\[23, Lemma 2.2\]).
If $`G`$ is $`\mathrm{Fn}(I,2)`$-generic over $`V`$ then $`\mathrm{𝒫}(\mathrm{})V`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`\mathrm{𝒫}(\mathrm{})`$ in $`V[G]`$.
###### Proof.
Let $`\stackrel{~}{X}`$ be an $`\mathrm{Fn}(I,2)`$-name for a subset of $`\mathrm{}`$. It is a well-known fact about $`\mathrm{Fn}(I,2)`$ that there is a countable subset $`J`$ of $`I`$ such that $`\stackrel{~}{X}`$ is completely determined by $`\mathrm{Fn}(J,2)`$. This means that for every $`p\mathrm{Fn}(I,2)`$ and every $`n\mathrm{}`$ we have $`pn\stackrel{~}{X}`$ (or $`pn\stackrel{~}{X}`$) if and only if $`pJ`$ does.
For every $`p\mathrm{Fn}(J,2)`$ define $`X_p=\{n:pn\stackrel{~}{X}\}`$; the countable family of these $`X_p`$ is as required. ∎
The factoring lemma for Cohen forcing (\[18, p. 255\]) implies that for every subset $`J`$ of $`I`$ the subalgebra $`A_J=\mathrm{𝒫}(\mathrm{})V[GJ]`$ is $`\mathrm{}_0`$-ideal in the final $`\mathrm{𝒫}(\mathrm{})`$. Using the fact, seen in the proof above, that names for subsets of $`\mathrm{}`$ are essentially countable one can verify that $`A_{\scriptscriptstyle \mathrm{𝒥}}=_{J\mathrm{𝒥}}A_J`$ for every chain $`\mathrm{𝒥}`$ of subsets of $`I`$ of uncountable cofinality. This shows that in the Cohen model $`\mathrm{𝒫}(\mathrm{})`$ has many $`\mathrm{}_0`$-ideal subalgebras and also that the family of these subalgebras is closed under unions of chains of uncountable cofinality.
What we call $`\mathrm{}_0`$-ideal is called ‘good’ in and in the term $`\sigma `$-subalgebra is used. In the latter paper it is also shown that if $`F:B[B]^\mathrm{}_0`$ witnesses the weak Freese-Nation property of $`B`$ then every subalgebra that is closed under $`F`$ is an $`\mathrm{}_0`$-ideal subalgebra; therefore an algebra with the weak Freese-Nation property has many $`\mathrm{}_0`$-ideal subalgebras and the family of these subalgebras is closed under directed unions.
We are naturally interested in Boolean algebras with many $`\mathrm{}_0`$-ideal subalgebras. Most of our results only require that there are many $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-ideal subalgebras.
###### Definition 2.3.
We will say that a Boolean algebra $`B`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal if the set of $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-ideal subalgebras of $`B`$ contains an $`\mathrm{}_1`$-cub of $`[B]^\mathrm{}_1`$. That is, there is a family $`\mathrm{𝒜}`$ consisting of $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-ideal subalgebras of $`B`$ such that every subset of size $`\mathrm{}_1`$ is contained in some member of $`\mathrm{𝒜}`$ and the union of each chain from $`\mathrm{𝒜}`$ of cofinality $`\omega _1`$ is again in $`\mathrm{𝒜}`$.
We leave to the reader the verification that $`\mathrm{𝒫}(\mathrm{})`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra if and only if $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra (but see Corollary 5.12). It is also worth noting that $`\mathrm{𝒫}(\omega _1)`$ is not $`(\mathrm{}_1,\mathrm{}_0)`$-ideal (see \[14, Proposition 5.3\]).
Since the definition of $`(\mathrm{}_1,\mathrm{}_0)`$-ideal requires that we have some $`\mathrm{}_1`$-cub consisting of $`\mathrm{}_0`$-ideal subalgebras, it is a relatively standard fact that every $`\mathrm{}_0`$-covering elementary substructure of size $`\mathrm{}_1`$ of a suitable $`H(\theta )`$ induces an $`\mathrm{}_0`$-ideal subalgebra. We shall use the following lemma throughout this paper, not always mentioning it explicitly — it is an instance of the rule-of-thumb that says: if $`X,\mathrm{𝒜}M`$, where $`M`$ is suitably closed and $`\mathrm{𝒜}`$ some sort of cub in $`\mathrm{𝒫}(X)`$, then $`XM\mathrm{𝒜}`$.
###### Lemma 2.4.
Let $`B`$ be an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra, let $`\theta `$ be a suitably large cardinal and let $`M`$ be an $`\mathrm{}_0`$-covering elementary substructure of size $`\mathrm{}_1`$ of $`H(\theta )`$ that contains $`B`$. Then $`BM`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B`$.
###### Proof.
Note first that $`M`$ contains an $`\mathrm{}_1`$-cub $`\mathrm{𝒜}`$ as in Definition 2.3: it must contain a solution to the equation
$`x`$ is an $`\mathrm{}_1`$-cub in $`[B]^\mathrm{}_1`$ that consists of $`\mathrm{}_0`$-ideal subalgebras.
Let $`f:\omega _1\mathrm{𝒜}M`$ be a surjection, not necessarily from $`M`$. Because $`M`$ is $`\mathrm{}_0`$-covering we can find, for every $`\alpha \omega _1`$, a countable element $`X_\alpha `$ of $`M`$ that contains $`f[\alpha ]`$. Consider the equation
$$x\mathrm{𝒜}\text{ and }(X_\alpha \mathrm{𝒜})x.$$
This equation has a solution in $`H(\theta )`$ and hence in $`M`$; we may take $`A_\alpha \mathrm{𝒜}M`$ such that $`(X_\alpha \mathrm{𝒜})A_\alpha `$. Thus we construct an increasing chain $`A_\alpha :\alpha <\omega _1`$ in $`\mathrm{𝒜}M`$ that is cofinal in $`\mathrm{𝒜}M`$. It follows that $`(\mathrm{𝒜}M)=_{\alpha <\omega _1}A_\alpha `$ belongs to $`\mathrm{𝒜}`$. Now check carefully that $`BM=(\mathrm{𝒜}M)`$ — use that $`\mathrm{𝒜}`$ is unbounded in $`[B]^\mathrm{}_1`$. ∎
The remarks preceding Definition 2.3 show that an algebra with the weak Freese-Nation property is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal. The converse is almost true — the difference is that we do not require closure under countable unions. In the notation used after Lemma 2.2 the family $`\mathrm{𝒜}=\{A_J:J[I]^\mathrm{}_1\}`$ witnesses that in the Cohen model $`\mathrm{𝒫}(\mathrm{})`$ is always $`(\mathrm{}_1,\mathrm{}_0)`$-ideal. However $`\mathrm{𝒜}`$ is *not* closed under unions of countable chains. Indeed, in one finds the theorem that if $`V`$ satisfies the $`\mathrm{𝖦𝖢𝖧}`$ and the instance $`(\mathrm{}_{\omega +1},\mathrm{}_\omega )(\mathrm{}_1,\mathrm{}_0)`$ of Chang’s conjecture then after adding one dominating real $`d`$ and then $`\mathrm{}_\omega `$ Cohen reals $`\mathrm{𝒫}(\mathrm{})`$ does not have the weak Freese-Nation property — as $`V[d]`$ still satisfies the $`\mathrm{𝖦𝖢𝖧}`$ the final model is a ‘Cohen model’.
Many properties of $`\mathrm{𝒫}(\mathrm{})`$ that hold in the Cohen model can be derived from the weak Freese-Nation property — see for example — and many of these can be derived from the fact that $`\mathrm{𝒫}(\mathrm{})`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal. It is not our intention to duplicate the effort of ; we will concentrate on topological applications. However, to give the flavour, and because we shall use the result a few times, we consider Kunen’s theorem from that in the Cohen model the Boolean algebra $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ does not have a chain of order type $`\omega _2`$.
It is quite straightforward to show that an algebra with the weak Freese-Nation property does not have any well-ordered chains of order type $`\omega _2`$; with a bit more effort the same can be said of $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebras.
###### Proposition 2.5.
An $`(\mathrm{}_1,\mathrm{}_0)`$-ideal Boolean algebra does not have any chains of order type $`\omega _2`$.
###### Proof.
Assume that $`\{c_\alpha :\alpha <\omega _2\}`$ is an increasing chain in $`B`$ and let $`\mathrm{𝒜}`$ be as in Definition 2.3. Recursively construct a chain $`\{A_\alpha :\alpha \omega _1\}`$ in $`\mathrm{𝒜}`$ and an increasing sequence $`\{\gamma _\alpha :\alpha \omega _1\}`$ of ordinals in $`\omega _2`$ as follows. Let $`\gamma _0=0`$ and, given $`A_\beta `$ and $`\gamma _\beta `$ for $`\beta <\alpha `$, let $`A=_{\beta <\alpha }A_\alpha `$ and $`\gamma =sup_{\beta <\alpha }\gamma _\beta `$. Choose $`A_\alpha \mathrm{𝒜}`$ such that $`c_\gamma A_\alpha `$ and such that for every $`aA`$, *if* there is a $`\beta `$ with $`ac_\beta `$ *then* there is a $`\beta `$ such that $`ac_\beta `$ *and* $`c_\beta A_\alpha `$; let $`\gamma _\alpha `$ be the first $`\gamma `$ for which $`c_\gamma A_\alpha `$.
In the end set $`A=\{A_\alpha :\alpha \omega _1\}`$ and $`\lambda =sup\{\gamma _\alpha :\alpha \omega _1\}`$. Now we have a contradiction because although $`c_\lambda ^{}A`$ should be countably generated it is not. Indeed, let $`C`$ be a countable subset of $`c_\lambda ^{}A`$; by construction we have for every $`cC`$$`\beta <\lambda `$ such that $`cc_\beta `$. Let $`\gamma `$ be the supremum of these $`\beta `$’s; then $`\gamma +1<\lambda `$ and so $`c_{\gamma +1}<c_\lambda `$ but no $`cC`$ is above $`c_{\gamma +1}`$. ∎
A proof using elementary substructures runs as follows: let $`MH(\theta )`$ be $`\mathrm{}_0`$-covering and of cardinality $`\mathrm{}_1`$, where $`\theta `$ is suitably large, and assume that $`B`$ and the chain $`\{c_\alpha :\alpha <\omega _2\}`$ belong to $`M`$. Next let $`\delta `$ be the ordinal $`M\omega _2`$; observe that $`\mathrm{cf}\delta =\mathrm{}_1`$: if $`\mathrm{cf}\delta `$ were countable then, because $`M`$ is $`\mathrm{}_0`$-covering, $`\delta `$ would be the supremum of an element of $`M`$ and hence in $`M`$. Consider the element $`c_\delta `$. By Lemma 2.4 there is a countable subset $`T`$ of $`c_\delta ^{}BM`$ that is cofinal in it. There is then (at least) one element $`a`$ of $`T`$ such that $`\{\alpha <\delta :c_\alpha a\}`$ is cofinal in $`\delta `$. However, $`\delta `$ is a solution to
$$x\omega _2\text{ and }a<c_x$$
hence there must be a solution $`\beta `$ in $`M`$ but then $`\beta <\delta `$ and $`a<c_\beta `$, so that $`\{\alpha <\delta :c_\alpha a\}`$ is not cofinal in $`\delta `$.
The reader is invited to supply a proof of the following proposition, which was established in for algebras with the weak Freese-Nation property.
###### Proposition 2.6.
An $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra contains no $`\mathrm{}_2`$-Lusin families.∎
An $`\mathrm{}_2`$-Lusin family is a subset $`A`$ of pairwise disjoint elements with the following property: for every $`x`$ at least one of the sets $`\{aA:ax\}`$ or $`\{aA:ax=0\}`$ has size less than $`\mathrm{}_2`$.
### In the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal
We have already indicated that in the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra. We state it as a separate theorem for future reference.
###### Theorem 2.7.
Let $`V`$ be a model of $`\mathrm{𝖢𝖧}`$ and let $`\kappa `$ be any cardinal. If $`G`$ is generic on $`\mathrm{Fn}(\kappa ,2)`$ then in $`V[G]`$ the algebra $`\mathrm{𝒫}(\mathrm{})`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal.∎
As it is clear that $`\omega _2`$ is the union of an increasing sequence of $`\mathrm{}_1`$-sized subsets the following Proposition, which is Steprāns’ Lemma 2.3, now follows.
###### Proposition 2.8.
In the $`\mathrm{}_2`$-Cohen model there is an increasing sequence of $`\mathrm{}_0`$-ideal subalgebras of $`\mathrm{𝒫}(\mathrm{})`$, each of size $`\mathrm{}_1`$, which is continuous at limits of uncountable cofinality and whose union is all of $`\mathrm{𝒫}(\mathrm{})`$.∎
### Generalizing $`R_\omega `$
The following definition generalizes Parivičenko’s property $`R_\omega `$. After the definition we discuss it more fully and indicate why it is the best possible generalization of $`R_\omega `$. The new property is actually a cardinal invariant which somehow quantifies some, but not all, of the strength of $`\mathrm{𝖬𝖠}_{\mathrm{countable}}`$ — see Proposition 2.12 and Remark 4.6.
###### Definition 2.9.
For a Boolean algebra $`B`$, say that a subset $`A`$ is $`\mathrm{}_0`$-ideal complete, if for any two countable subsets $`S`$ and $`T`$ of $`A`$ with $`ST`$ there is a $`bBA`$ such that $`b^{}A`$ is generated by $`S`$ and $`b^{}A`$ is generated by $`T`$. We will let $`𝔪_c(B)`$ denote the minimum cardinality of a subset of $`B`$ that is not $`\mathrm{}_0`$-ideal complete. Also $`𝔪_c`$ denotes $`𝔪_c\left(\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}\right)`$.
A remark about the previous definition might be in order. In the definition of $`\mathrm{}_0`$-ideal completeness the set $`A`$ is divided into three subsets: $`A_S`$, the set of elements $`a`$ for which there is a finite subset $`F`$ of $`S`$ such that $`aF`$; the set $`A_T`$, defined similarly, and $`A_r`$, the rest of $`A`$. The element $`b`$ must effect the same division of $`A`$: we demand that $`A_S=\{aA:a<b\}`$, $`A_T=\{aA:a<b^{}\}`$ and $`A_r=\{aA:ab`$ and $`ab^{}\}`$. Observe that one can also write $`A_r=\{aA:ba0`$ and $`b^{}a0\}`$; one says that $`b`$ *reaps* the set $`A_r`$. We see that every subset of size less than $`𝔪_c(B)`$ can always be reaped; we shall come back to this in Section 4.
Thus Parovičenko’s property $`R_\omega `$ has become the statement that countable subsets are $`\mathrm{}_0`$-ideal complete, in other words that $`𝔪_c(B)>\mathrm{}_0`$.
###### Remark 2.10.
In Definition 2.9 we explicitly do not exclude the possibility that $`S`$ or $`T`$ is finite or even empty. Thus if $`𝔪_c(B)>\mathrm{}_0`$ then there is no countable strictly increasing sequence with $`1`$ as its supremum: for let $`S`$ be such a sequence and take $`T=\{0\}`$, then there must apparently be a $`b<1`$ such that $`a<b`$ for all $`aS`$.
###### Remark 2.11.
In the case of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ one cannot relax the requirements on $`S`$ and $`T`$: consider a Hausdorff gap; this is a pair of increasing sequences $`\{a_\alpha :\alpha \omega _1\}`$ and $`\{b_\alpha :\alpha \omega _1\}`$ such that $`a_\alpha b_\beta =0`$ for all $`\alpha `$ and $`\beta `$, and for which there is no $`x`$ such that $`a_\alpha x`$ for all $`\alpha `$ and $`b_\beta x^{}`$ for all $`\beta `$. Thus there are cases with $`|S|=|T|=\mathrm{}_1`$ where no $`b`$ can be found.
In an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra with $`𝔪_c(B)>\mathrm{}_0`$ this can be sharpened, as follows. By recursion one can construct a strictly increasing chain $`s_\alpha :\alpha <\delta `$ in $`B`$ with $`0<s_\alpha <1`$ for all $`\alpha `$, until no further choices can be made. Because $`B`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal this must stop before $`\omega _2`$ and because $`𝔪_c(B)>\mathrm{}_0`$ we have $`\mathrm{cf}\delta =\mathrm{}_1`$. Thus we have a situation where no $`b`$ be found with $`|S|=\mathrm{}_1`$ and $`|T|=1`$ (take $`T=\{0\}`$). This shows that if $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal then the cardinal number $`𝔱`$ (see ) is equal to $`\mathrm{}_1`$.
The following proposition shows why we are interested in $`𝔪_c`$.
###### Proposition 2.12.
$`\mathrm{𝖬𝖠}_{\mathrm{countable}}`$ implies that $`𝔪_c=𝔠`$.
###### Proof.
Let $`A`$, $`S`$ and $`T`$ be given, where, without loss of generality, we assume that $`S`$ and $`T`$ are increasing sequences of length $`\omega `$ and $`|A|<𝔠`$. There is a natural countable poset that produces an infinite set $`b`$ such that $`s<b`$ and $`t<b^{}`$ for all $`sS`$ and $`tT`$: it consists of triples $`p,s,t`$, where $`p\mathrm{Fn}(\omega ,2)`$, $`sS`$, $`tT`$ and $`st\mathrm{dom}(p)`$. The ordering is $`p,s,tq,u,v`$ iff $`pq`$, $`su`$, $`tv`$ and if $`n\mathrm{dom}(p)\mathrm{dom}(q)`$ then $`p(n)=1`$ if $`nu`$ and $`p(n)=0`$ if $`nv`$.
It is relatively straightforward to determine a family $`𝒟`$ of fewer than $`𝔠`$ dense sets so that any $`𝒟`$-generic filter produces an element $`b`$ as required. ∎
It is well-known that $`\mathrm{𝖬𝖠}_{\mathrm{countable}}`$ holds in any extension by a ccc finite-support iteration whose length is the final value of the continuum and hence in any model obtained by adding $`𝔠`$ or more Cohen reals.
So in the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra in which $`𝔪_c`$ is $`𝔠`$. Note that this is then consistent with most cardinal arithmetic. However if only $`\mathrm{}_2`$ Cohen reals are added then this provides our characterizations of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ and $`\mathrm{}^{}`$ (see also the results 5.3 through 5.5).
###### Theorem 2.13.
In the $`\mathrm{}_2`$-Cohen model the algebra $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is characterized by the properties of being an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal Boolean algebra of cardinality $`𝔠`$ in which $`𝔪_c`$ has value $`𝔠`$. ∎
The proof is quite straightforward: we use Proposition 2.8 to express any algebra with the properties of the Theorem as the union of a $`\omega _2`$-chain of $`(\mathrm{}_1,\mathrm{}_0)`$-ideal subalgebras and we apply $`𝔪_c=𝔠`$ to construct an isomorphism between it and $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ by recursion. This result and its proof admit a topological reformulation that is quite appealing.
###### Theorem 2.14.
In the $`\mathrm{}_2`$-Cohen model $`\mathrm{}^{}`$ is the unique compact space that is expressible as the limit of an inverse system $`<\{X_\alpha :\alpha <\omega _2\},\{f_\alpha ^\beta :\alpha <\beta <\omega _2\}>`$ such that
1. each $`X_\alpha `$ is a compact zero-dimensional space of weight less than $`𝔠`$;
2. for each limit $`\lambda <\omega _2`$, $`X_\lambda `$ is equal to $`\begin{array}{c}lim\end{array}_{\beta <\lambda }X_\beta `$ and $`f_\alpha ^\lambda =\begin{array}{c}lim\end{array}_{\alpha <\beta <\lambda }f_\alpha ^\beta `$;
3. for each $`\alpha <\beta <\omega _2`$, $`f_\alpha ^\beta `$ sends zero-set subsets of $`X_\beta `$ to zero-sets of $`X_\alpha `$ (i.e. clopen sets are sent to $`G_\delta `$-sets;
4. for each $`\alpha <\omega _2`$ and each pair, $`C_0`$, $`C_1`$ of disjoint cozero-sets of $`X_\alpha `$ (possibly empty), there are a $`\beta <\omega _2`$ and a clopen subset $`b`$ of $`X_\beta `$ such that
$$f_\alpha ^\beta (b)=X_\alpha C_0\text{ and }f_\alpha ^\beta (X_\beta b)=X_\alpha C_1.\mathit{}$$
###### Remark 2.15.
It is our (subjective) feeling that the $`(\mathrm{}_1,\mathrm{}_0)`$-ideal property together with $`𝔪_c`$ captures the essence of the behaviour of $`\mathrm{𝒫}(\mathrm{})`$ and $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ in the Cohen model. By Theorem 2.13 this is certainly the case for the $`\mathrm{}_2`$-Cohen model. Evidence in support of our general feeling will be provided in the next section, where we will derive a number of results from “$`\mathrm{𝒫}(\mathrm{})`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal” that were originally derived in the Cohen model. Apparently it is unknown whether these properties characterize $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ in Cohen models with $`𝔠>\mathrm{}_2`$.
### Other cardinals
We may generalize Definition 2.3 to cardinals other than $`\mathrm{}_1`$: we can call a Boolean algebra $`(\kappa ,\mathrm{}_0)`$-ideal if the family of $`\kappa `$-sized $`\mathrm{}_0`$-ideal subalgebras contains a $`\kappa `$-cub, meaning a subfamily closed under unions of chains of length at most $`\kappa `$ (but of uncountable cofinality). Similarly we can define $`B`$ to be $`(,\mathrm{}_0)`$-ideal if it is $`(\kappa ,\mathrm{}_0)`$-ideal for every (regular) $`\kappa `$ below $`|B|`$.
The discussion after Lemma 2.2 establishes that every Boolean algebra with the weak Freese-Nation property $`(,\mathrm{}_0)`$-ideal and in any Cohen model the algebra $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(,\mathrm{}_0)`$-ideal. One can also prove a suitable version of Lemma 2.4.
###### Lemma 2.16.
Let $`B`$ be an $`(\kappa ,\mathrm{}_0)`$-ideal algebra, let $`\theta `$ be a suitably large cardinal and let $`M`$ be an elementary substructure of size $`\kappa `$ of $`H(\theta )`$ that contains $`B`$. Then $`BM`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B`$, *provided* $`M`$ can be written as $`_{\alpha <\kappa }M_\alpha `$, where $`M_\beta :\beta \alpha M_{\alpha +1}`$ for all $`\alpha `$.∎
In applications one also needs $`M`$ to be $`\mathrm{}_0`$-covering; this is possible only if the structure $`([\kappa ]^\mathrm{}_0,)`$ has cofinality $`\kappa `$. This accounts for the assumption $`\mathrm{cf}[\kappa ]^\mathrm{}_0=\kappa `$ in Theorem 5.6.
## 3. The Axiom “$`\mathrm{𝒫}(\mathrm{})`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra”
Throughout this section we assume that $`\mathrm{𝒫}(\mathrm{})`$ is an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra and show how useful this can be as an axiom in itself. We fix an $`\mathrm{}_1`$-cub $`\mathrm{𝒜}`$ in $`\left[\mathrm{𝒫}(\mathrm{})\right]^\mathrm{}_1`$ that consists of $`\mathrm{}_0`$-ideal subalgebras.
### Mappings onto cubes
In order to avoid additional definitions we state, in the rest of this section, some of the results in their topological, rather than Boolean algebraic, formulations. We shall also use elementary substructures to our advantage; we shall use the phrase ‘by elementarity’ to indicate that a judicious choice of equation would give the desired result.
The first result we present is due to Baumgartner and Weese .
###### Theorem 3.1.
If $`X`$ is a compact space with a countable dense set $`D`$ such that every infinite subset of $`D`$ contains a converging subsequence, then $`X`$ does not map onto $`[0,1]^{\omega _2}`$.
###### Proof.
If $`f`$ were a mapping of $`X`$ onto $`[0,1]^{\omega _2}`$ then $`f[D]`$ would be a countable dense subset of $`[0,1]^{\omega _2}`$ with the same property as $`D`$. Therefore we are done once we show that $`[0,1]^{\omega _2}`$ has no countable dense subset every infinite subset of which contains a converging sequence. So we take a countable dense subset of $`[0,1]^{\omega _2}`$, which we identify with $`\mathrm{}`$, and exhibit an infinite subset of it that does not contain a converging sequence.
To this end we fix a suitably large cardinal $`\theta `$ and consider an $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-covering elementary substructure $`M`$ of $`H(\theta )`$. We put $`\delta =M\omega _2`$ and let $`c=\mathrm{}\pi _\delta ^{}\left[[\frac{1}{4},1]\right]`$ and $`d=\mathrm{}\pi _\delta ^{}\left[[0,\frac{3}{4}]\right]`$, where, generally, $`\pi _\alpha `$ denotes the projection onto the $`\alpha `$-th coordinate.
Let $`CM`$ be a countable set such that $`c^{}M`$ is generated by $`C_1=c^{}C`$; similarly choose a countable element $`D`$ of $`M`$ for $`d`$ and put $`D_1=d^{}D`$. This can be done because $`M`$ is $`\mathrm{}_0`$-covering.
For $`xC`$ let $`S_x=\{\alpha :x^{}\pi _\alpha ^{}\left[[\frac{1}{4},1]\right]\}`$. Observe that if $`\delta S_x`$ then $`S_x`$ is cofinal in $`\omega _2`$ because, apparently, there is then in $`M`$ no solution $`\eta `$ to $`(\eta \omega _2)(\alpha S_x)(\alpha <\eta )`$ in $`M`$ and hence not in $`H(\theta )`$ either. It follows that $`C_1`$ is contained in the set
$$C_2=\{xC:S_x\text{ is cofinal in }\omega _2\},$$
which, by elementarity, is in $`M`$. We define $`T_y`$, for $`yD`$, in an analogous way and find the set
$$D_2=\{yD:T_y\text{ is cofinal in }\omega _2\},$$
which is in $`M`$ and which contains $`D_1`$. We claim that the ideal generated by $`C_2D_2`$ does not contain a cofinite subset of $`\mathrm{}`$.
Indeed, take $`x_1`$, …, $`x_k`$ in $`C_2`$ and $`y_1`$, …, $`y_k`$ in $`D_2`$. We can find distinct $`\alpha _1`$, …, $`\alpha _k`$, $`\beta _1`$, …, $`\beta _k`$ with $`\alpha _iS_{x_i}`$ and $`\beta _iT_{y_i}`$ for all $`i`$. The set $`U=_{i=1}^k\left(\pi _{\alpha _i}^{}\left[[0,\frac{1}{4})\right]\pi _{\beta _i}^{}\left[(\frac{3}{4},1]\right]\right)`$ is disjoint from $`_{i=1}^k(x_iy_i)`$ and its intersection with $`\mathrm{}`$ is infinite.
Because $`C_2D_2`$ is countable we can, by elementarity, find an infinite subset $`a`$ of $`\mathrm{}`$ in $`M`$ that is almost disjoint from every one of its elements. Now if $`a`$ had an infinite converging subset then, again by elementarity, it would have one, $`b`$ say, that belongs to $`M`$. However, if $`b^{}c`$ then $`b^{}x`$ for some $`xC_1`$, which is impossible; likewise $`b^{}d`$ is impossible. It follows that $`\pi _\delta [b]`$ does not converge in $`[0,1]`$. ∎
###### Remark 3.2.
A careful study of the proof of Theorem 3.1 shows how one can reach $`\delta `$ by a traditional recursion. Build an increasing sequence $`A_\alpha :\alpha <\omega _1`$ in $`\mathrm{𝒜}`$ and a sequence $`\delta _\alpha :\alpha <\omega _1`$ in $`\omega _2`$ by doing the following at successor steps. Enumerate $`A_\alpha `$ as $`a_\beta :\beta <\omega _1`$ and choose, whenever possible, a subset $`b_\beta `$ of $`a_\beta `$ that converges in $`[0,1]^\mathrm{}_2`$. Next choose, for each $`\beta <\omega _1`$, a subset $`d_\beta `$ as follows: let $`C=\{\gamma <\beta :S_{a_\gamma }`$ is cofinal in $`\omega _2\}`$ and $`D=\{\gamma <\beta :T_{a_\gamma }`$ is cofinal in $`\omega _2\}`$ (here $`S_x`$ and $`T_x`$ are defined as in the proof); as in the proof we can find a nonzero $`d_\beta `$ in $`(CD)^{}`$. Let $`A_{\alpha +1}`$ be an element of $`\mathrm{𝒜}`$ that contains $`A_\alpha \{b_\beta \}_{\beta \omega _1}\{d_\beta \}_{\beta \omega _1}`$ and choose $`\delta _{\alpha +1}`$ so large that $`supS_a<\delta _{\alpha +1}`$ or $`supT_a<\delta _{\alpha +1}`$ whenever $`aA_{\alpha +1}`$ and $`S_a`$ or $`T_a`$ is bounded in $`\omega _2`$. The rest of the proof is essentially the same.
The next result, from , provides a nice companion to Theorem 3.1.
We prove the result for the case $`𝔠=\mathrm{}_2`$ only — basically the same proof will work when $`𝔠=\mathrm{}_n`$ for some $`n\omega `$. For larger values of $`𝔠`$ we need assumptions like to push the argument through.
###### Theorem 3.3 ($`2^\mathrm{}_1=𝔠=\mathrm{}_2`$).
If $`X`$ is compact, separable and of cardinality greater than $`𝔠`$ then $`X`$ maps onto $`I^𝔠`$.
###### Proof.
Suppose that $`X`$ is compact and that $`\mathrm{}`$ is dense in $`X`$. Fix a suitably large cardinal $`\theta `$ and construct an increasing sequence $`M_\alpha :\alpha <\omega _2`$ of $`\mathrm{}_1`$-sized elementary substructures of $`H(\theta )`$ that are $`\mathrm{}_0`$-covering and that always $`M_\beta :\beta <\alpha M_{\alpha +1}`$; put $`M=_{\alpha <\omega _2}M_\alpha `$. Furthermore by the cardinality assumptions we can ensure that $`M^\omega `$ and $`M^{\omega _1}`$ are subsets of $`M`$.
Fix any point $`x`$ in $`XM`$ (because $`|M|<|X|`$). For each $`\alpha <\omega _2`$ let $`\mathrm{}_\alpha =\{F\mathrm{}:FM_\alpha `$ and $`x\mathrm{cl}F\}`$. Because $`|\mathrm{}_\alpha |<𝔠`$ we have $`\mathrm{}_\alpha M`$ and so by elementarity there is a point $`x_\alpha XM`$ such that $`\mathrm{}_\alpha =\{F\mathrm{}:FM_\alpha `$ and $`x_\alpha \mathrm{cl}F\}`$. Fix a function $`f_\alpha :X[0,1]`$ so that $`f_\alpha (x_\alpha )=0`$ and $`f_\alpha (x)=1`$, and set $`a_\alpha =\{n:f_\alpha (n)<\frac{1}{4}\}`$ and $`b_\alpha =\{n:f_\alpha (n)>\frac{3}{4}\}`$. There is a $`g(\alpha )<\omega _2`$ such that $`x_\alpha `$, $`f_\alpha `$, $`a_\alpha `$ and $`b_\alpha `$ belong to $`M_{g(\alpha )}`$. Finally, fix a cub $`C`$ in $`\omega _2`$ such that $`\alpha <\lambda `$ implies $`g(\alpha )<\lambda `$ whenever $`\lambda C`$. Set $`S=\{\lambda C:\mathrm{cf}\lambda =\omega _1\}`$.
Now apply the Pressing-Down lemma to find a stationary set $`TS`$ and a $`\beta \omega _2`$ so that, for every $`\lambda T`$, each of $`a_\lambda ^{}M_\lambda `$, $`a_\lambda ^{}M_\lambda `$, $`b_\lambda ^{}M_\lambda `$ and, $`b_\lambda ^{}M_\lambda `$ is generated by a countable subset of $`M_\beta `$.
By induction on $`\lambda T`$ we prove that
$$\{(a_\alpha ,b_\alpha ):\alpha T\lambda +1\}$$
is a dyadic family. In fact, if $`H`$ and $`K`$ are disjoint finite subsets of $`T`$ then
$$\underset{\alpha H}{}a_\alpha \underset{\alpha K}{}b_\alpha $$
is not in the ideal generated by $`\mathrm{}_\beta `$. Let $`\lambda =\mathrm{max}(HK)`$ and suppose first that $`\lambda K`$. Put $`y=_{\alpha H}a_\alpha _{\alpha K\{\lambda \}}b_\alpha `$; then $`y`$ is not contained in any member of $`\mathrm{}_\beta `$.
Assume there is an $`I\mathrm{}_\beta `$ such that $`yb_\lambda I`$; then $`yI`$ belongs to $`M_\lambda b_\lambda ^{}`$ and hence it is contained in a $`cM_\beta b_\lambda ^{}`$. Because $`cb_\lambda `$ we have $`f_\lambda [c][0,\frac{3}{4}]`$ and so $`x\mathrm{cl}c`$ whence $`c\mathrm{}_\beta `$. We have a contradiction since it now follows that $`yIc\mathrm{}_\beta `$.
Next suppose $`\lambda H`$ and put $`y=_{\alpha H\{\lambda \}}a_\alpha _{\alpha K}b_\alpha `$; again, $`y`$ is not contained in any element of $`\mathrm{}_\beta `$. Assume there is $`I\mathrm{}_\beta `$ such that $`ya_\lambda I`$; now $`yI`$ belongs to $`M_\lambda a_\lambda ^{}`$ and hence it is contained in a $`cM_\beta a_\lambda ^{}`$. Because $`ca_\lambda `$ we have $`f_\lambda [c][\frac{1}{4},1]`$ and so $`x_\lambda \mathrm{cl}c`$; because $`cM_\lambda `$ this means $`x\mathrm{cl}c`$ whence $`c\mathrm{}_\beta `$. Again we have a contradiction because we have $`yIc\mathrm{}_\beta `$.
It now follows that $`\mathrm{}\{f_\lambda :\lambda T\}`$ is a continuous map from $`X`$ into $`I^T`$ and that the image of $`X`$ contains $`\{0,1\}^T`$, which in turn can be mapped onto $`[0,1]^T`$. ∎
This result is optimal: in Fedorčuk constructed, in the $`\mathrm{}_2`$-Cohen model, a separable compact space of cardinality $`𝔠=2^\mathrm{}_1`$ that does not map onto $`I^𝔠`$ because its weight is $`\mathrm{}_1`$.
### The size of sequentially compact spaces
The next result arose in the study of compact sequentially compact spaces (see for the applications). Recall that a filter(base) of sets in a space $`X`$ is said to *converge* to a point if every neighbourhood of the point contains an element of the filter(base).
###### Lemma 3.4.
If $`X`$ is a regular space and $`\mathrm{}X`$ has the property that every infinite subset contains a converging sequence then for each ultrafilter $`u`$ on $`\mathrm{}`$ that converges to some point of $`X`$ there is an $`\mathrm{}_1`$-sized filter subbase, $`v`$, that converges (to the same point).
###### Proof.
Let $`u`$ be an ultrafilter on $`\mathrm{}`$ that converges to a point $`x`$ of $`X`$. Let $`MH(\theta )`$ be any $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-covering model such that $`u`$, $`X`$ and $`x`$ are in $`M`$. We shall prove that $`v=Mu`$ also converges to $`x`$.
Since $`M`$ is $`\mathrm{}_0`$-covering and $`uM`$, there is an increasing chain $`\{u_\alpha :\alpha \omega _1\}`$ of countable subsets of $`u`$ such that each $`u_\alpha `$ is a member of $`M`$ and $`uM=\{u_\alpha :\alpha \omega _1\}`$. For each $`\alpha \omega _1`$, there is an $`a_\alpha \mathrm{}`$ such that $`a_\alpha M`$ and $`a_\alpha U`$ is finite for each $`Uu_\alpha `$. By the assumption on the embedding of $`\mathrm{}`$ in $`X`$, we may assume that $`a_\alpha `$ converges to a point $`x_\alpha X`$. Observe that for each $`b\mathrm{𝒫}(\mathrm{})M`$ we have $`bu`$ if and only if $`a_\alpha ^{}b`$ for uncountably many $`\alpha `$.
Suppose that $`x`$ is an element of some open subset $`W`$ of $`X`$. Let $`\{b_n:n\omega \}M\mathrm{𝒫}(\mathrm{})`$ generate $`(W\mathrm{})^{}M`$. Since $`u`$ converges to $`x`$, the set $`W\mathrm{}`$ is a member of $`u`$. Therefore $`\mathrm{}b_n`$ is a member of $`u`$ for each $`n`$, hence there is an $`\alpha `$ such that $`\{\mathrm{}b_n:n\omega \}u_\alpha `$. It follows, then, that $`a_\beta `$ is almost disjoint from each $`b_n`$ for all $`\beta \alpha `$. Thus, for $`\alpha \beta <\omega _1`$ we have $`a_\beta (W\mathrm{})^{}`$, which means that $`Wa_\beta `$ is infinite for each $`\beta \alpha `$. It follows that $`\{x_\beta :\alpha \beta <\omega _1\}`$ is contained in the closure of $`W`$. Since $`W`$ was an arbitrary neighbourhood of $`x`$ and $`X`$ is regular, it follows that there is an $`\alpha ^{}`$ such that $`\{x_\beta :\alpha ^{}\beta <\omega _1\}`$ is contained in $`W`$. Since $`W`$ is open, it follows that $`a_\beta `$ is almost contained in $`W`$ whenever $`\alpha ^{}\beta <\omega _1`$.
Now suppose that $`\{c_n:n\omega \}M\mathrm{𝒫}(\mathrm{})`$ generates $`(W\mathrm{})^{}M`$. By the above, it follows that, whenever $`\alpha ^{}\beta <\omega _1`$, there is an $`n`$ such that $`a_\beta `$ is almost contained in $`c_n`$. Fix $`n`$ such that $`a_\beta `$ is almost contained in $`c_n`$ for uncountably many $`\eta `$. As we observed above, it follows that $`c_nu`$. Therefore, as required, we have shown that $`W`$ contains a member of $`v`$. ∎
###### Theorem 3.5.
If $`X`$ is a regular space in which $`\mathrm{}`$ is dense and every subset of $`\mathrm{}`$ contains a converging sequence, then $`X`$ has cardinality at most $`2^\mathrm{}_1`$.
###### Proof.
Each point of $`X`$ will be the unique limit point of some filter base on $`\mathrm{}`$ of cardinality $`\mathrm{}_1`$. ∎
Compare this theorem with Theorem 3.1, which draws the conclusion that $`X`$ cannot be mapped onto $`[0,1]^\mathrm{}_2`$. In fact if $`2^\mathrm{}_1<2^\mathrm{}_2`$ then Theorem 3.1 becomes a consequence of Theorem 3.5.
### $`\mathrm{}^{}`$ minus a point
It was shown by Gillman in , assuming $`\mathrm{𝖢𝖧}`$, that for every point $`u`$ of $`\mathrm{}^{}`$ one can partition in $`\mathrm{}^{}\{u\}`$ into two open sets, each of which has $`u`$ in its closure. Clearly this show that $`\mathrm{}^{}\{u\}`$ is not $`C^{}`$-embedded in $`\mathrm{}^{}`$. Here we present Malykhin’s result, from , that establishes the complete opposite.
###### Theorem 3.6 ($`𝔪_c>\mathrm{}_1`$).
$`\mathrm{}^{}`$ minus a point is $`C^{}`$-embedded in $`\mathrm{}^{}`$.
###### Proof.
Assume that $`\mathrm{}^{}\{u\}`$ is not $`C^{}`$-embedded; so there is a continuous function $`f:\mathrm{}^{}\{u\}[0,1]`$ such that $`u`$ is simultaneously a limit point of $`f^{}(0)`$ and $`f^{}(1)`$.
Fix an increasing sequence $`\{c_n:n\omega \}`$ in $`\mathrm{𝒫}(\mathrm{})u`$ such that in the case that $`u`$ is not a $`P`$-point every member of $`u`$ meets some $`c_n`$ in an infinite set. Now define $`\mathrm{}=\{a\{c_n\}_n^{}:a^{}f^{}(0)\}`$ and $`\mathrm{𝒥}=\{a\{c_n\}_n^{}:a^{}f^{}(1)\}`$. The ideals $`\mathrm{}`$ and $`\mathrm{𝒥}`$ are $`P`$-ideals: if $`I`$ is a countable subset of $`\mathrm{}`$ then apply $`𝔪_c>\mathrm{}_0`$ to find $`a\{c_n\}_n^{}`$ with $`Ia^{}`$. Because $`au`$ the function $`f`$ is defined on all of $`a^{}`$; it then follows that there is a $`ba`$ such that $`Ib^{}`$ and $`b^{}f^{}(0)`$.
Claim 1. If $`Uu`$ then there is $`a\mathrm{}`$ such that $`aU`$ (similarly there is $`b\mathrm{𝒥}`$ with $`bU`$).
Proof. For every $`n`$ the set $`a_n=Uc_n`$ belongs to $`u`$, hence $`a_n^{}`$ meets $`f^{}(0)`$ and there is a subset $`b_n`$ of $`a_n`$ with $`b_n^{}f^{}(0)`$ — here we use the well-known fact that $`f^{}(0)`$ is regularly closed. Now take an infinite set $`a`$ such that $`b^{}_{mn}a_m`$ for all $`n`$; then $`aU`$ and $`a^{}f^{}(0)`$.
Let $`M`$ be an $`\mathrm{}_0`$-covering elementary substructure of $`H(\theta )`$, of size $`\mathrm{}_1`$, that contains $`u`$, $`f`$ and $`\{c_n:n\omega \}`$.
Claim 2. If $`b\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is such that $`\mathrm{}Mb^{}`$ (or $`\mathrm{𝒥}Mb^{}`$) then there is $`UuM`$ such that $`U^{}b`$.
Proof. Let $`Cb^{}M`$ be a countable cofinal set and choose for every $`cC`$, whenever possible, an $`i_c\mathrm{}M`$ such that $`ci_c`$. Let $`IM`$ be a countable subset of $`\mathrm{}`$ that contains all the possible $`i_c`$; because $`I`$ is countable there is $`i\mathrm{}M`$ such that $`a<i`$ for all $`aI`$. Note that $`ib^{}M`$, hence there is $`cC`$ such that $`i<c`$; it follows that $`\mathrm{}Mc^{}`$. Note that in $`M`$ there is no solution to “$`x\mathrm{}`$ and $`xc`$” hence there is none in $`H(\theta )`$; it follows that $`\mathrm{}c^{}`$. But this implies that $`cu`$.
The claim implies that if $`b\mathrm{𝒫}(\mathrm{})`$ meets every $`UuM`$ in an infinite set then there are $`I\mathrm{}M`$ and $`J\mathrm{𝒥}M`$ that meet $`b`$ in an infinite set. This in turn implies that the closed set $`F=\{U^{}:UuM\}`$ is contained in $`\mathrm{cl}f^{}(0)\mathrm{cl}f^{}(1)`$. The inequality $`𝔪_c>\mathrm{}_1`$ implies that $`uM`$ does not generate an ultrafilter, so that $`F`$ consists of more than one point. This contradicts our assumption that $`u`$ is the only point in $`\mathrm{cl}f^{}(0)\mathrm{cl}f^{}(1)`$. ∎
### A first-countable nonremainder
The final result in this section is due to M. Bell . He produced a compact first countable space which is not a continuous image of $`\mathrm{}^{}`$ (equivalently: not a remainder of $`\mathrm{}`$). We will show that such a space can be taken to be a subspace of the following space, which is an image of $`\mathrm{}^{}`$. The space is, in hindsight, easy to describe. In the first version of this paper we started out with a generalization of Alexandroff’s doubling procedure; the referee rightfully pointed out that we were simply working with the square of the Alexandroff double of the Cantor set. In private correspondence, Bell points out that his original space is not embeddable in the square of the Alexandroff double.
###### Definition 3.7.
Let $`\mathrm{𝔻}`$ be the Alexandroff double of the Cantor set, i.e., $`\mathrm{𝔻}=\mathrm{}\times 2`$, toplogized as follows: all points of $`\mathrm{}\times \{1\}`$ are isolated and basic neighbourhoods of a point $`x,0`$ is of the form $`(U\times 2)\left\{x,1\right\}`$, where $`U`$ is a neighbourhood of $`x`$ in $`\mathrm{}`$. It is well-known that this results in a compact first countable space.
We let $`\mathrm{𝕂}=\mathrm{𝔻}\times \mathrm{𝔻}`$. We shall show that $`\mathrm{𝕂}`$ is a continuous image of $`\mathrm{}^{}`$ and that it contains a closed subspace that is *not* a continuous image of $`\mathrm{}^{}`$.
In proving that $`\mathrm{𝕂}`$ is a continuous image of $`\mathrm{}^{}`$ we use results from . We let $`W=\{k,l\mathrm{}^2:l2^k\}`$ and we let $`\pi :W\mathrm{}`$ be the projection on the first coordinate. A compact space is called an *orthogonal* image of $`\mathrm{}^{}`$ if there is a continuous map $`f:W^{}X`$ such that the diagonal map $`\beta \pi f:W^{}\mathrm{}^{}\times X`$ is onto. Theorem 2.5 of states that products of $`𝔠`$ (or fewer) orthogonal images of $`\mathrm{}^{}`$ are again orthogonal images of $`\mathrm{}^{}`$. Thus the following proposition more than shows that $`\mathrm{𝕂}`$ is a continuous image of $`\mathrm{}^{}`$.
###### Proposition 3.8.
The space $`\mathrm{𝔻}`$ is an orthogonal image of $`\mathrm{}^{}`$.
###### Proof.
Let $`\{q_l:l\mathrm{}\}`$ be a countable dense subset of $`\mathrm{}`$ and define $`f:W\mathrm{}`$ by $`f(k,l)=q_l`$; observe that $`\beta f`$ maps $`W^{}`$ onto $`\mathrm{}`$ and that $`\beta f(u)=q_l`$ for all $`u`$ in $`\{k,l:2^kl\}^{}`$. This readily implies that $`\beta \pi \beta f`$ maps $`W^{}`$ onto $`\mathrm{}`$.
A minor modification of the usual argument that nonempty $`G_\delta `$-subsets of $`W^{}`$ have nonempty interior lets us associate with every $`x\mathrm{}`$ a subset $`A_x`$ of $`W`$ that meets all but finitely many of the vertical lines $`V_k=\{k,l:l2^k\}`$ and such that $`\beta f[A_x^{}]=\{x\}`$. Now define $`g:WW`$ by $`g(k+1,2l)=g(k+1,2l+1)=k,l`$ (and $`g(0,0)=0,0`$); observe that $`B_x=g^{}[A_x]`$ meets all but finitely many $`V_k`$ in at least two points, so that we may split it into two parts, $`C_x`$ and $`D_x`$, each of which meets all but finitely many $`V_k`$.
Now we turn the map $`\beta f\beta g:W^{}\mathrm{}`$ into a map from $`W^{}`$ to $`\mathrm{𝔻}`$: every point of $`D_x^{}`$ will be mapped to $`x,1`$ and the points $`u`$ of $`W^{}_xD_x^{}`$ will be mapped to $`<(\beta f\beta g)(u),0>`$. It is straightforward to verify that the map $`h`$ thus obtained witnesses that $`\mathrm{𝔻}`$ is an orthogonal image of $`\mathrm{}^{}`$. ∎
###### Theorem 3.9 ($`2^\mathrm{}_1=𝔠`$).
The space $`\mathrm{𝕂}`$ has a compact subspace $`X`$ that is not an image of $`\mathrm{}^{}`$.
###### Proof.
We obtain $`X`$ by removing a (suitably chosen) set of isolated points from $`\mathrm{𝕂}`$. We enumerate $`\mathrm{}`$ as $`\{r_\alpha :\alpha <𝔠\}`$ and we use our assumption $`2^\mathrm{}_1=𝔠`$ to enumerate the family $`[\omega _1]^\mathrm{}_1`$ as $`\{A_\alpha :\alpha 𝔠\}`$ with cofinal repetitions. The set of isolated points that we keep is $`\{<r_\alpha ,1,r_\beta ,1>:\alpha A_\beta `$ or $`\beta A_\alpha \}`$. Furthermore, we let $`U_\alpha `$ be intersection of $`X`$ with the clopen ‘cross’
$$\mathrm{𝔻}\times \left\{r_\alpha ,1\right\}\left\{r_\alpha ,1\right\}\times \mathrm{𝔻}.$$
Note that then for all $`\alpha `$ one has $`A_\alpha =\{\xi \omega _1:U_\xi U_\alpha =\mathrm{}\}`$.
Now suppose that $`f`$ is a mapping of $`\mathrm{}^{}`$ onto $`X`$ and for each $`\alpha 𝔠`$ fix a representative $`a_\alpha \mathrm{}`$ for $`f^{}[U_\alpha ]`$. Observe that thus $`A_\alpha =\{\xi \omega _1:a_\xi a_\alpha =^{}\mathrm{}\}`$. Fix an $`\mathrm{}_1`$-sized $`\mathrm{}_0`$-ideal subalgebra $`B`$ of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ that contains $`\{a_\xi :\xi \omega _1\}`$.
For each $`bB`$ let $`S_b=\{\xi :a_\xi <b\}`$ and pick $`\alpha 𝔠`$ such that both $`S_bA_\alpha `$ and $`S_bA_\alpha `$ have cardinality $`\mathrm{}_1`$ whenever $`S_b`$ has cardinality $`\mathrm{}_1`$. Now $`Ba_\alpha ^{}`$ is countably generated and it contains the uncountable set $`\{a_\xi :\xi A_\alpha \}`$; it follows that there is a $`b<a_\alpha ^{}`$ such that $`S_b`$ is uncountable. But now pick any $`\xi S_bA_\alpha `$. Then $`a_\xi ^{}b^{}a_\alpha ^{}`$ yet $`a_\xi a_\alpha ^{}\mathrm{}`$ — a clear contradiction. ∎
## 4. Other cardinal invariants
In this section we relate $`𝔪_c(B)`$ to other known cardinal invariants of Boolean Algebras; we have already connected $`𝔪_c`$ to the idea of reaping. We formalize this idea in the following definition, which is analogous to the cardinal $`𝔯`$ in $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ (see or ).
###### Definition 4.1.
A subset $`A`$ of a Boolean algebra $`B`$ is *reaped* by the element $`bB`$, if $`b`$ and its complement meet every non-zero element of $`A`$. The cardinal $`𝔯(B)`$ is defined as the minimum cardinality of a subset $`A`$ of $`B`$ that is not reaped by any element of $`B`$.
Our discussion after Definition 2.9 therefore establishes the inequality $`𝔯(B)𝔪_c(B)`$. The other half, so to speak, of $`𝔪_c`$ is provided by the proper analogue, for arbitrary Boolean algebras, of the cardinal number $`𝔡`$.
In Van Douwen showed that $`𝔡`$ is equal to the number $`𝔡_2`$ from the following definition.
###### Definition 4.2.
If $`D\omega ^\omega `$ and $`A[\omega ]^\mathrm{}_0`$ then $`D`$ is said to *dominate on $`A`$* if for each $`g\omega ^\omega `$ there are $`dD`$ and $`aA`$ such that $`g(n)<d(n)`$ for each $`na`$. The cardinal $`𝔡_2`$ is defined as
$$𝔡_2=\mathrm{min}\left\{|A|+|D|:A[\omega ]^\mathrm{}_0\text{}D\omega ^\omega \text{ and }D\text{ dominates on }A\right\}.$$
To find a natural analogue of the cardinal invariant $`𝔡_2`$ in a general Boolean algebra we proceed along the lines of Rothberger’s work on the cardinals $`𝔟`$ and $`𝔡`$. For this we say that an ideal $`\mathrm{}`$ in a Boolean algebra $`B`$ is *co-generated* by a set $`S`$ if $`\mathrm{}=S^{}`$. We will say that $`\mathrm{}`$ is countably co-generated if there is a countably infinite set that co-generates it but, in order to avoid cumbersome consideration of cases, no finite set co-generates it. The cardinal invariant $`𝔡`$ is naturally equal to the minimum cardinal of a cofinal subset of any countably co-generated non-principal ideal in $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ and likewise $`𝔟`$ is the cardinal of an unbounded subset of a countably co-generated non-principal ideal in $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ — this is so because any countably co-generated ideal in $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is naturally isomorphic to the ideal in $`\mathrm{𝒫}(\omega \times \omega )/\mathrm{𝑓𝑖𝑛}`$ that is generated by the set of the form
$$L_f=\{m,n:m\omega ,nf(m)\},$$
where $`f\omega ^\omega `$.
###### Proposition 4.3.
If $`\mathrm{}\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is a countably co-generated ideal then $`𝔡_2`$ is equal to the minimum cardinality of a subset $`\mathrm{𝒥}`$ of $`\mathrm{}`$ such that no member of $`\mathrm{}`$ meets every member of $`\mathrm{𝒥}`$.
###### Proof.
Given $`A`$ and $`D`$ let $`\mathrm{𝒥}`$ be the set of graphs $`ga`$, where $`gD`$ and $`aA`$. Observe that $`ga`$ is disjoint from $`L_f`$ iff $`f(n)<g(n)`$ for all $`na`$.
Conversely, given $`\mathrm{𝒥}`$ construct for each $`J\mathrm{𝒥}`$ a function $`g_J`$ with domain $`a_J=\{m:(n)(m,nJ)\}`$ whose graph is contained in $`J`$. Clearly if $`L_fJ=\mathrm{}`$ then $`g_J(n)>f(n)`$ for $`na_J`$. ∎
It is clear that this proof uses a lot of the underlying structure of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$; one cannot hope to do something similar for arbitrary Boolean algebras.
###### Definition 4.4.
For a Boolean algebra $`B`$, let $`𝔡_2(B)`$ be the minimum cardinal $`\kappa `$ such that whenever $`\mathrm{}`$ is a countably co-generated ideal of $`B`$ and $`A`$ is a subset of $`\mathrm{}`$ of cardinality less than $`\kappa `$, there is a $`b\mathrm{}`$ that meets each member of $`A`$.
We get the following characterization of $`𝔪_c(B)`$, with the immediate corollary that $`𝔪_c=\mathrm{min}\{𝔯,𝔡\}`$.
###### Theorem 4.5.
If $`B`$ is a Boolean algebra with $`𝔪_c(B)>\mathrm{}_0`$ then $`𝔪_c(B)`$ is the minimum of $`𝔯(B)`$ and $`𝔡_2(B)`$.
###### Proof.
We have already seen that $`𝔯(B)𝔪_c(B)`$.
Let $`\mathrm{}`$ be co-generated by the countable set $`S`$ and let $`\mathrm{𝒥}`$ be a subset of $`\mathrm{}`$ of cardinality less than $`𝔪_c(B)`$. By our assumption on $`𝔪_c(B)`$ the set $`A=SJ`$ has cardinality less than $`𝔪_c(B)`$ as well so that it is $`\mathrm{}_0`$-ideal complete. There is therefore an element $`b`$ of $`B`$ such that $`s<b`$ for all $`sS`$ and such that $`b`$ reaps $`\mathrm{𝒥}`$; then $`b^{}`$ is an element of $`\mathrm{}`$ that meets all elements of $`\mathrm{𝒥}`$. Thus $`𝔡_2(B)𝔪_c(B)`$.
Next let $`A`$ be a subset of $`B`$, of size less than both $`𝔯(B)`$ and $`𝔡_2(B)`$. Let $`S`$ and $`T`$ be countable subsets of $`A`$ with $`ST`$ and divide $`A`$ into three subsets: $`A_S`$, the set of elements $`a`$ for which there is a finite subset $`F`$ of $`S`$ such that $`aF`$; the set $`A_T`$, defined similarly, and $`A_r`$, the rest of $`A`$. Applying $`𝔪_c(B)>\mathrm{}_0`$ we find for each $`aA_r`$ a nonzero element $`b_a`$ below $`a`$ that is in $`(A_SA_T)^{}`$ and we find $`bB`$ such that $`Sb^{}`$ and $`Tb^{}`$. Because $`|A|<𝔡_2(B)`$ we can find $`b_1<b`$ and $`b_2<b^{}`$ such that
* $`b_1S`$ and for all $`aA_r`$ if $`bb_a0`$ then $`b_1b_a0`$, and
* $`b_2T`$ and for all $`aA_r`$ if $`b^{}b_a0`$ then $`b_2b_a0`$
Because $`|A|<𝔯(B)`$ we can find $`c_1<b_1`$ that reaps all possible $`b_1b_a`$ and likewise we can find $`c_2<b_2`$. Finally then $`d=bc_1^{}c_2`$ is as required in the definition of $`\mathrm{}_0`$-ideal completeness. ∎
###### Remark 4.6.
We can now see that $`𝔪_c=𝔠`$ does not imply $`\mathrm{𝖬𝖠}_{\mathrm{countable}}`$; indeed, it is well-known that in the Laver model $`\mathrm{𝖬𝖠}_{\mathrm{countable}}`$ fails but that also $`𝔡=𝔯=𝔠`$.
## 5. General structure of $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebras
In this section we explore the general structure of $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebras. It is straightforward to check that “being an $`\mathrm{}_0`$-ideal subalgebra of” is a transitive relation.
###### Proposition 5.1.
If $`A`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B`$ and $`B`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`C`$, then $`A`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`C`$.∎
We have already mentioned Parovičenko’s theorem that under $`\mathrm{𝖢𝖧}`$ the algebra $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is the unique $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra with $`𝔪_c=𝔠`$. This leads us to the following definition.
###### Definition 5.2.
A Boolean algebra $`B`$ is Cohen-Parovičenko if $`B`$ is $`(,\mathrm{}_0)`$-ideal and $`𝔪_c(B)=𝔠`$.
In the special case when $`𝔠=\mathrm{}_2`$ we have a convenient characterization in terms of well-orderings at our disposal.
###### Proposition 5.3.
If $`𝔠=\mathrm{}_2`$ then an algebra $`B`$ of cardinality $`𝔠`$ is Cohen-Parovičenko if and only if for each enumeration $`B=\{b_\alpha :\alpha \omega _2\}`$ the set of $`\lambda \omega _2`$ for which $`B_\lambda =\{b_\alpha :\alpha \lambda \}`$ is both an $`\mathrm{}_0`$-ideal and an $`\mathrm{}_0`$-ideal complete subalgebra is closed and unbounded in $`\omega _2`$.∎
The reason for adding the prefix ‘Cohen’ is contained in the following proposition, which together with Theorem 5.5 gives a ‘factorization’ of Steprāns’ characterization of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ (Theorem 2.13). The proposition itself combines Theorem 2.7 and Proposition 2.12.
###### Proposition 5.4.
In the Cohen model $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is Cohen-Parovičenko.∎
The following theorem combines Parivičenko’s and Steprāns’ theorems into one. We have been unable to find any sort of similar result in the case that $`𝔠>\mathrm{}_2`$.
###### Theorem 5.5.
If $`𝔠\mathrm{}_2`$, then all Cohen-Parovičenko Boolean algebras of cardinality $`𝔠`$ are pairwise isomorphic.∎
To show that this theorem is never vacuous we now construct $`(,\mathrm{}_0)`$-ideal algebras with prescribed $`𝔪_c`$-numbers, including $`𝔠`$. Thus we see that the idealness of an algebra has no bearing on its $`𝔪_c`$-number. By contrast, a careful inspection of \[13, Proposition 2.3 and Corollary 2.4\] will reveal that if $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(\kappa ,\mathrm{}_0)`$-ideal and $`\mathrm{cf}[\kappa ]^\mathrm{}_0=\kappa `$ then $`𝔪_c>\kappa `$, thus showing that $`𝔪_c=𝔠`$ in case $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(,\mathrm{}_0)`$-ideal and there are cofinally many $`\kappa `$ below $`𝔠`$ with $`\mathrm{cf}[\kappa ]^\mathrm{}_0=\kappa `$.
###### Theorem 5.6.
For each regular cardinal $`\kappa 𝔠`$ such that $`\mathrm{cf}[\kappa ]^\mathrm{}_0=\kappa `$, there is an algebra of cardinality $`𝔠`$ that is $`(,\mathrm{}_0)`$-ideal and has $`\kappa `$ as its $`𝔪_c`$-number. Thus, if $`𝔠`$ is regular then there is Cohen-Parovičenko algebra of cardinality $`𝔠`$.
We split the construction into two propositions.
###### Proposition 5.7.
There is a $`(,\mathrm{}_0)`$-ideal algebra $`B`$ of size $`𝔠`$ with $`𝔪_c(B)\kappa `$.
###### Proof.
We obtain $`B`$ as the direct limit of a sequence $`B_\xi :\xi <\mu `$, where $`\mu `$ is the ordinal $`𝔠\kappa `$ if $`\kappa <𝔠`$ and $`\mu =𝔠`$ otherwise.
We begin by letting $`B_0`$ be the two-element algebra. At limit stages $`\xi `$ set $`B_\xi =\underset{}{\mathrm{lim}}_{\eta <\xi }B_\eta `$. Carry an enumeration, $`\{S_\xi ,T_\xi :\xi \mu \}`$ with cofinal repetitions, of pairs of countably infinite subsets of $`B`$ so that $`S_\xi T_\xi B_\xi `$ and $`S_\xi T_\xi `$ for all $`\xi `$. For simplicity, assume that $`S_\xi `$ and $`T_\xi `$ are strictly increasing sequences (if infinite) or singletons (if finite, where an empty set may always be replaced by $`\{0\}`$).
To construct $`B_{\xi +1}`$ from $`B_\xi `$ first take the completion $`\stackrel{~}{B}_\xi `$ of $`B_\xi `$ and in it we define $`s_\xi `$ and $`t_\xi `$ by $`s_\xi =S_\xi `$ and $`t_\xi ^{}=T_\xi `$; note that $`s_\xi t_\xi `$. There are two cases to consider.
If $`s_\xi <t_\xi `$ then we let $`B_{\xi +1}`$ be the subalgebra of $`\stackrel{~}{B}_\xi ^2`$ generated by the diagonal $`\{b,b:bB_\xi \}`$ and the element $`b_\xi =s_\xi ,t_\xi `$. Observe that $`s_\xi ,t_\xi `$ does exactly what is required in Definition 2.9 — with $`A=B_\xi `$, $`S=S_\xi `$ and $`T=T_\xi `$. Also observe that $`B_\xi `$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B_{\xi +1}`$: a typical element $`b`$ of $`B_{\xi +1}`$ looks like $`(b_\xi a_0)(b_\xi ^{}a_1)`$ and from this it is easily seen that the countable set $`A_b=\{(sa_0)(ta_1)(a_0a_1):sS_\xi ,tT_\xi \}`$ generates $`b^{}B_\xi `$: if $`ab`$ then $`aa_0a_1`$ and so we have to cover $`aa_1^{}`$ (which is below $`b_\xi a_0`$), $`aa_0^{}`$ (which is below $`b_\xi ^{}a_1`$) and $`aa_0a_1`$.
If $`s_\xi =t_\xi `$ then this still works if $`S_\xi `$ and $`T_\xi `$ are both infinite or both finite but not if, say, $`S_\xi `$ is infinite and $`T_\xi `$ is finite, for then we seek an element $`b_\xi `$ such that $`s<b_\xi <t_\xi `$ for all $`sS_\xi `$ — note that in this case $`T_\xi =\{t_\xi ^{}\}`$ and that $`t_\xi `$ belongs to $`B_\xi `$. To remedy this we take the Stone space $`X_\xi `$ of $`B_\xi `$ and consider the closed set $`C_\xi =s_\xi S_\xi `$. We let $`B_{\xi +1}`$ be the clopen algebra of the subspace $`Y_\xi =\left(X_\xi \times \{0\}\right)\left(C_\xi \times \{1\}\right)`$ of $`X_\xi \times \{0,1\}`$. Observe that $`b_\xi =s_\xi \times \{0\}`$ does what we want: for every $`sS_\xi `$ we have $`s<b_\xi <t_\xi `$, because $`t_\xi `$ now corresponds to $`b_\xi \left(C_\xi \times \{1\}\right)`$. A typical element $`b`$ of $`B_{\xi +1}`$ now looks like $`(b_\xi a_0)(c_\xi a_1)(t_\xi ^{}a_2)`$, where $`c_\xi =C_\xi \times \{1\}`$ — because $`C_\xi s_\xi `$ we have $`bB_\xi `$ iff we can take $`a_0=a_1`$. As above, we can verify that $`A_b=\{(sa_0)(t_\xi a_1)(t_\xi ^{}a_2):sS_\xi \}`$ generates $`b^{}B_\xi `$. If $`ab`$ then $`t_\xi ^{}at_\xi a_2`$, so we concentrate on $`a^{}=at_\xi =(ab_\xi )(ac_\xi )`$. Now $`a^{}a_1^{}b_\xi `$, so there is an $`sS_\xi `$ above $`a^{}a_1^{}`$. For this $`s`$ we have $`a(sa_0)(t_\xi a_1)(t_\xi ^{}a_2)`$.
We shall refer to $`\xi `$ as of type $`0`$ if we simply adjoin $`s_\xi ,t_\xi `$; the other $`\xi `$ will be of type $`1`$.
It is straightforward to check that $`𝔪_c(B)\kappa `$: if $`|A|<\kappa `$ then $`AB_\eta `$ for some $`\eta `$ and if $`S`$ and $`T`$ are countable subsets of $`A`$ with $`ST`$ then there is a $`\xi `$ above $`\eta `$ with $`S,T=S_\xi ,T_\xi `$; the element $`b_\xi `$ is as required for $`A`$, $`S`$ and $`T`$.
We show that $`B`$ is $`(,\mathrm{}_0)`$-ideal by showing that $`MB`$ is $`\mathrm{}_0`$-ideal in $`B`$ whenever $`M`$ is an elementary substructure of $`H(𝔠^+)`$ with $`<S_\xi ,T_\xi :\xi <\mu >`$ and $`<B_\xi :\xi <\mu >`$ both in $`M`$, and with $`|M|`$ less than $`𝔠`$ and regular.
Let $`bBM`$, take $`eb^{}M`$, and fix the $`\delta `$ and $`\xi `$ for which $`bB_{\delta +1}B_\delta `$ and $`eB_{\xi +1}B_\xi `$ respectively.
Claim 1. If $`\xi \delta `$ then there is an $`aA_b`$ with $`ea`$.
Proof. If $`\xi <\delta `$ then $`eB_\delta `$ and we are done.
If $`\xi >\delta `$ we consider two cases. If $`\xi `$ is of type $`0`$ and $`e=(b_\xi e_0)(b_\xi ^{}e_1)`$ then $`e_0b^{}b_\xi ^{}`$, hence $`e_0b^{}t`$ for some $`t`$ in $`T_\xi `$; likewise $`e_1b^{}s`$ for some $`s`$ in $`S_\xi `$. But then $`e(e_0t^{})(e_1s^{})b`$, where the middle element belongs to $`B_\xi `$; it follows that there is an $`aA_b`$ with $`ea`$.
If $`\xi `$ is of type $`1`$ and $`e=(b_\xi e_0)(c_\xi e_1)(t_\xi ^{}e_2)`$ then $`t_\xi ^{}e_2`$ belongs to $`B_\xi `$, so we concentrate on the other parts of $`e`$ — and we assume $`e_0,e_1t_\xi `$. Observe that $`b_\xi e_0t_\xi e_0b`$: use the fact that $`bB_\xi `$. Next $`e_1b^{}b_\xi `$, so that there is $`sS_\xi `$ with $`e_1b^{}s`$ and hence $`c_\xi e_1s^{}e_1b`$. We see that $`e(t_\xi e_0)(s^{}e_1)(t_\xi ^{}e_2)b`$, where the middle element belongs to $`B_\xi `$; again we can find our $`aA_b`$ with $`ea`$.
This claim essentially takes of the case $`\delta M`$: by our obvious inductive assumption we have for every $`aA_b`$ a countable generating set $`C_a`$ for $`a^{}M`$. By the claim the countable set $`C_b=_{aA_b}C_a`$ generates $`b^{}M`$ (note that $`\xi M`$, so $`\xi \delta `$).
To fully finish the proof we must show what to do if $`\delta M`$. The set $`C_b`$ still takes care of the $`e`$ with $`\xi \delta `$. The following two claims show what to add to $`C_b`$ in order to take care of the $`e`$ with $`\xi =\delta `$.
Claim 2. If $`\delta `$ is of type $`0`$, $`e=(b_\delta e_0)(b_\delta ^{}e_1)`$ and $`b=(b_\delta a_0)(b_\delta ^{}a_1)`$ then there are $`c_0C_{a_0}`$ and $`c_1C_{a_1}`$ such that $`e(b_\delta c_0)(b_\delta ^{}c_1)b`$.
Proof. Simply observe that $`b_\delta e_ia_i^{}M`$ for $`i=0`$$`1`$.
We see that we must add $`\{(b_\delta c_0)(b_\delta ^{}c_1):c_0C_{a_0},c_1C_{a_1}\}`$ to $`C_b`$.
Claim 3. If $`\delta `$ is of type $`1`$ and $`e=(b_\delta e_0)(c_\delta e_1)(t_\delta ^{}e_2)`$ and $`b=(b_\delta a_0)(c_\delta a_1)(t_\delta ^{}a_2)`$ Then there are $`c_0C_{a_0}`$, $`c_1C_{a_1}`$ and $`c_2C_{a_2}`$ such that $`e(b_\delta c_0)(c_\delta c_1)(t_\delta ^{}c_2)b`$.
Proof. Simply observe that $`b_\delta e_0a_0^{}M`$, $`c_\delta e_1a_1^{}M`$ and $`t_\delta ^{}e_2a_2^{}M`$.
Now add $`\{(b_\delta c_0)(c_\delta c_1)(t_\delta ^{}c_2):c_0C_{a_0},c_1C_{a_1},c_2C_{a_2}\}`$ to $`C_b`$. ∎
Note that if $`\kappa =𝔠`$ we are done: the algebra $`B`$ is $`(,\mathrm{}_0)`$-ideal with $`𝔪_c(B)=𝔠`$. In the case where $`\kappa <𝔠`$ we use the cofinality assumption to find a subalgebra of $`B`$ with the right properties.
###### Proposition 5.8.
If $`\kappa <𝔠`$ then the algebra $`B`$ constructed in the proof of Proposition 5.7 contains an algebra $`A`$ of cardinality $`𝔠`$ with $`𝔪_c(A)=\kappa `$.
###### Proof.
We fix a cofinal subfamily $`\{Y_\alpha :\alpha <\kappa \}`$ of $`[\kappa ]^\mathrm{}_0`$ with $`Y_\alpha \alpha `$ for all $`\alpha `$. We also assume that, for every $`\alpha <\kappa `$, all ordered pairs $`S,T`$ with $`S,TB_{𝔠\alpha }`$ occur in the list $`\{S_\xi ,T_\xi :𝔠\alpha \xi <𝔠(\alpha +1)\}`$. This enables us to choose, recursively, $`\lambda _\alpha [𝔠\alpha ,𝔠(\alpha +1))`$ such that $`S_{\lambda _\alpha }=\{b_{\lambda _\beta }:\beta Y_\alpha \}`$ and $`T_{\lambda _\alpha }=\{0\}`$. Note that then $`b_{\lambda _\beta }<b_{\lambda _\alpha }`$ whenever $`\beta Y_\alpha `$. In what follows we abbreviate $`b_{\lambda _\alpha }`$ by $`p_\alpha `$.
Because the $`Y_\alpha `$ form a cofinal family in $`[\kappa ]^\mathrm{}_0`$, the family $`\{p_\alpha :\alpha <\kappa \}`$ is $`\mathrm{}_0`$-directed, i.e., if $`F\kappa `$ is countable then there is an $`\alpha `$ such that $`p_\beta <p_\alpha `$ for all $`\beta F`$. It follows that $`I=\{b:(\alpha )(bp_\alpha )\}`$ is a $`P`$-ideal. We set $`F=\{b^{}:bI\}`$ and consider the subalgebra $`A=IF`$ of $`B`$. It is clear that $`𝔪_c(A)\kappa `$: no element of $`A`$ reaps the family $`\{p_\alpha :\alpha <\kappa \}`$.
To show $`𝔪_c(A)\kappa `$ we take a subalgebra $`D`$ of $`A`$ of size less than $`\kappa `$ and countable subsets $`S`$ and $`T`$ of $`D`$ with $`ST`$; we assume $`S`$ and $`T`$ are increasing sequences. Also, fix $`\alpha <\kappa `$ such that $`DB_{𝔠\alpha }`$ and for every $`dD`$ there is $`\beta <\alpha `$ with $`dp_\beta `$ or $`d^{}p_\beta `$. If some member of $`S`$ or $`T`$ belongs to $`F`$ then any $`bB`$ that witnesses this instance of $`\mathrm{}_0`$-ideal completeness of $`D`$ in $`B`$ automatically belongs to $`A`$.
In the other case, when $`STI`$, we can assume that $`Y_\alpha `$ contains, for every $`aST`$, a $`\beta `$ such that $`ap_\beta `$; but then $`STp_\alpha ^{}`$. Also note that $`p_\alpha `$ meets every nonzero element of $`B_{\lambda _\alpha }`$ and hence of $`D`$. Now choose $`\zeta [𝔠\alpha ,𝔠(\alpha +1))`$ such that $`S_\zeta =S\{p_\alpha ^{}\}`$ and $`T_\zeta =T`$. Let $`dDb_\zeta ^{}`$; there is an $`sS`$ with $`dsp_\alpha ^{}`$, then $`ds^{}p_\alpha ^{}`$ and so $`ds^{}=0`$ whence $`ds`$. We see that $`S`$ generates $`b_\zeta ^{}D`$ and, similarly, that $`T`$ generates $`b_\zeta ^{}D`$.
We finish by showing that $`A`$ is $`(,\mathrm{}_0)`$-ideal. The notation $`b^{}`$ will always mean the set computed in $`B`$. Let $`M`$ be any elementary substructure of $`H(𝔠^+)`$ of regular size less than $`𝔠`$ such that $`B_\xi :\xi <\mu `$ and $`\lambda _\alpha :\alpha \kappa `$ are members of $`M`$; this ensures that $`MB`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B`$. We shall show that for any $`b`$ in $`B`$, the ideal $`b^{}(MA)`$ is countably generated; we denote the countable generating set, when found, by $`b^{M,A}`$.
Fix $`\delta <\mu `$ so that $`bB_{\delta +1}B_\delta `$ and assume we have found $`a^{M,A}`$ for all $`aB_\delta `$. By Claim 5 in the proof of Proposition 5.7 the set $`_{aA_b}a^{M,A}`$ takes care of all $`eb^{}(MA)`$, except possibly those in $`B_{\delta +1}B_\delta `$. In particular we can set $`b^{M,A}=_{aA_b}a^{M,A}`$ when $`\delta M`$.
Thus we are left with the case where $`\delta M`$. If there is an $`eb^{}MF`$ then $`C_bF`$ generates $`b^{}MA`$, where $`C_b`$ is as defined in the proof of Proposition 5.7. In the other case, where $`b^{}MF=\mathrm{}`$, we add
$$\{(b_\delta c_0)(b_\delta ^{}c_1):c_0a_0^{M,A},c_1a_1^{M,A}\}$$
to $`b^{M,A}`$ if $`\delta `$ is of type $`0`$ and we add
$$\{(b_\delta c_0)(c_\delta c_1)(t_\delta ^{}c_2):c_0a_0^{M,A},c_1a_1^{M,A},c_2a_2^{M,A}\}$$
if $`\delta `$ is of type $`1`$.
Indeed, if $`eb^{}(MA)`$ then $`e`$ belongs to $`IM`$ and hence so do $`eb_\delta `$ and $`eb_\delta ^{}`$. Note that $`eb_\delta a_0`$ and $`eb_\delta ^{}a_1`$ so that there are $`c_0a_0^{M,A}`$ and $`c_1a_1^{M,A}`$ with $`eb_\delta c_0a_0`$ and $`eb_\delta ^{}c_1a_1`$ respectively.
If $`\delta `$ is of type $`1`$ then we observe that $`eb_\delta `$, $`ec_\delta `$ and $`et_\delta ^{}`$ all belong to $`IM`$ and are below $`a_0`$, $`a_1`$ and $`a_2`$ respectively. ∎
### Mapping $`F`$-spaces onto $`\beta \mathrm{}`$
Every compact $`F`$-space contains a copy of $`\beta \mathrm{}`$: it follows straight from the definition of $`F`$-space that the closure of a countably infinite relatively discrete subset is homeomorphic to $`\beta \mathrm{}`$. Thus, in a manner of speaking, $`\beta \mathrm{}`$ is a minimal $`F`$-space. Bell has asked whether $`\beta \mathrm{}`$ is also minimal in the mapping-onto sense: does every infinite compact zero-dimensional $`F`$-space map onto $`\beta \mathrm{}`$? The ease with which $`\beta \mathrm{}`$ can be embedded into such a space belies the dual difficulty in constructing an embedding of $`\mathrm{𝒫}(\mathrm{})`$ into its algebra of clopen sets. Indeed, we show by means of a Cohen-Parovičenko algebra that such an embedding does not always exist. Before that we prove that Bell’s question has a positive answer if the Continuum Hypothesis is assumed.
###### Proposition 5.9 ($`\mathrm{𝖢𝖧}`$).
Every infinite compact zero-dimensional $`F`$-space maps onto $`\beta \mathrm{}`$.
###### Proof.
It suffices to prove that $`\mathrm{𝒫}(\mathrm{})`$ will embed into $`B`$ where $`B`$ is infinite and has no $`(\omega ,\omega )`$-gaps. Fix any sequence $`\{b_n:n\omega \}`$ of pairwise disjoint non-zero elements of $`B`$. Let $`\{a_\alpha :\alpha \omega _1\}`$ be an enumeration of $`\mathrm{𝒫}(\mathrm{})`$ so that $`a_n=\{n\}`$ for each $`n\omega `$. Inductively choose elements $`b_\alpha B`$ so that the mapping $`a_\alpha b_\alpha `$ lifts to an isomorphism from the algebra generated by $`\{a_\beta :\beta \alpha \}`$. If $`a_\alpha `$ is in the algebra generated by its predecessors then there is nothing to do. Otherwise, by the inductive hypothesis, the ideal $`\mathrm{}`$ generated by $`\{b_\beta :a_\beta <a_\alpha \}`$ is disjoint from the ideal $`\mathrm{𝒥}`$ generated by $`\{b_\beta :a_\beta a_\alpha =0\}`$. Since $`B`$ has no $`(\omega ,\omega )`$-gaps, there is a $`b_\alpha B`$ such that $`\mathrm{}b_\alpha ^{}`$ and $`\mathrm{𝒥}b_\alpha ^{}`$. ∎
###### Theorem 5.10.
It is consistent that there is an infinite compact zero-dimensional $`F`$-space that does not map onto $`\beta \mathrm{}`$.
###### Proof.
It is consistent with $`𝔠=\mathrm{}_2`$ that $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ contains an $`\omega _2`$-chain, this happens, e.g., if $`\mathrm{𝖬𝖠}`$ holds. But now let $`B`$ be the Cohen-Parovičenko algebra from Theorem 5.6. Clearly $`S(B)`$ is a compact zero-dimensional $`F`$-space. Assume that $`f`$ maps $`S(B)`$ onto $`\beta \mathrm{}`$ and let $`b_n=f^{}(n)`$ for each $`n`$. Now let $`\mathrm{}`$ be the ideal in $`B`$ generated by $`\{b_n:n\omega \}`$. By the forthcoming Corollary 5.12 $`B/\mathrm{}`$ is still $`(\mathrm{}_1,\mathrm{}_0)`$-ideal and so, by Proposition 2.5, does not contain an $`\omega _2`$-chain. However $`B/\mathrm{}`$ is isomorphic to the algebra of clopen subsets of the closed set $`X_nb_n=f^{}(\mathrm{}^{})`$ and certainly does contain $`\omega _2`$-chains. ∎
This proof does not work in the $`\mathrm{}_2`$-Cohen model, where $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is *the* Cohen-Parovičenko algebra. We therefore ask, also in the hope of establishing the consistency with $`\neg \mathrm{𝖢𝖧}`$ of a yes answer to Bell’s question, the following.
###### Question 1.
Is it true in the $`\mathrm{}_2`$-Cohen model that every compact zero-dimensional $`F`$-space does map onto $`\beta \mathrm{}`$?
### Quotient algebras
Under $`\mathrm{𝖢𝖧}`$ one can use Parovičenko’s theorem to find many copies of $`\mathrm{}^{}`$ inside of $`\mathrm{}^{}`$: the proof usually boils down to showing that a quotient of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ by some ideal is isomorphic to $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$. The same can be done in the Cohen model because many quotients of Cohen-Parovičenko algebras are again Cohen-Parovičenko. First we consider quotients by small ideals.
###### Lemma 5.11.
If $`B`$ is a Boolean algebra, $`A`$ is an $`\mathrm{}_0`$-ideal subalgebra, and $`\mathrm{}`$ is an ideal which is generated by $`\mathrm{}A`$, then $`A/\mathrm{}`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B/\mathrm{}`$.
###### Proof.
Fix any $`bB`$ and fix a cofinal sequence $`\{a_n:n\omega \}b^{}A`$. Let $`cA`$ be such that $`c/\mathrm{}<b/\mathrm{}`$, which means that $`cb`$ is covered by some member $`d`$ of $`\mathrm{}A`$. It follows then that $`cd<b`$. Hence there is $`n`$ such that $`cd<a_n<b`$. But now it follows that $`c/\mathrm{}<a_n/\mathrm{}`$. ∎
###### Corollary 5.12.
If $`\mathrm{}`$ is an $`\mathrm{}_1`$-generated ideal in a $`(\kappa ,\mathrm{}_0)`$-ideal Boolean algebra $`B`$, then $`B/\mathrm{}`$ is also a $`(\kappa ,\mathrm{}_0)`$-ideal Boolean algebra.∎
Another interesting consequence is that $`\omega _1^{}`$ is not the image of the Stone space of an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra.
###### Corollary 5.13.
The algebra $`P(\omega _1)/\mathrm{𝑓𝑖𝑛}`$ cannot be embedded into an $`(\mathrm{}_1,\mathrm{}_0)`$-ideal algebra.
###### Proof.
This proceeds much as the proof of Theorem 5.10 since $`P(\omega _1)/\mathrm{𝑐𝑡𝑏𝑙𝑒}`$ is a quotient of $`P(\omega _1)/\mathrm{𝑓𝑖𝑛}`$ by an $`\mathrm{}_1`$-generated ideal and contains $`\omega _2`$-chains. ∎
###### Corollary 5.14 ($`\neg \mathrm{𝖢𝖧}`$).
If $`B`$ is Cohen-Parovičenko and $`\mathrm{}`$ is an $`\mathrm{}_1`$-generated ideal then $`B/\mathrm{}`$ is again Cohen-Parovičenko.
###### Proof.
It remains only to prove that $`𝔪_c(B/\mathrm{})=𝔠`$. To do so, fix countable subsets $`S`$ and $`T`$ of $`B`$ so that $`st\mathrm{}`$ for each $`sS`$ and $`tT`$. Since $`S`$ and $`T`$ are countable it is routine to inductively remove from each member of $`S`$ and $`T`$ some member of $`\mathrm{}`$ so as to ensure that $`st=0`$ for each $`sS`$ and $`tT`$. Now suppose that $`A`$ is a subalgebra of $`B`$ that contains $`S`$ and $`T`$ and has cardinality less than $`𝔠`$. We may assume that $`A`$ contains a generating set for $`\mathrm{}`$. Since $`B`$ is Cohen-Parovičenko there is a $`bB`$ such that $`S`$ generates $`b^{}A`$ and $`T`$ generates $`b^{}A`$. Now suppose $`a/\mathrm{}`$ is below $`b/\mathrm{}`$, i.e., $`ab\mathrm{}`$. Since $`A\mathrm{}`$ generates $`\mathrm{}`$, there is a $`cA\mathrm{}`$ such that $`ab<c`$. Thus $`ac<b`$ and so there is a finite join, $`s`$, of members of $`S`$ such that $`ac<s`$. It follows that $`\{s/\mathrm{}:sS\}`$ generates $`(b/\mathrm{})^{}A/\mathrm{}`$. Similarly $`(b/\mathrm{})^{}A/\mathrm{}`$ is generated by $`\{t/\mathrm{}:tT\}`$. ∎
Another situation that occurs frequently is that one has a lifting for the ideal $`\mathrm{}`$, this is a Boolean homomorphism $`l:B/\mathrm{}B`$ with the property that $`l(b/\mathrm{})/\mathrm{}=b/\mathrm{}`$. In dual terms this means that the closed set $`F=S(B)\{i^{}:i\mathrm{}\}`$ is a retract of $`S(B)`$. The retraction $`r`$ and the lifting $`l`$ are connected by the formula $`l(b/\mathrm{})=r^{}[b^{}F]`$ for each $`bB`$.
###### Theorem 5.15.
If $`\mathrm{}`$ is an ideal on $`B`$ for which there is a lifting $`l:B/\mathrm{}B`$ then for each $`\mathrm{}_0`$-ideal subalgebra $`A`$ of $`B`$ such that $`l[A/\mathrm{}]A`$ the quotient $`A/\mathrm{}`$ is an $`\mathrm{}_0`$-ideal subalgebra of $`B/\mathrm{}`$. Therefore, if $`B`$ is an $`(\kappa ,\mathrm{}_0)`$-ideal Boolean algebra, then so is $`B/\mathrm{}`$.
###### Proof.
Let $`A`$ be an $`\mathrm{}_0`$-ideal subalgebra of $`B`$ such that $`l[A/\mathrm{}]A`$. Fix any $`bB`$. We will show that $`(b/\mathrm{})^{}A/\mathrm{}`$ is countably generated. In fact, suppose that $`\{a_n:n\omega \}`$ generates $`l(b/\mathrm{})^{}A`$. Fix any $`xA`$ such that $`x/\mathrm{}<b/\mathrm{}`$. By assumption, $`x^{}=l(x/\mathrm{})`$ is also a member of $`A`$. Furthermore $`l(x/\mathrm{})l(b/\mathrm{})`$, hence there is an $`n`$ such that $`x^{}a_n`$. Clearly then $`x/\mathrm{}=x^{}/\mathrm{}a_n/\mathrm{}`$. ∎
## 6. Other remainders and applications to $`\mathrm{}^{}`$
We say that a zero-dimensional space $`K`$ is Cohen-Parovičenko if its algebra of clopen sets is Cohen-Parovičenko. In this section we are interested in identifying which remainders of $`\sigma `$-compact locally compact spaces can be Cohen-Parovičenko; by ‘remainder’ we mean the Čech-Stone remainder $`\beta XX`$ — commonly denoted by $`X^{}`$. We then apply this information and the results of the previous section to the study of $`\mathrm{}^{}`$ under the assumption that it is Cohen-Parovičenko. We are motivated by the somewhat classical results about $`\mathrm{}^{}`$ that are known to follow from $`\mathrm{𝖢𝖧}`$ (see ). The predisposition of this section is to assume that $`\mathrm{}^{}`$ is Cohen-Parovičenko and to determine how this affects the structure of $`\mathrm{}^{}`$ and of other remainders.
In what follows, whenever $`X`$ is a zero-dimensional compact space, we write $`𝔪_c(X)`$, $`𝔯(X)`$ and $`𝔡_2(X)`$ for the values that these functions have on the Boolean algebra $`\mathrm{CO}(X)`$ of clopen subsets of $`X`$. We first prove a lemma concerning the behaviour of $`𝔡_2`$ and $`𝔯`$ under continuous mappings.
###### Lemma 6.1.
If $`f:XY`$ is an open continuous surjection then $`𝔡_2(Y)𝔡_2(X)`$ and $`𝔯(Y)𝔯(X)`$.
###### Proof.
Let $`\mathrm{}`$ be an ideal of $`\mathrm{CO}(Y)`$, co-generated by the family $`\{c_n:n\omega \}`$, and let $`A`$ be a subfamily of $`\mathrm{}`$ of size less than $`𝔡_2(X)`$. In $`\mathrm{CO}(X)`$ we can find an element $`b`$ such that $`bf^{}[c_n]=\mathrm{}`$ for all $`n`$ and $`bf^{}[a]\mathrm{}`$ for all $`aA`$. Because the map $`f`$ is open the set $`f[b]`$ is clopen, it also belongs to $`\mathrm{}`$ and it meets every element of $`A`$.
Next let $`\mathrm{𝒞}`$ be a family of clopen subsets of $`X`$, of size less than $`𝔯(Y)`$. Because $`f`$ is open the family $`\{f[c]:c\mathrm{𝒞}\}`$ consists of clopen sets and so we can find $`b\mathrm{CO}(Y)`$ that reaps it. Then $`f^{}[b]`$ reaps the family $`\mathrm{𝒞}`$. ∎
Our first result is somewhat surprising. It implies that if $`\mathrm{𝖢𝖧}`$ fails then most remainders are not Cohen-Parovičenko. Recall that a space is *basically disconnected* if each cozero-set has clopen closure — dually: the algebra of clopen subsets is countably complete. Unless stated otherwise the spaces we are considering are all zero-dimensional.
###### Proposition 6.2.
Let $`X`$ be the topological sum of countably many compact spaces that are not basically disconnected. If its remainder $`X^{}`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal then $`𝔡=\mathrm{}_1`$.
###### Proof.
Write $`X=_{n\omega }X_n`$ and fix for each $`n`$ an infinite family $`\{a(n,m):m\omega \}`$ of pairwise disjoint clopen sets of $`X_n`$ so that $`D_n`$, the closure of their union, is not open.
Assume that $`M`$ is an $`\mathrm{}_0`$-covering elementary substructure of some $`H(\theta )`$ of size $`\mathrm{}_1`$ that contains $`X`$ and the family $`\{a(n,m):n,m\omega \}`$. We show that $`M\mathrm{}^{\mathrm{}}`$ is cofinal in $`\mathrm{}^{\mathrm{}}`$. Let $`f:\mathrm{}\mathrm{}`$ be a strictly increasing function not in $`M`$. We find $`gM`$ such that $`fg`$.
Let $`b=\{a(n,m):mf(n)\}`$; observe that $`b`$ is also not in $`M`$. We take a countable subfamily $`C`$ of $`M\mathrm{CO}(X)`$ that is cofinal in $`b^{}M`$ — this means that $`cb`$ is compact for all $`cC`$ and that whenever $`dM`$ and $`db`$ is compact there is $`cC`$ such that the difference $`dc`$ is compact.
There are two cases to consider. If there is a $`cC`$ such that the set $`I_c=\{n:(m)(ca(n,m)\mathrm{})\}`$ is infinite then we are done. Indeed, define $`hM`$ by $`h(n)=\mathrm{min}\{m:ca(n^+,m)\mathrm{}\}`$, where $`n^+=\mathrm{min}\{ln:lI_c\}`$. Because $`cb`$ is compact there is an $`l`$ such that $`ca(n,m)=\mathrm{}`$ whenever $`nl`$ and $`mf(n)`$. It follows that for $`nl`$ we have $`h(n)=h(n^+)>f(n^+)>f(n)`$. Now define $`g`$ by $`g(n)=\mathrm{max}\{h(n),f(n)\}`$; this $`g`$ belongs to $`M`$ because it is a finite modification of $`h`$ and it is as required.
In the other case, where $`I_c`$ is finite for all $`c`$, we may assume $`CM`$: indeed, take a countable $`DM`$ with $`CD`$ and replace $`C`$ by the set of elements $`d`$ of $`D`$ for which $`I_d`$ is finite. By subtracting a compact part from each $`c`$ we can also assume that every $`c`$ is disjoint from every $`a(n,m)`$.
But now from an enumeration $`\{c_n:n\omega \}`$ of $`C`$ (that is in $`M`$) we define the clopen set
$$c=\{X_n(c_0\mathrm{}c_n):n\omega \}.$$
It follows that $`c`$ is in $`Mb^{}`$. Now for each $`n`$, $`D_nc`$ is empty, but $`D_n`$ is not equal to $`X_nc`$ since $`D_n`$ is not open. Therefore, there is some $`dMb^{}`$ such that $`cd`$ and, for each $`n`$, $`dX_nc`$ is not empty. It follows that $`dc_k`$ is not compact for any $`k`$ and so $`\{c_k:k\omega \}`$ is not a generating set for $`b^{}M`$. Therefore this case does not occur. ∎
###### Theorem 6.3 ($`\neg \mathrm{𝖢𝖧}`$).
If $`X=_{n\omega }X_n`$ is the topological sum of countably many compact spaces that are not basically disconnected then the remainder of $`X`$ is not Cohen-Parovičenko.
###### Proof.
Choose $`x_nX_n`$ for all $`n`$ and observe that $`D=\mathrm{cl}\{x_n:n\omega \}X^{}`$ is homeomorphic to $`\mathrm{}^{}`$. The map that send $`X_n`$ to the point $`x_n`$ induces an open retraction from $`X^{}`$ onto $`D`$. It follows that $`𝔡𝔡_2(X^{})`$, so if $`\mathrm{CO}(X^{})`$ is $`(\mathrm{}_1,\mathrm{}_0)`$-ideal then, by Proposition 6.2, we get $`𝔡_2(X^{})=\mathrm{}_1<𝔠`$. ∎
###### Remark 6.4.
Clearly it follows from the previous result that if $`\mathrm{}^{}`$ is Cohen-Parovičenko and $`\mathrm{𝖢𝖧}`$ fails then $`\left(\omega \times (\omega +1)\right)^{}`$ is not homeomorphic to $`\omega ^{}`$. Using this fact and tracking the location of both clopen and nowhere dense $`P`$-set copies of $`\mathrm{}^{}`$ in their remainders, one can easily show that, in addition, $`\omega \times (\omega +1)`$ and $`\omega \times (\omega ^2+1)`$ do not have homeomorphic remainders either.
It is also worth mentioning the following result since it has already found applications in Functional Analysis, see .
###### Corollary 6.5 ($`\neg `$CH).
If $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is Cohen-Parovičenko and $`C`$ is a non-compact cozero set in $`\mathrm{}^{}`$, then the closure of $`C`$ is not a retract of $`\mathrm{}^{}`$.
###### Proof.
Let $`C`$ be a non-compact cozero subset of $`\mathrm{}^{}`$. It follows that $`C`$ is a countable union of compact open subsets of $`\mathrm{}^{}`$ and, as is well-known, that the closure of $`C`$ is just its Čech-Stone compactification. Now if the closure were a retract of $`\mathrm{}^{}`$, then its clopen algebra would be an $`(\mathrm{}_1,\mathrm{}_0)`$-algebra. The boundary of $`C`$, which is homeomorphic to $`\beta CC`$, is a $`G_\delta `$-set in the closure of $`C`$; hence its clopen algebra is also an $`(\mathrm{}_1,\mathrm{}_0)`$-algebra by Lemma 5.11.
However, no clopen subset of $`\mathrm{}^{}`$ is basically disconnected so by Proposition 6.2 we have $`𝔡=\mathrm{}_1`$. But we assumed that $`𝔪_c`$ and hence $`𝔡`$ was equal to $`𝔠`$. ∎
With the previous results in mind it is tempting to hope that for $`\sigma `$-compact locally compact $`X`$ and $`Y`$, if $`X^{}`$ and $`Y^{}`$ were homeomorphic then $`X`$ and $`Y`$ would be homeomorphic-modulo-compact-sets in some sense. For example, we do not know if $`(\omega \times 2^\omega )^{}`$ and $`(\omega \times 2^{\omega _1})^{}`$ are homeomorphic in the Cohen model.
We do however know of other spaces whose remainder is Cohen-Parovičenko. The proof of this fact is a rather interesting use of the basic results we have developed about Cohen-Parovičenko algebras. Recall that the Gleason cover or absolute of a compact space $`X`$ is denoted by $`E(X)`$ and that $`E(X)`$ is just the Stone space of the complete Boolean algebra of regular open subsets of $`X`$. We write $`E_\kappa `$ for $`\omega \times E(2^\kappa )`$.
###### Lemma 6.6.
Let $`X=_{n\omega }X_n`$ be the topological sum of basically disconnected compact spaces. Then $`𝔡_2(X^{})=𝔡`$.
###### Proof.
Lemma 6.1 gives us $`𝔡𝔡_2(X^{})`$. To prove the other inequality we take an ideal in $`\mathrm{CO}(X^{})`$ that is co-generated by a strictly increasing sequence. Translating this into $`\mathrm{CO}(X)`$ we get an increasing sequence $`C_k:k\omega `$ of clopen sets in $`X`$ such that for all $`k`$ the difference $`C_{k+1}C_k`$ is not compact (for convenience we assume $`C_0=\mathrm{}`$), and the ideal $`\mathrm{}`$ in $`\mathrm{CO}(X)`$ consisting of those sets $`D`$ for which every intersection $`DC_k`$ is compact.
For every $`n,k\omega `$ put $`a(n,k)=X_n(C_{k+1}C_k)`$. This transforms the $`C_k`$ into an increasing sequence $`c_k:k\omega `$ of subsets of the countable set $`A=\{a(n,k):n,k\omega \}`$, where $`c_k=\{a(n,l):n\omega ,l<k\}`$. To every $`D\mathrm{}`$ corresponds the set $`x_D=\{aA:Da\mathrm{}\}`$; observe that $`x_Dc_k`$ is finite for each $`k`$. Because each $`X_n`$ is basically disconnected the sets $`D_n=\mathrm{cl}_ka(n,k)`$ are clopen (maybe empty); we put $`Y=X_nD_n`$.
Let $`\mathrm{𝒥}`$ be a subfamily of $`\mathrm{}`$ of size less than $`𝔡`$ and consisting of non-compact sets. Fix an infinite subset $`d`$ of $`A`$ such that $`dc_k`$ is finite for all $`k`$ and such that $`dx_D`$ is infinite whenever $`D\mathrm{𝒥}`$ and $`x_D`$ is infinite. Finally put $`C=Yd`$. Clearly $`CC_k`$ is compact for every $`k`$. If $`D\mathrm{𝒥}`$ and $`x_D`$ is finite then $`DY`$ is not compact; if $`x_D`$ is infinite then $`Dd`$ is not compact. ∎
###### Remark 6.7.
A similar result does not hold for $`𝔯`$. Indeed, consider the space $`E_\kappa `$; a clopen set in its remainder $`E_\kappa ^{}`$ is determined by a clopen set of $`E_\kappa `$ itself. In turn a clopen subset of $`E_\kappa `$ is determined by a regular open subset of $`\omega \times 2^\kappa `$ and it is well-known that such a regular open set depends on at most countably many coordinates. Thus, if $`\mathrm{𝒞}`$ is a family of fewer than $`\kappa `$ many clopen sets in $`E_\kappa ^{}`$ then we can find an $`\alpha \kappa `$ such that no element of $`\mathrm{𝒞}`$ depends on $`\alpha `$. But this means that the clopen set $`\omega \times \pi _\alpha ^{}(0)`$ (or rather the clopen subset of $`E_\kappa ^{}`$ determined by it) reaps the family $`\mathrm{𝒞}`$. We deduce that $`𝔯(E_\kappa ^{})\kappa `$ and hence that, for example, $`𝔯(E_𝔠^{})>𝔯`$ in models where $`𝔠>𝔯`$.
###### Theorem 6.8.
For each cardinal $`\kappa 𝔠`$, the remainder of $`E_\kappa `$ is Cohen-Parovičenko iff $`\mathrm{}^{}`$ is Cohen-Parovičenko.
###### Proof.
We start out by observing two partial equivalences.
Claim 1. $`𝔡=𝔠`$ iff $`𝔡_2(E_\kappa ^{})=𝔠`$.
Proof. By Lemma 6.6 we know that $`𝔡_2(E_\kappa ^{})=𝔡`$ for all $`\kappa `$.
Claim 2. The algebra $`\mathrm{CO}(E_\kappa ^{})`$ is $`(,\mathrm{}_0)`$-ideal iff $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is.
Proof. This follows by applying Theorem 5.15 twice. First: $`\mathrm{}^{}`$ is easily seen to be a retract of $`E_\kappa ^{}`$, so $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(,\mathrm{}_0)`$-ideal if $`\mathrm{CO}(E_\kappa ^{})`$ is. Second: $`\beta E_\kappa `$ is a separable extremally disconnected compact space and hence can be embedded as retract in $`\mathrm{}^{}`$, so that $`\mathrm{CO}(\beta E_\kappa ^{})`$ is $`(,\mathrm{}_0)`$-ideal if $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is and, by Corollary 5.12, so is the clopen algebra of $`E_\kappa ^{}`$.
We would be done if we could also prove $`𝔯(E_\kappa ^{})=𝔯`$ but by Remark 6.7 we know that this cannot be done. We circumvent this difficulty by showing that $`𝔯𝔡`$ if $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(,\mathrm{}_0)`$-ideal. This will follow from the following technical lemma, which is in the spirit of Proposition 2.3 of , whose content was explained just before Theorem 5.6. ∎
###### Lemma 6.9.
If $`\kappa <𝔡`$ and $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ is $`(\kappa ,\mathrm{}_0)`$-ideal then also $`\kappa <𝔯`$.
###### Proof.
It suffices to show that if $`MH(\theta )`$ is $`\mathrm{}_0`$-covering, of size $`\kappa `$ and such that $`M\mathrm{𝒫}(\mathrm{})`$ is $`\mathrm{}_0`$-ideal in $`\mathrm{𝒫}(\mathrm{})`$ then there is an $`r\mathrm{𝒫}(\mathrm{})`$ that reaps $`M[\mathrm{}]^\mathrm{}_0`$. By Van Douwen’s characterization of $`𝔡`$ (see Definition 4.2) there is $`f{}_{}{}^{\mathrm{}}\mathrm{}`$ such that for every $`xM[\mathrm{}]^\mathrm{}_0`$ and every $`gM{}_{}{}^{\mathrm{}}\mathrm{}`$ there is an $`nx`$ such that $`f(n)g(n)`$. Fix a countable subset $`C`$ of $`M{}_{}{}^{\mathrm{}}\mathrm{}`$ such that for every subset $`a`$ of $`\mathrm{}\times \mathrm{}`$ with $`aL_f`$ there is $`cC`$ such that $`aL_c`$. As we can assume $`CM`$ and because $`M`$ knows that $`C`$ is countable we can find $`gM{}_{}{}^{\mathrm{}}\mathrm{}`$ such that $`c<^{}g`$ for all $`cC`$. We claim that $`r=\{n:f(n)g(n)\}`$ is as required.
Now let $`xM[\mathrm{}]^\mathrm{}_0`$; the choice of $`f`$ implies that $`rx`$ is infinite. To show that $`r^{}x`$ is infinite consider $`a=L_g(x\times \omega )`$. Clearly there is no $`cC`$ with $`a^{}L_c`$, hence $`aL_f`$ is infinite; this gives infinitely many $`n`$ with $`g(n)>f(n)`$. ∎
###### Remark 6.10.
Many of the foregoing consequences of $`\mathrm{}^{}`$ being Cohen-Parovičenko do need the assumption of $`\neg \mathrm{𝖢𝖧}`$. For example, it is shown in that a homeomorphism between nowhere dense $`P`$-set subsets of $`\mathrm{}^{}`$ can be lifted to a homeomorphism on $`\mathrm{}^{}`$. In addition, Steprāns proves that all $`P`$-points can be taken to one another by autohomeomorphisms of $`\mathrm{}^{}`$ in the Cohen model (and it appears that only the assumption that $`𝔪_c=\mathrm{}_2=𝔠`$ is used). However we can provide the following elegant contrasting result.
###### Proposition 6.11.
If $`𝔠=\mathrm{}_2`$ and if $`\mathrm{}^{}`$ is Cohen-Parovičenko then there are two $`P`$-sets in $`\mathrm{}^{}`$, of character $`\mathrm{}_1`$ and $`\mathrm{}_2`$ respectively, that are both homeomorphic to $`\mathrm{}^{}`$.
###### Proof.
Using Theorem 6.8 we see that $`E_𝔠^{}`$ is Cohen-Parovičenko. We may therefore apply Theorem 5.5 to deduce that $`\mathrm{}^{}`$ and $`E_𝔠^{}`$ are homeomorphic. Now fix one point $`x`$ in $`E(2^𝔠)`$; the set $`\left(\omega \times \{x\}\right)^{}`$ is a $`P`$-set of character $`𝔠`$ in $`E_𝔠^{}`$ and clearly homeomorphic to $`\mathrm{}^{}`$.
To get a $`P`$-set of character $`\mathrm{}_1`$ we take a strictly decreasing chain $`a_\alpha :\alpha <\omega _1`$ of clopen sets in $`\mathrm{}^{}`$ whose intersection $`A`$ is nowhere dense in $`\mathrm{}^{}`$ — see Remark 2.11 for the construction. Clearly then $`A`$ is a $`P`$-set of character $`\mathrm{}_1`$. The ideal $`\mathrm{}`$ generated by the family $`\{a_\alpha ^{}:\alpha <\omega _1\}`$ is $`\mathrm{}_1`$-generated, so by Corollary 5.14 the algebra $`(\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛})/\mathrm{}`$ is Cohen-Parovičenko and hence isomorphic to $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$. Its Stone space is $`A`$, which consequently is homeomorphic to $`\mathrm{}^{}`$. ∎
## 7. Problems
### Other reals
The Cohen model is probably the most intensively investigated model of $`\neg \mathrm{𝖢𝖧}`$ of all; this may explain our success in extracting key features of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ and $`\mathrm{}^{}`$ in that model. It would be of great interest if a similar thing could be done for other familiar models of $`\neg \mathrm{𝖢𝖧}`$.
The Laver and Sacks (also side-by-side) models are particular favourites of the authors but the Random real model seems the most likely candidate for a successful investigation.
### Characterizing $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$
Theorems 2.13 and 5.5 lead one to hope that there is a characterization of $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ in any Cohen model. As first steps on the way to such a result we ask the following questions.
###### Question 2.
Is, in the $`\mathrm{}_3`$-Cohen model, $`\mathrm{𝒫}(\mathrm{})/\mathrm{𝑓𝑖𝑛}`$ the unique Cohen-Parovičenko algebra?
Or, more generally.
###### Question 3.
If $`𝔠=\mathrm{}_3`$ are then all Cohen-Parovičenko algebras of cardinality $`𝔠`$ isomorphic? |
warning/0001/math-ph0001011.html | ar5iv | text | # The kernel of Fock representations of Wick algebras with braided operator of coefficients
## 1. Introduction
The problem of positivity of the Fock space inner product is central in the study of the Fock representation of Wick algebras (see , , , ). The paper presents several conditions on the coefficients of the Wick algebra for the Fock inner product to be positive. If the operator of coefficients of the Wick algebra $`T`$ satisfies the braid condition and the norm restriction $`T1`$, then, as proved in , the Fock inner product is positive. Moreover if $`1<T<1`$, it was shown in that the Fock inner product is strictly positive. In this article we prove that, for braided $`T`$ with $`T1`$, the kernel of the Fock inner product coincides with the largest quadratic Wick ideal. In particular this implies that, for $`1<T1`$, the Fock inner product is strictly positive definite, and the Fock representation is faithful.
This article is organized as follows. In sec. 2 we present definitions of Wick algebras and the Fock representation and show that, in the braided case, the kernel of the Fock representation is generated by the kernel of the Fock inner product. In sec. 3 we prove that if the operator $`T`$ is braided and $`T1`$, then the kernel of the Fock inner product coincides with the two-sided ideal generated by $`\mathrm{ker}(1+T)`$. In sec. 4 we combine results obtained in sec. 2 and sec. 3 to examine the $`C^{}`$-representability of certain Wick algebras or their quotients. All results are illustrated by examples of different kinds of $`q_{ij}`$-CCR.
## 2. Preliminaries
For more detailed information about Wick algebras and the Fock representation we refer the reader to . In this section we present only the basic definitions and properties.
1. The notion of a $``$-algebra allowing Wick ordering (Wick algebra) was presented in the paper as a generalization of a wide class of $``$-algebras, including the twisted CCR and CAR algebras (see ), the $`q`$-CCR (see ) algebra, etc.
###### Definition 1.
Let $`𝕁=𝕁_d=\{1,2,\mathrm{},d\}`$, $`T_{ij}^{kl}C`$, $`i,j,k,l𝕁`$, be such that $`T_{ij}^{kl}=\overline{T_{ji}^{lk}}`$. The Wick algebra with the set of coefficients $`\{T_{ij}^{kl}\}`$ is denoted $`W(T)`$, and is a $``$-algebra, defined by generators $`a_i`$, $`a_i^{}`$, $`i𝕁`$, which satisfy the basic relations:
$$a_i^{}a_j=\delta _{ij}1+\underset{k,l=1}{\overset{d}{}}T_{ij}^{kl}a_la_k^{}$$
###### Definition 2.
Monomials of the form $`a_{i_1}a_{i_2}\mathrm{}a_{i_m}a_{j_1}^{}a_{j_2}^{}\mathrm{}a_{j_k}^{}`$ are called Wick ordered monomials.
It was proved in that the Wick ordered monomials form a basis for $`W(T)`$.
Let $`=e_1,\mathrm{},e_d`$. Consider the full tensor algebra over $`,^{}`$, denoted by $`𝒯(,^{})`$. Then
$$W(T)𝒯(,^{})/e_i^{}e_j\delta _{ij}1T_{ij}^{kl}e_le_k^{}.$$
To study the structure of Wick algebras, and the structure of the Fock representation, it is useful to introduce the following operators on $`^n:=\underset{n}{\underset{}{\mathrm{}}}`$ (see ):
$`T:`$ $`,Te_ke_l={\displaystyle \underset{i,j}{}}T_{ik}^{lj}e_ie_j,T=T^{},`$
$`T_i:`$ $`^n^n,T_i=\underset{i1}{\underset{}{1\mathrm{}1}}T\underset{ni1}{\underset{}{1\mathrm{}1}},`$
$`R_n:`$ $`^n^n,R_n=1+T_1+T_1T_2+\mathrm{}+T_1T_2\mathrm{}T_{n1},`$
$`P_n:`$ $`^n^n,P_2=R_2,P_{n+1}=(1P_n)R_{n+1}.`$
In this article we suppose that the operator $`T`$ is contractive, i.e., $`T1`$, and satisfies the braid condition, i.e., on $`^{\mathrm{\hspace{0.17em}3}}`$ the equality $`T_1T_2T_1=T_2T_1T_2`$ holds. It follows from the definition of $`T_i`$ that then $`T_iT_j=T_jT_i`$ if $`\left|ij\right|2`$, and for the braided $`T`$ one has $`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1}`$.
###### Remark 1.
These conditions hold for such well-known algebras as $`q_{ij}`$-CCR, $`\mu `$\- CCR, $`\mu `$-CAR (see ).
The Fock representation of a Wick $``$-algebra is determined by a vector $`\mathrm{\Omega }`$ such that $`a_i^{}\mathrm{\Omega }=0`$ for all $`i=1,\mathrm{},d`$ (see ).
###### Definition 3 (The Fock representation).
The representation $`\lambda _0`$, acting on the space $`𝒯()`$ by formulas
$`\lambda _0(a_i)e_{i_1}\mathrm{}e_{i_n}`$ $`=e_ie_{i_1}\mathrm{}e_{i_n},n\{0\},`$
$`\lambda _0(a_i^{})1`$ $`=0,`$
where the action of $`\lambda _0(a_i^{})`$ on the monomials of degree $`n1`$ is determined inductively using the basic relations, is called the Fock representation.
Note that the Fock representation is not a $``$-representation with respect to the standard inner product on $`𝒯()`$. However, it was proved in that there exists a unique Hermitian sesquilinear form $`<,>_0`$ on $`𝒯()`$ such that $`\lambda _0`$ is a $``$-representation on $`(𝒯(),<,>_0)`$. This form is called the Fock inner product on $`𝒯()`$.
The subspaces $`^n`$, $`^m`$, $`nm`$, are orthogonal with respect to $`<,>_0`$, and on $`^n`$ we have the following formula (see ):
$$<X,Y>_0=<X,P_nY>,n2.$$
So, the positivity of the Fock inner product is equivalent to the positivity of operators $`P_n`$, $`n2`$, and $`=_{n2}\mathrm{ker}P_n`$ determines the kernel of the Fock inner product. It was noted in that the Fock representation is the GNS representation associated with the linear functional $`f`$ on a Wick algebra such that $`f(1)=1`$ and, for any Wick ordered monomial, $`f(a_{i_1}\mathrm{}a_{i_n}a_{j_1}^{}\mathrm{}a_{j_m}^{})=0`$. Then for any $`X,Y𝒯()`$ we have (see ):
$$<X,Y>_0=f(X^{}Y).$$
2. In the following proposition we describe the kernel of the Fock representation of a Wick algebra with braided operator $`T`$ in terms of the Fock inner product.
###### Proposition 1.
Let $`W(T)`$ be the Wick algebra with braided operator $`T`$, and let the Fock representation $`\lambda _0`$ be positive (i.e., the Fock inner product is positive definite). Then $`\mathrm{ker}\lambda _0=𝒯(^{})+𝒯()^{}`$.
###### Proof.
First, we show that $`X\mathrm{ker}P_m`$ implies $`X\mathrm{ker}\lambda _0`$. Indeed, let $`Y^n`$; then
$$\lambda _0(X)Y=XY.$$
Note that for braided $`T`$ we have the following decomposition (see and sec. 3 for more details):
$$P_{n+m}=P(D_m)(P_m\mathrm{𝟏}_n),$$
where
$`P(D_m)=\stackrel{~}{R}_{n+m}\stackrel{~}{R}_{n+m1}\mathrm{}\stackrel{~}{R}_{m+1},`$
$`\stackrel{~}{R}_k=1+T_{k1}+T_{k2}T_{k1}+\mathrm{}+T_1T_2\mathrm{}T_{k1},k2.`$
Then
$$P_{n+m}(\lambda _0(X)Y)=P_{n+m}(XY)=P(D_m)(P_mXY)=0,$$
and $`\lambda _0(X)=0`$ on $`(𝒯(),<,>_0)`$. Therefore $`\mathrm{ker}\lambda _0`$, and since $`\mathrm{ker}\lambda _0`$ is a $``$-ideal,
$$𝒯(^{})+𝒯()^{}\mathrm{ker}\lambda _0.$$
(1)
To prove the converse inclusion, we need a formula determining the action of $`\lambda _0(X^{})`$ on $`𝒯()`$ for any $`X^k`$, $`k`$. For $`k=1`$, $`X=e_i`$, $`i=1,\mathrm{},d`$, it was proved in that:
$$\lambda _0(e_i^{})Y=\mu (e_i^{})R_nY,Y^n,$$
where $`\mu (e_i^{}):𝒯()𝒯()`$ is the annihilation operator:
$$\mu (e_i^{})e_{i_1}e_{i_2}\mathrm{}e_{i_n}=\delta _{ii_1}e_{i_2}\mathrm{}e_{i_n}.$$
Then, using the definition of $`P_n`$, it is easy to see that, for $`X^n`$ and $`Y^n`$,
$$\lambda _0(X^{})Y=<X,P_nY>=<X,Y>_0.$$
Let now
$$Z=\underset{i=1}{\overset{n}{}}Y_iX_i^{}+\underset{j=n+1}{\overset{l}{}}Y_jX_j^{}\mathrm{ker}\lambda _0,$$
where $`Y_i𝒯()`$, $`i=1,\mathrm{},l`$,
$$X_i^m,i=1,\mathrm{},n,X_j^{n_j},n_j>m,j=n+1,\mathrm{},l.$$
Now (1) implies that we can suppose that the elements $`X_i`$ are linearly independent modulo $``$. Denote by $`\{\widehat{X}_i,i=1,\mathrm{},n\}^m`$ a family dual to the $`\{X_i,i=1,\mathrm{},n\}`$ with respect to $`<,>_0`$, i.e., such that
$$<X_i,P_m\widehat{X}_j>=<X_i,\widehat{X}_j>_0=\delta _{ij},i,j=1,\mathrm{},n.$$
Since, for any $`j=n+1,\mathrm{},l`$ and $`i=1,\mathrm{},n`$,
$$\lambda _0(X_j^{})\widehat{X}_i=0,$$
we have, in $`(𝒯(),<,>_0)`$,
$$0=\lambda _0(Z)\widehat{X}_i=Y_i,i=1,\mathrm{},n,$$
which implies $`Y_i`$, $`i=1,\mathrm{},n`$. The proof can be completed by evident induction. ∎
###### Remark 2.
In particular, we have shown, for braided $`T`$, and for any $`X^n`$ and $`Y\mathrm{ker}P_m`$, that
$$XY\mathrm{ker}P_{n+m}.$$
By similar arguments, $`YX\mathrm{ker}P_{n+m}`$, i.e., $`=_{n2}\mathrm{ker}P_n`$ is a two-sided ideal in $`𝒯()`$.
The two-sided ideal $`𝒥𝒯()`$ is called a Wick ideal (see ) if it satisfies the following condition:
$$𝒯\left(^{}\right)𝒥𝒥𝒯\left(^{}\right).$$
(2)
If $`𝒥`$ is generated by some subspace of $`^n`$, then $`𝒥`$ is called a homogeneous Wick ideal of degree $`n`$.
We show that for Wick algebras with braided operator of coefficients, $``$ is a Wick ideal .
###### Proposition 2.
Let $`T`$ satisfy the braid condition, and $`=_{n2}\mathrm{ker}P_n`$; then
$$^{}+^{}.$$
(3)
###### Proof.
Note that conditions (2) and (3) are equivalent (see ). To prove the proposition, it is sufficient to show that, if $`X\mathrm{ker}P_n`$ for some $`n2`$, then for any $`i=1,\mathrm{},d`$,
$$e_i^{}X\mathrm{ker}P_{n1}+\mathrm{ker}P_n^{}.$$
Indeed, for any $`X^n`$, we have the following formula (see ):
$$e_i^{}X=\mu (e_i^{})R_nX+\mu (e_i^{})\underset{k=1}{\overset{d}{}}T_1T_2\mathrm{}T_n(Xe_k)e_k^{}.$$
Then for $`X\mathrm{ker}P_n`$, we have
$$P_{n1}\mu (e_i^{})R_nX=\mu (e_i^{})(1P_{n1})R_nX=\mu (e_i^{})P_nX=0.$$
Note that, for braided $`T`$, for any $`k=2,\mathrm{},n`$,
$$T_k(T_1T_2\mathrm{}T_n)=(T_1T_2\mathrm{}T_n)T_{k1},$$
which implies that
$$(1P_n)(T_1T_2\mathrm{}T_n)=(T_1T_2\mathrm{}T_n)(P_n1).$$
Then for any $`k=1,\mathrm{},d`$,
$`P_n\mu (e_i^{})T_1T_2\mathrm{}T_n(Xe_k)`$ $`=\mu (e_i^{})(1P_n)T_1T_2\mathrm{}T_n(Xe_k)`$
$`=\mu (e_i^{})T_1T_2\mathrm{}T_n(P_n1)(Xe_k)=0.\text{ }\text{}`$
For Wick algebras with braided $`T`$, the largest homogeneous ideal of degree $`n`$ is generated by $`\mathrm{ker}R_n`$ (see and ), i.e., the condition $`\mathrm{ker}R_n\{0\}`$ is necessary and sufficient for the existence of homogeneous Wick ideals. In the following proposition we show that the same is true for arbitrary Wick ideals.
###### Theorem 1.
If $`𝒥𝒯()`$ is a non-trivial Wick ideal, then there exists $`n2`$ such that $`\mathrm{ker}R_n\{0\}`$.
###### Proof.
For any $`X𝒯()`$, by $`\mathrm{deg}X`$ we denote the highest degree of its homogeneous components. Let $`Y𝒥`$ be of minimal degree.
$$Y=Y_1+Y_2+\mathrm{}+Y_k,Y_i^{n_i},i=1,\mathrm{},k,n_i\{0\}.$$
Suppose that $`\mathrm{deg}Y2`$: then for any $`i=1,\mathrm{},d`$, we have
$$e_i^{}Y=\underset{j=1}{\overset{k}{}}\mu (e_i^{})R_{n_j}Y_j+\mu (e_i^{})\underset{j=1}{\overset{k}{}}\underset{l=1}{\overset{d}{}}\stackrel{~}{T}_{n_j}(Y_je_l)e_l^{},$$
where we put $`R_0=1`$, $`R_1=1`$, and
$$\stackrel{~}{T}_k=\{\begin{array}{cc}T_1T_2\mathrm{}T_k,\hfill & k2,\hfill \\ T,\hfill & k=1,\hfill \\ 1,\hfill & k=0.\hfill \end{array}$$
Then condition (3) implies that for any $`i=1,\mathrm{},d`$,
$$\underset{j=1}{\overset{k}{}}\mu (e_i^{})R_{n_j}Y_j𝒥.$$
Since the degrees of these elements are less than the degree of $`Y`$, we conclude that
$$\underset{j=1}{\overset{k}{}}\mu (e_i^{})R_{n_j}Y_j=0,i=1,\mathrm{},d,$$
and the independence of the Wick ordered monomials then implies
$$\mu (e_i^{})R_{n_j}Y_j=0,i=1,\mathrm{},d.$$
Let $`Y_k`$ be the highest homogeneous component of $`Y`$; then, by our assumption, $`\mathrm{deg}Y_k2`$, and $`_{i=1}^de_i\mu (e_i^{})R_{n_k}Y_k=R_{n_k}Y_k=0`$, i.e., $`Y_k\mathrm{ker}R_{n_k}`$.
To complete the proof, note that if $`X=\beta +_{i=1}^d\alpha _ie_i𝒥`$, then for any $`j`$, we have
$$e_j^{}X=\alpha _j+\beta e_j^{}+\underset{i=1}{\overset{d}{}}\alpha _i\underset{k,l=1}{\overset{d}{}}T_{ji}^{kl}e_le_k^{},$$
and (3) implies $`\alpha _j=0`$, $`j=1,\mathrm{},d`$, $`\beta =0`$. ∎
## 3. The structure of $`\mathrm{ker}P_n`$
In this section, we show that for Wick algebras with braided $`T`$ satisfying the condition $`1<T1`$, the Fock representation is faithful, and for $`1T1`$, the kernel of the Fock representation is generated by the largest quadratic Wick ideal (the largest quadratic Wick ideal is the largest homogeneous Wick ideal of degree $`2`$).
To do this we need some properties of quasimultiplicative maps on the Coxeter group $`S_n`$ (for more detailed information we refer the reader to ).
1. Consider $`S_{n+1}`$ as a Coxeter group, i.e., a group defined as follows: $`S_{n+1}=\sigma _i:\sigma _i^2=e,\sigma _i\sigma _j=\sigma _j\sigma _i,|ij|2,\sigma _i\sigma _{i+1}\sigma _i=\sigma _{i+1}\sigma _i\sigma _{i+1},i=1,\mathrm{},n`$. In order to study the invertibility of $`P_n`$ for any family of operators $`\{T_i,i=1\mathrm{},n,TB(𝒦)\}`$, satisfying the conditions
$$T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1},T_i^{}=T_i,1T_i1,$$
where $`𝒦`$ is a separable Hilbert space, we may define (as in ) the function
$$\varphi :S_{n+1}B(𝒦)$$
by the formulas
$$\varphi (e)=1,\varphi (\sigma _i)=T_i,$$
(4)
$$\varphi (\pi )=T_{i_1}\mathrm{}T_{i_k},$$
(5)
where $`\pi =\sigma _{i_1}\mathrm{}\sigma _{i_k}`$ is a reduced decomposition. It was shown in that
$$P_{n+1}=P(S_{n+1})=\underset{\sigma S_{n+1}}{}\varphi (\sigma ).$$
Denote by $`S`$ the set of generators of $`S_{n+1}`$ as a Coxeter group. Consider, for any $`JS`$, the set
$$D_J=\{\sigma S_{n+1}\left|\sigma s\right|=\left|\sigma \right|+1,sJ\}.$$
Let $`W_J`$ be a Coxeter group, generated by $`J`$. Then $`S_{n+1}=D_JW_J`$ (see ), and $`P_{n+1}=P(D_J)P(W_J)`$ (see ). Using the equalities $`P_{n+1}^{}=P_{n+1}`$, $`P(W_J)^{}=P(W_J)`$, we obtain $`P_{n+1}=P(W_J)P(D_J)^{}`$, where for all $`MS_{n+1}`$,
$$P(M)=\underset{\sigma M}{}\varphi (\sigma ).$$
In what follows we use a quasimultiplicative analogue of the Euler-Solomon formula (see \[2, Lemma 2.6\]):
$$\underset{\begin{array}{c}JS\\ JS,J\mathrm{}\end{array}}{}(1)^{\left|J\right|}P(D_J)=(1)^{\left|S\right|}1+\varphi (\sigma _0^{(n+1)})P(S_{n+1}),$$
(6)
where $`\sigma _0^{(n+1)}`$ is the unique element of $`S_{n+1}`$ with maximal possible length of the reduced decomposition.
###### Remark 3.
1. The element $`\sigma _0^{(n+1)}`$ of the group $`S_{n+1}`$ has the form
$$\sigma _0^{(n+1)}=(\sigma _1\mathrm{}\sigma _n)(\sigma _1\mathrm{}\sigma _{n1})\mathrm{}(\sigma _1\sigma _2)\sigma _1.$$
Set $`U_n=\varphi (\sigma _0^{(n+1)})`$: then
$$U_n=(T_1T_2\mathrm{}T_n)(T_1T_2\mathrm{}T_{n1})\mathrm{}(T_1T_2)T_1.$$
2. It is easy to see that the operator $`U_n`$ is selfadjoint, and, taking adjoints, we can rewrite (6) in the following form:
$$\underset{JS,J\mathrm{}}{}(1)^{\left|J\right|}P(D_J)^{}=(1)^{n+1}1+U_nP_{n+1}.$$
(7)
3. Note also that, for all $`JS`$, the group $`W_J`$ is isomorphic to $`S_k`$ for some $`k<n`$, or to the direct product of some such groups.
2. In what follows we shall use the following properties of the operator $`U_n`$.
###### Proposition 3.
$`\mathrm{ker}P_{n+1}`$ is invariant with respect to the action of $`U_n`$.
###### Proof.
First we show that for all $`JS`$,
$$P(D_J^{}):\mathrm{ker}P_{n+1}\mathrm{ker}P_{n+1}.$$
It can be easily obtained from the equality
$$P_{n+1}P(D_J)^{}=P(D_J)P(W_J)P(D_J)^{}=P(D_J)P_{n+1}.$$
Then by (7), we have
$$U_n(1)^n1:\mathrm{ker}P_{n+1}\mathrm{ker}P_{n+1}.\text{ }\text{}$$
###### Proposition 4.
Let operators $`\{T_i,i=1,\mathrm{},n\}`$ satisfy the braid condition $`T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1}`$, $`i=1,\mathrm{},n1`$, and $`T_iT_j=T_jT_i`$, $`\left|ij\right|2`$. Then
$$T_kU_n=U_nT_{n+1k},k=1,\mathrm{}n.$$
(8)
###### Proof.
1. For $`n=1`$ the equality is evident.
2. Suppose that (8) holds for any $`nm`$. Note that
$$U_{m+1}=T_1T_2\mathrm{}T_{m+1}U_m.$$
Then, for $`1<km+1`$, we have
$`T_kU_{m+1}`$ $`=T_k(T_1T_2\mathrm{}T_{m+1})U_m`$
$`=T_1T_2\mathrm{}T_{k2}T_kT_{k1}T_kT_{k+1}\mathrm{}T_{m+1}U_m`$
$`=T_1T_2\mathrm{}T_{k2}T_{k1}T_kT_{k1}T_{k+1}\mathrm{}T_{m+1}U_m`$
$`=(T_1T_2\mathrm{}T_{m+1})T_{k1}U_m`$
$`=T_1T_2\mathrm{}T_{m+1}U_mT_{m+1(k1)}`$
$`=U_{m+1}T_{m+2k}.`$
In particular, for $`k=m+1`$ we have $`T_{m+1}U_{m+1}=U_{m+1}T_1`$. Then taking adjoints, we obtain the required equality for $`k=1`$. ∎
3. Now we can formulate the main result of this paper.
###### Theorem 2.
Let $`W(T)`$ be a Wick algebra with braided operator $`T`$ satisfying the norm bound $`T1`$. Then for any $`n2`$, we have
$$\mathrm{ker}P_{n+1}=\underset{k+l=n1}{}^k\mathrm{ker}(1+T)^l=\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k).$$
###### Proof.
In fact, we shall prove the following: Let $`T_1,T_2,\mathrm{},T_nB(𝒦)`$, where $`𝒦`$ is a finite-dimensional Hilbert space, be selfadjoint contractions satisfying the relations
$$T_iT_{i+1}T_i=T_{i+1}T_iT_{i+1},1in1,T_iT_j=T_jT_i,\left|ij\right|2.$$
Then
$$\mathrm{ker}P_{n+1}=\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k).$$
(9)
(It follows trivially from the decomposition $`P_{n+1}=P(D_{\{k\}})(1+T_k)`$ that $`_{k=1}^n\mathrm{ker}(1+T_k)\mathrm{ker}P_{n+1}`$).
We proceed using induction.
The case $`n=2`$.
In this case $`P_2=1+T`$.
The case $`nn+1`$.
It follows from $`P_{n+1}=P(W_J)P(D_J)^{}`$ that
$$P(D_J)^{}:\mathrm{ker}P_{n+1}\mathrm{ker}P(W_J),$$
i.e., $`\mathrm{ran}(P(D_J)^{}|_{\mathrm{ker}P_{n+1}})\mathrm{ker}P(W_J)`$. Moreover, it is obvious that for any $`JS`$, $`\mathrm{ker}P(W_J)\mathrm{ker}P_{n+1}`$. Therefore, by (7), we have the following inclusion:
$$\mathrm{ran}(U_n(1)^n1)|_{\mathrm{ker}P_{n+1}}\underset{JS,J\mathrm{}}{}\mathrm{ker}P(W_J).$$
Since, for $`JS`$, the group $`W_J=W_{J_1}\times \mathrm{}\times W_{J_k}`$, where $`W_{J_l}S_{n_l}`$ with $`n_l<n+1`$, we have a decomposition into the product of pairwise commuting selfadjoint operators
$$P(W_J)=P(W_{J_1})\mathrm{}P(W_{J_k}).$$
Therefore
$$\mathrm{ker}P(W_J)=\underset{l=1}{\overset{k}{}}\mathrm{ker}P(W_{J_l})\underset{i=1}{\overset{n}{}}\mathrm{ker}(1+T_i),$$
where the last inclusion is obtained from the assumption of induction. So,
$$\mathrm{ran}(U_n(1)^n1)|_{\mathrm{ker}P_{n+1}}\underset{i=1}{\overset{n}{}}\mathrm{ker}(1+T_i).$$
(10)
Consider the operator $`1U_n^2`$. Since $`U_n=U_n^{}:\mathrm{ker}P_{n+1}\mathrm{ker}P_{n+1}`$, then
$$\mathrm{ker}P_{n+1}=\mathrm{ran}(1U_n^2)|_{\mathrm{ker}P_{n+1}}+\mathrm{ker}(1U_n^2)|_{\mathrm{ker}P_{n+1}}.$$
Moreover, since $`\mathrm{ran}(1U_n^2)\mathrm{ran}(U_n(1)^n1)`$, using (10), we have the inclusion
$$\mathrm{ker}P_{n+1}\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_i)+\mathrm{ker}(1U_n^2)|_{\mathrm{ker}P_{n+1}}.$$
To finish the proof it remains only to show that
$$\mathrm{ker}(1U_n^2)\mathrm{ker}P_{n+1}\underset{i=1}{\overset{n}{}}\mathrm{ker}(1+T_i).$$
To this end, we may present $`1U_n^2`$ in the form
$`1U_n^2`$ $`=1T_1T_2\mathrm{}T_nU_{n1}^2T_n\mathrm{}T_2T_1`$
$`=(1T_1^2)+T_1(1T_2^2)T_1`$
$`+\mathrm{}`$
$`+T_1T_2\mathrm{}T_{n1}(1T_n^2)T_{n1}\mathrm{}T_2T_1`$
$`+T_1T_2\mathrm{}T_n(1U_{n1}^2)T_n\mathrm{}T_2T_1.`$
Since $`T1`$ implies that $`T_i1`$, $`i=1,\mathrm{}n`$, and $`U_k1`$, $`k2`$, then we have a sum of non-negative operators, and $`v\mathrm{ker}(1U_n^2)`$ implies that
$`T_1^2v=v,`$
$`T_2^2T_1v=T_1v,`$
$`\mathrm{}`$
$`T_n^2T_{n1}\mathrm{}T_2T_1v=T_{n1}\mathrm{}T_2T_1v.`$
However, $`T_kU_n=U_nT_{n+1k}`$ implies that $`T_kU_n^2=U_n^2T_k`$, and, consequently,
$$T_k:\mathrm{ker}(1U_n^2)\mathrm{ker}(1U_n^2),k=1,\mathrm{},n.$$
Moreover, since the restriction of $`T_1`$ to $`\mathrm{ker}(1U_n^2)`$ is an involution,
$$\mathrm{ran}(T_1)|_{\mathrm{ker}(1U_n^2)}=\mathrm{ker}(1U_n^2),$$
and, for any $`v\mathrm{ker}(1U_n^2)`$, we have $`T_2^2v=v`$. By the same arguments, we obtain that
$$v\mathrm{ker}(1U_n^2),T_i^2v=v,i=1,\mathrm{},d.$$
Let now $`v\mathrm{ker}(1U_n^2)\mathrm{ker}P_{n+1}`$; then, for any $`k=1,\mathrm{},n`$,
$`P_{n+1}T_kv`$ $`=P(D_{\{k\}})(1+T_k)T_kv`$
$`=P(D_{\{k\}})(T_k+T_k^2)v`$
$`=P(D_{\{k\}})(1+T_k)v=P_{n+1}v=0.`$
Therefore $`T_k`$ maps $`\mathrm{ker}(1U_n^2)\mathrm{ker}P_{n+1}`$ onto itself for any $`k=1,\mathrm{},n`$. This fact implies that, for any $`\sigma S_{n+1}`$, we have $`\varphi (\sigma )v\mathrm{ker}(1U_n^2)\mathrm{ker}P_{n+1}`$, and
$$k=1,\mathrm{},n,\sigma S_{n+1},(1T_k^2)\varphi (\sigma )v=0.$$
For convenience, we fix the set $`S_{n+1}`$, and set $`v_i:=\varphi (\pi _i)v`$ for $`\pi _iS_{n+1}`$ ($`\pi _1:=\mathrm{id}`$ and $`v_1=v`$). Then the condition $`P_{n+1}v=0`$ takes the form
$$\underset{k=1}{\overset{n!}{}}v_k=0.$$
(11)
Finally, for any pair $`ij`$ there exist generators $`\sigma _{i_1},\mathrm{},\sigma _{i_m}S`$ such that
$$\pi _j=\sigma _{i_1}\mathrm{}\sigma _{i_m}\pi _i$$
and
$$v_j=T_{i_1}\mathrm{}T_{i_m}v_i.$$
Note that, if $`v_k=T_rv_l`$ for some $`r=1,\mathrm{},n`$, then $`T_r^2v_k=v_k`$ implies that $`v_kv_l\mathrm{ker}(1+T_r)`$. Therefore, for any $`ij`$,
$$v_iv_j\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k).$$
In particular, for any $`j=2,\mathrm{},n!`$,
$$v_1v_j\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k).$$
Then from (11), we have
$$n!v=n!v_1\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k),$$
and therefore
$$v\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k).\text{ }\text{}$$
###### Remark 4.
Evidently the proof does not depend on the dimension of $`𝒦`$. Indeed, in the case when $`𝒦`$ is infinite-dimensional, the linear subspace in (9) is replaced by its closure. I.e., if $`𝒦`$ is a separable Hilbert space and $`\{T_i,i=1,\mathrm{},n\}`$ are selfadjoint contractions satisfying the braid conditions, then
$$\mathrm{ker}P_{n+1}=\overline{\underset{k=1}{\overset{n}{}}\mathrm{ker}(1+T_k)}.$$
As a corollary we have an improved version of the result of Bożejko and Speicher (see ).
###### Proposition 5.
If the operator $`T`$ satisfies the braid condition, and $`1<T1`$, then $`P_n>0`$, $`n2`$, i.e., the Fock inner product is strictly positive, and the Fock representation acts in the whole space $`𝒯()`$.
###### Proof.
Recall that if $`T`$ is braided and $`T1`$ then $`P_n0`$ (see ). It remains only to show that $`\mathrm{ker}P_n=\{0\}`$ for $`1<T1`$. This fact trivially follows from our theorem since in this case $`\mathrm{ker}(1+T)=\{0\}`$. ∎
## 4. Corollaries and examples
We summarize the results obtained above in the following proposition.
###### Proposition 6.
If $`W(T)`$ is a Wick algebra with braided operator of coefficients $`T`$ satisfying the norm bound $`T1`$, then the following three statements hold.
1. The kernel of the Fock representation is generated by the largest quadratic Wick ideal. In particular, if $`1<T1`$, then the Fock representation is faithful.
2. For any $`n2`$ we have the inclusion $`_n_2`$.
3. If $`1<T1`$, then $`W(T)`$ has no non-trivial Wick ideals.
###### Example 1.
Consider the $`q`$-CCR algebra based on a Hilbert space $``$ and the relations
$`a_i^{}a_i`$ $`=1+qa_ia_i^{},i=1,\mathrm{},d,`$
$`a_i^{}a_j`$ $`=qa_ja_i^{},ij,0<q<1.`$
We pick an orthogonal basis $`(e_i)`$ in $``$, and then $`T`$ is determined on this basis by the formulas
$$Te_ie_j=qe_je_i,T<1.$$
It is evident that $`T`$ is braided. Then by the proposition, we cannot have any Wick ideals in $`W(T)`$.
It was proved in that for braided $`T`$ satisfying the norm bound $`T<1`$, the Fock representation is bounded. Therefore we may consider the $`C^{}`$-algebra generated by operators of the Fock representation.
Recall that a $``$-algebra is called $`C^{}`$-representable if it can be realized as a $``$-subalgebra of a certain $`C^{}`$-algebra (see for example ). Combining the results of Theorem 2 and Proposition 5, we obtain the following statement.
###### Proposition 7.
If $`W(T)`$ is a Wick algebra with braided operator of coefficients $`T`$ satisfying the norm bound $`T<1`$, then $`W(T)`$ is $`C^{}`$-representable.
Suppose that, in the case of braided $`T`$ with $`T=1`$ and $`\mathrm{ker}(1+T)\{0\}`$, the Fock representation is bounded. Then Theorem 2 implies that the quotient $`W(T)/_2`$ is $`C^{}`$-representable.
###### Example 2.
Consider the following type of $`q_{ij}`$-CCR (see ):
$`a_i^{}a_i`$ $`=1+q_ia_ia_i^{},i=1,\mathrm{},d,0<q_i<1,`$
$`a_i^{}a_j`$ $`=\lambda _{ij}a_ja_i^{},ij,\left|\lambda _{ij}\right|=1,\lambda _{ij}=\overline{\lambda }_{ij}.`$
The corresponding $`T`$ is braided, $`T=1`$, and
$$\mathrm{ker}(1+T)=<a_ja_i\lambda _{ij}a_ia_j,i<j>.$$
Moreover, the Fock representation of this algebra is bounded. Then, as noted above, the $``$-algebra generated by the relations
$`a_i^{}a_i`$ $`=1+q_ia_ia_i^{},i=1,\mathrm{},d,0<q_i<1,`$
$`a_i^{}a_j`$ $`=\lambda _{ij}a_ja_i^{},ij,\left|\lambda _{ij}\right|=1,\lambda _{ij}=\overline{\lambda }_{ij},`$
$`a_ja_i`$ $`=\lambda _{ij}a_ia_j,i<j`$
is $`C^{}`$-representable.
A description of the irreducible representations of these relations can be found for example in \[8, sec. 2.4\].
Note that, if $`T=1`$, then the operators of the Fock representation can be unbounded.
###### Example 3.
Consider the following Wick algebra:
$`a_i^{}a_i`$ $`=1+a_ia_i^{},i=1,\mathrm{},d,`$
$`a_i^{}a_j`$ $`=qa_ja_i^{},ij,1<q<1.`$
The corresponding $`T`$ is determined by the formulas
$$Te_ie_i=e_ie_i,Te_je_i=qe_ie_j,ij,i=1,\mathrm{},d.$$
It is easy to see that $`T`$ is braided and $`1<T1`$. So, the Fock representation of this algebra is faithful. Note that, if we consider the complement of $`𝒯()`$ with respect to the Fock inner product, then the operators of the Fock representation are unbounded.
For the definition and properties of representations of $``$-algebras by unbounded operators, see for example .
Unbounded representations of Wick algebras will be considered in more detail later.
Acknowledgements. Yu. Samoĭlenko and Dan. Proskurin express their gratitude to Prof. Vasyl L. Ostrovskyĭ and Stanislav Popovych for their attention and helpful discussions. We also thank Brian Treadway for proofreading and polishing. |
warning/0001/physics0001035.html | ar5iv | text | # What Dimensions Do the Time and Space Have: Integer or Fractional?
## I Introduction
The problem concerning the nature of space and time is one of the most interesting problems of the modern physics. Are the space and time continuous? Why is time irreversible? What dimensions do space and time have? How is the nature of time in the equations of modern physics is reflected? Different approaches (quantum gravity, irreversible thermodynamics, synergetics and others) provide us with different answers to these questions. In this paper the hypothesis about a nature of time and space based on an ideas of the fractal geometry is offered. The corresponding mathematical methods this hypothesis makes use of are based on using the idea about fractional dimensions (FD) as the main characteristics of time and space and in connection with this the generalization of the Riemann-Liouville fractional derivatives are introduced. The method and theory are developed to describe dynamics of functions defined on multifractal sets of time and space with FD.
Following , we will consider both time and space as an only material fields existing in the Universe and generating all other physical fields. Assume that each of them consists of a continuous, but not differentiable bounded set of small elements. Let us suppose a continuity, but not a differentiability, of sets of small time intervals (from which time consist) and small space intervals (from which space consist). First, let us consider set of small time intervals $`S_t`$ (for the set of small space intervals the way of reasoning is similar). Let time be defined on multifractal set of such intervals (determined on the carrier of a measure $`_t^n`$ ). Each of intervals of this set (further we use the approximation in which the description of each multifractal interval of these sets will be characterized by middle time moment $`t`$ and refer to each of these intervals as ”points”) is characterized by global fractal dimension (FD) $`d_t(𝐫(t),t)`$, and for different intervals FD are different( because of the time dependence and spatial coordinates dependence of $`d_t`$ ). For multifractal sets $`S_t`$ (or $`S_r`$) each set is characterized by global FD of this set and by local FD of this set( the characteristics of local FD of time and space sets in this paper we do not research). In this case the classical mathematical calculus or fractional (say, Riemann - Liouville) calculus can not be applied to describe a small changes of a continuous function of physical values $`f(t)`$, defined on time subsets $`S_t`$, because the fractional exponent depends on the coordinates and time. Therefore, we have to introduce integral functionals (both left-sided and right-sided) which are suitable to describe the dynamics of functions defined on multifractal sets (see ). Actually, this functionals are simple and natural generalization the Riemann-Liouville fractional derivatives and integrals:
$$D_{a+,t}^df(t)=\left(\frac{d}{dt}\right)^n_a^t\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(tt^{})^{d(t^{})n+1}}$$
(1)
$$D_{b,t}^df(t)=(1)^n\left(\frac{d}{dt}\right)^n_t^b\frac{f(t^{})dt^{}}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)^{d(t^{})n+1}}$$
(2)
where $`\mathrm{\Gamma }(x)`$ is Euler’s gamma function, and $`a`$ and $`b`$ are some constants from $`[0,\mathrm{})`$. In these definitions, as usually, $`n=\{d\}+1`$ , where $`\{d\}`$ is the integer part of $`d`$ if $`d0`$ (i.e. $`n1d<n`$) and $`n=0`$ for $`d<0`$. Functions under the integral sign we will consider to be generalized functions defined on the space of finite functions . Similar expressions can be written down for GFD of functions $`f(𝐫,t)`$ with respect to spatial variables $`𝐫`$, with $`f(𝐫,t)`$ being defined on the elements of set $`S_𝐫`$ whose dimension is $`d_𝐫`$.
For an arbitrary $`f(t)`$ it is useful to expand the generalized function $`1/(tt^{})^{\epsilon (t^{})}`$ under the integral sign in (1)-(2) into a power series in $`\epsilon (t^{})`$ when $`d=n+\epsilon ,\epsilon +0`$ and write
$`D_{a+,t}^df(t)=\left({\displaystyle \frac{d}{dt}}\right)^n{\displaystyle _a^t}{\displaystyle \frac{f(t^{})}{\mathrm{\Gamma }(nd(t^{}))(tt^{})}}`$ (3)
$`\times \left(1+\epsilon (t^{})\mathrm{ln}(tt^{})+\mathrm{}\right)dt^{}`$ (4)
$`D_{b,t}^df(t)=(1)^n\left({\displaystyle \frac{d}{dt}}\right)^n{\displaystyle _t^b}{\displaystyle \frac{f(t^{})}{\mathrm{\Gamma }(nd(t^{}))(t^{}t)}}`$ (5)
$`\times \left(1+\epsilon (t^{})\mathrm{ln}(t^{}t)+\mathrm{}\right)dt^{}`$ (6)
Taking into account that all functions here are real functions and $`1/t=P(1/t)\pm \pi i\delta (t)`$, singular integrals here can be defined through the rule
$$_0^t\frac{f(t^{})}{tt^{}}𝑑t^{}=af(t)$$
(7)
where $`a`$ is a real regularization factor. A good agreement of (3)-(5) with the exact values given by expressions (1)-(2) can be obtained at large time by fitting the value of $`a`$.
Instead of usual integrals and usual partial derivatives, in the frames of multifractional time hypothesis it is necessary to use GFD operators to describe small alteration of physical variables. These functionals reduce to ordinary integrals and derivatives if space and time dimensions are taken to be integer, and coincide with the Riemann-Liouville fractional operators if $`d_i=const`$. If fractional dimension can be represented as $`d_i=n+\epsilon _i(𝐫(t),t),|\epsilon |1`$, it is also possible to reduce GFD to ordinary derivatives of integer order. Here we show this only for the case when $`d=1\epsilon <1`$
$`D_{0+}^{1+\epsilon }f(t)`$ $`=`$ $`{\displaystyle \frac{}{t}}{\displaystyle _0^t}{\displaystyle \frac{\epsilon (\tau )f(\tau )d\tau }{\mathrm{\Gamma }(1+\epsilon (\tau ))(t\tau )^{1\epsilon (\tau )}}}`$ (8)
$``$ $`{\displaystyle \frac{}{t}}{\displaystyle _0^t}{\displaystyle \frac{\epsilon (\tau )f(\tau )d\tau }{(t\tau )^{1\epsilon (\tau )}}}`$ (9)
Though for $`\epsilon 0`$ the last integral is well defined and is real-valued, expanding it in power series in $`\epsilon `$ leads to singular integrals like (7)
$$A=_0^t\frac{\epsilon (t^{})f(t^{})}{tt^{}}𝑑t^{}$$
To regularize this integral we will consider it to be defined on the space of finite main functions $`\phi (t^{})/2\pi i`$ and take the real part of the common regularization procedure
$$A=a\epsilon (t)f(t)$$
(10)
Thus we obtain
$$D_{0+}^{1+\epsilon }f(t)=\frac{}{t}f(t)+\frac{}{t}\left[a\epsilon (𝐫(t),t)f(t)\right]$$
(11)
where $`a`$ is a regularization parameter. For the sake of independence of GFD from this constant it is useful in the following to choose $`\beta _i`$ (on which $`\epsilon `$ depends linearily) proportional to $`a^1`$. It can be shown that for large $`t`$ the exact expressions for the terms in (1)-(2) proportional to $`\epsilon `$ are very close to the approximate expression given by (11) provided a special choice for the parameter $`\alpha `$ is is made ($`t=t_0+(tt_0),tt_0t_0,\alpha \mathrm{ln}t\mathrm{ln}t_0`$)
## II Equations of physical theories in multifractal time and space
Equations describing dynamics of physical fields, particles and so on can be obtained from the principle of minimum of fractional dimensions functionals. To do this, introduce functionals of fractional dimensions of space and time $`F_\alpha (\mathrm{}|d_\alpha (𝐫)),\alpha =t,𝐫`$. These functionals are quite similar to the free energy functionals,but now it is fractional dimension (FD) that plays the role of an order parameter (see also ). Assume further that FD $`d_\alpha `$ is determined by the Lagrangian densities $`L_{\alpha ,i},(i=1,2,\mathrm{},\alpha =t,𝐫)`$ of all the fields $`\psi _{\alpha ,i}`$, describing the particles and $`\mathrm{\Phi }_{\alpha ,i}`$ describing the interactions in the point (r)
$$d_\alpha =d_\alpha [L_{\alpha ,i}(𝐫,t)]$$
(12)
Equations that govern $`d_\alpha `$ behavior can be found by minimizing this functional and lead to the Euler’s equations written down in terms of GFD defined in (1)-(2)
$$D_{+,L_{\alpha ,i}(x)}^{d_\alpha }d_\alpha D_{,x}^{d_\alpha }D_{+,L_{\alpha ,i}^{}(x)}^{d_\alpha }d_\alpha =0$$
(13)
Substitution in this equation GFD for usual derivatives and specifying the choice for $`F`$ dependence on $`d_\alpha `$ and relations between $`d_\alpha `$ and $`L_\alpha `$ (the latter can correspond to the well known quantum field theory Lagrangians) makes possible to write down the functional dependence $`F[L]`$ in the form ($`a,b,c`$ are unknown functions of $`L`$ or constants, $`L_0`$ is infinitely large density of the measure carrier $`^n`$ energy)
$`F(\mathrm{}|d_\alpha )`$ $`=`$ $`{\displaystyle }dL_\alpha \{{\displaystyle \frac{1}{2}}[a(L_\alpha ){\displaystyle \frac{d_\alpha }{L_\alpha }}]^2`$ (14)
$`+`$ $`{\displaystyle \frac{b(L_\alpha )}{2}}(L_\alpha L_{\alpha ,0})d_\alpha ^2+c(L_\alpha )d_\alpha \}`$ (15)
or
$`F(\mathrm{}|d_\alpha )`$ $`=`$ $`{\displaystyle }d^4L_\alpha \{{\displaystyle \frac{1}{2}}[a(L_\alpha ){\displaystyle \frac{d_\alpha }{L_\alpha }}]^2`$ (16)
$`+`$ $`{\displaystyle \frac{b(L_\alpha )}{2}}(L_\alpha L_{\alpha ,0})d_\alpha ^2+{\displaystyle \frac{1}{4}}c(L_\alpha )d_\alpha ^4\}`$ (17)
The equations that determine the value of fractional dimension follow from taking the variation of (14)-(16) and read
$$\frac{}{L}\left(a(L)\frac{d_{t,\alpha }}{L}\right)+b(L)(LL_0)d_\alpha +c(L)d_\alpha ^2=0$$
(18)
or
$$\frac{}{L_\alpha }\left(a(L_\alpha )\frac{d_{t,\alpha }}{L_\alpha }\right)+b(L_\alpha )(L_\alpha L_{0,\alpha })d_\alpha ^2+c(L_\alpha )d_\alpha ^4=0$$
(19)
For nonstationary processes one have to substitute the time derivative of $`d_\alpha `$ into the right-hand side of Eqs.(18)-(19). Neglecting the diffusion of $`d_\alpha `$ processes in the space with energy densities given by the Lagrangians $`L`$ we can define $`L_\alpha L_{\alpha ,0}=\stackrel{~}{L}_\alpha L_{\alpha ,0}`$ with $`\stackrel{~}{L}_\alpha `$ having sense of over vacuum energy density and for the simplest case (18) gives ($`\alpha =t,L_{t,i}L_i`$)
$$d_t=\stackrel{~}{L}_t=1+\underset{i}{}\beta _iL_i(t,𝐫,\mathrm{\Phi }_i,\psi _i)$$
(20)
More complicated dependencies of $`d_\alpha `$ on $`L_{\alpha ,i}`$ are considered in . Note that relation (20) (and similar expression for $`d_𝐫`$ does not contain any limitations on the value of $`\beta _iL_i(t,𝐫,\mathrm{\Phi }_i,\psi _i)`$ unless such limitations are imposed on the corresponding Lagrangians, and therefore $`d_t`$ can reach any whatever high or small value.
The principle of fractal dimension minimum, consisting in the requirement for $`F_\alpha `$ variations to vanish under variation with respect to any field, in this theory produce the principle of energy minimum (for any type of fractional dimension dependency on the Lagrangian densities). It allows to receive Euler’s-like equations with generalized fractional derivatives for functions $`f(y(x),y^{}(x))`$, that describe behaviour of physical value $`f`$ depending on physical variables $`y`$ and their generalized fractional derivatives $`y^{}=D_{+,x}^{d_\alpha }f`$
$$\delta F_{t,y_i}\delta d_{t,y_i}=0$$
(21)
$$\delta _{y_i}d_\alpha (f)=\delta _{y_i}L_{\alpha ,i}(f)=0,\alpha =𝐫,t$$
(22)
$$D_{+,y_i(x)}^{d_\alpha }fD_{,x}^{d_\alpha }D_{+,y_i^{}(x)}^{d_\alpha }f=0$$
(23)
The boundary conditions will have the form
$$D_{+,y_i^{}(x)}^{d_\alpha }f|_{x_0}^{x_1}=0$$
(24)
In these equations the variables $`x`$ stand for either $`t`$ or $`𝐫`$ (the latter takes into account fractality of spatial dimensions), $`y_i=\{\mathrm{\Phi }_i,\psi _i\},(i=1,2,\mathrm{}),L_{\alpha ,i}`$ are the Lagrangian densities of the fields and particles. Here $`f`$ can be of any mathematical nature (scalar, vector, tensor, spinor, etc.), and modification of these equations for functions $`f`$ of more complicated structure does not encounter any principal difficulties. As Lagrangians $`L_{\alpha ,i}`$ one can choose any of the known in the theoretical physics Lagrangians of fields and their sums, taking into account interactions between different fields.
From Eq.(22) it is possible to obtain generalizations of all known equations of physics (Newton, Shroedinger, Dirac, Einstein equations and etc.), and the similar equations for fractional space dimensions ($`\alpha =𝐫`$). Such generalized equations extend the application of the corresponding theories for the cases when time and space are defined on multifractal sets, i.e. these equations would describe dynamics of physical values in the time and the space with fractional dimensions.The Minkowski-like space-time with fractional (fractal) dimensions for the case $`d_t1`$ can be defined on the flat continuous Minkowski space-time (that is, the measure carrier is the Minkowski space-time $`^4`$). These equations can be reduced to the well known equations of the physical theories for small energy densities, or, which is the same, for small forces ($`d_t1`$) if we neglect the corrections arising due to fractality of space and time dimensions (a number of examples from classical and quantum mechanics and general theory of relativity were considered in ).For statistical systems of many classical particles the GFD help to describe an influence of fractal structures arising in systems on behavior of distribution functions.
## III Generalized Newton Equations
Below we write down the modified Newton equations generated by the multifractal time field in the presence of gravitational forces only
$$D_{,t}^{d_t(r,t)}D_{+,t}^{d_t(r,t)}𝐫(t)=D_{+,r}^{d_r}\mathrm{\Phi }_g(𝐫(t))$$
(25)
$$D_{,r}^{d_r}D_{+,r}^{d_r}\mathrm{\Phi }_g(𝐫(t))+\frac{b_g^2}{2}\mathrm{\Phi }_g(𝐫(t))=\kappa $$
(26)
In (26) the constant $`b_g^1`$ is of order of the size of the Universe and is introduced to extend the class of functions on which generalized fractional derivatives concept is applicable. These equations do not hold in closed systems because of the fractality of spatial dimensions, and therefore we approximate fractional derivatives as $`D_{0+}^{d_𝐫}`$. The equations complementary to (25)-(26) will be given in the next paragraph. Now we can determine $`d_t`$ for the distances much larger than gravitational radius $`r_0`$ (for the problem of a body’s motion in the field of spherical-symmetric gravitating center) as (see (14) and for more details)
$$d_t1+\beta _g\mathrm{\Phi }_g$$
(27)
Neglecting the fractality of spatial dimensions and the contribution from the term with $`b_g^1`$, and taking $`\beta _g=2c^2`$), from the energy conservation law (approximate since our theory and mathematical apparatus apply only to open systems) we obtain
$`\left[1{\displaystyle \frac{2\gamma M}{c^2r}}\right]\left({\displaystyle \frac{r(t)}{t}}\right)^2`$ (28)
$`+\left[1{\displaystyle \frac{2\gamma M}{c^2r}}\right]r^2\left({\displaystyle \frac{\phi (t)}{t}}\right)^2{\displaystyle \frac{2mc^2}{r}}=2E`$ (29)
Here we used the approximate relation between generalized fractional derivative an usual integer-order derivative (13) with $`a=0.5`$ and notations corresponding to the conventional description of motion of mass $`m`$ near gravitating center $`M`$. The value $`a=0.5`$ follows from the regularization method used and alters if we change the latter. Eq.(28) differs from the corresponding equation in general theory of relativity by presence of additional term in the first square brackets. This term describes velocity alteration during gyration and is negligible while perihelium gyration calculations. If we are to neglect it, Eq.(28) reduces to the corresponding classical limit of equations of general relativity equation. For large energy densities (e.g., gravitational field at $`r<r_0`$) Eqs.(10) contain no divergences , since integrodifferential operators of generalized fractional diferentiation reduce to generalized fractional integrals (see (1)).
Note, that choosing for fractional dimension $`d_𝐫`$ in GFD $`D_{0+}^{d_𝐫}`$ Lagrangian dependence in the form $`L_{𝐫,i}L_{t,i}`$ gives for (22) additional factor of $`0.5`$ in square brackets in (28) and it can be compensating by fitting factor $`\beta _g`$.
## IV Fields arising due to the fractality of spatial dimensions (”temporal” fields)
If we are to take into account the fractality of spatial dimensions ($`d_x1,d_y1,d_z1`$), Eqs.(22)-(24), we arrive to a new class of equations describing certain physical fields (we shall call them ”temporal” fields) generated by the space with fractional dimensions. These equations are quite similar to the corresponding equations that appear due to fractality of time dimension and were given earlier. In Eqs.(13)-(16) we must take $`x=𝐫,/;\alpha =𝐫`$ and fractal dimensions $`d𝐫(t(𝐫),𝐫)`$ will obey (19) with $`t`$ being replaced by $`𝐫`$. For example, for time $`t(𝐫(t),t))`$ and potentials $`\mathrm{\Phi }_g(t(𝐫),𝐫)`$ and $`\mathrm{\Phi }_e(t(𝐫),𝐫)`$ (analogues of the gravitational and electric fields) the equations analogous to Newton’s will read (here spatial coordinates play the role of time)
$$D_{,𝐫}^{d_𝐫(𝐫,t)}D_{+,𝐫}^{d_𝐫(𝐫,t)}t(𝐫)=D_{+,t}^{d_t}\left(\mathrm{\Phi }_g(t(𝐫))+e_rm_r^1\mathrm{\Phi }_e(t(𝐫))\right)$$
(30)
$$D_{,t}^{d_t}D_{+,t}^{d_t}\mathrm{\Phi }_g(t(𝐫))+\frac{b_{gt}^2}{2}\mathrm{\Phi }_g(t(𝐫))=\kappa _r$$
(31)
$$D_{,t}^{d_t}D_{+,t}^{d_t}\mathrm{\Phi }_e(t(𝐫))+\frac{b_{et}^2}{2}\mathrm{\Phi }_e(t(𝐫))=e_r$$
(32)
These equations should be solved together with the generalized Newton equations (25)-(26) for $`𝐫(t(𝐫),𝐫)`$.
With the general algorithm proposed above, it is easy to obtain generalized equations for any physical theory in terms of GFD. From these considerations it also follows that for every physical field originating from the time with fractional dimensions there is the corresponding field arising due to the fractional dimension of space. These new fields were referred to as ”temporal fields” and obey Eqs.(19)-(21) with $`x=𝐫,\alpha =𝐫`$. Then the question arises, do these equations have any physical sense or can these new fields be discovered in certain experiments? I wont to pay attention on the next fact: if $`L_{t,i}L_{𝐫,i}`$ no new fields are generated. This is the case when fractal dimension of time and space $`d_{t,/br}`$ can not be divided on $`d_t+d_{\underset{¯}{\text{r}}}`$, the time and space fractal sets can not be divided too.The FD time and space are common and defined by value given by $`L_i`$ (the latter can be chosen in the form of usual Lagrangians in the known theories).
## V Can repulsive gravitational forces exist?
In general theory of relativity no repulsive gravitational forces are possible without a change of the Riemann space curvature (metric tensor changes). But in the frames of multifractal time and space model, even when we can neglect the fractality of spatial coordinates, from (30) it follows (for spherically-symmetric mass and electric charge distributions)
$$m_r\frac{^2t(𝐫)}{𝐫^2}=\frac{}{t}\underset{i}{}\mathrm{\Phi }_i(t)\frac{m_rk_r}{c^2t^2}\pm \frac{e^2}{ct^2}$$
(33)
with accuracy of the order of $`b^2`$. Here $`m_r`$ is the analogue of mass in the time space and corresponds to spatial inertia of object alteration with time changing (it is possible that $`m_r`$ coincides with ordinary mass up to a dimensional factor). Eq.(33) describes the change of the time flow velocity from space point to space point depending on the ”temporal” forces and indicates that in the presence of physical fields time does not flow uniformly in different regions of space, i.e the time flow is irregular and heterogeneous (see also Chapter 5 in ). Note, that introducing equations like (33) in the time space is connected with the following from our model consequences (see (22)-(24)) about equivalence of time and space and the possibility to describe properties of time (a real field generating all the other fields except ”temporal”) by the methods used to describe the characteristics of space. Below we will show that taking into consideration usual gravitational field in the presence of its ”temporal” analogues gives way to the existence of gravitational repulsion proportional to the third power of velocity. Indeed, the first term in the right-hand side of (33) is the analogue of gravity in the space of time (”temporal” field). Neglecting fractional corrections to the dimensions and taking into account both usual and ”temporal” gravitation, Newton equations have the form
$$m\frac{d^2𝐫}{dt^2}=𝐅_𝐫+𝐅_t=\frac{\gamma mM}{𝐫^2}+\frac{m_rk_r}{ct^2}\left(\frac{d𝐫}{dt}\right)^3$$
(34)
The criteria for the velocity, dividing the regions of attraction and repulsion reads
$$\left(\frac{d𝐫}{dt}\right)^3=\left(\frac{\gamma mM}{\underset{¯}{\text{r}}^2}\right)^1\frac{m_rk_r}{ct^2}$$
(35)
Here $`𝐫(t)`$ must also satisfy Eq.(25). Introducing gravitational radii $`r_0`$ and $`t_0`$ (the latter is the ”temporal” gravitational radius, similar to the conventional radius $`r_0`$), we can rewrite (28) as follows
$$\left|\frac{d𝐫}{dt}\right|=c\sqrt[3]{c\frac{t^2}{r^2}\frac{t_0}{r_0}}$$
(36)
In the last two expressions $`r`$ is the distance from a body with mass $`m`$ to the gravitating center, $`t`$ is the time difference between the points where the body and the gravitating center are situated, $`m_r=m/c,\kappa _r=\kappa _t/c`$. If we admit that $`r_0`$ and $`t_0`$ are related to each other as $`r_0=t_0/c`$, the necessary condition for the dominance of gravitation repulsion will be $`c<rt^1`$. It is not clear whether this criteria is only a formal consequence of the theory or it has something to do with reality and gravitational repulsion does exist in nature. What is doubtless, that in the frames of multifractal theory of time and space it is possible to introduce (though, may be, only formally) dynamic gravitational forces of repulsion (as well as repulsive forces of any other nature, including nuclear).
## VI The geometrization of all physical fields and forces
The multifractal model of time and space allows to consider the fractional dimensions of time $`d_t`$ and space $`d_𝐫`$ (or undivided FD $`d_{t𝐫}`$ as the source of all physical fields (see (14)) (including, in particular, the case when flat (not fractal) Minkowski space-time $`^4`$ is chosen as the measure carrier). From this point of view, all physical fields are consequences of fractionality (fractality) of time and space dimensions. So all the physical fields and forces are exist in considered model of multifractal geometry of time and space as far as the multifractal fields of time and space are exists. Within this point of view, all physical fields are real as far as our model of real multifractal fields of time and space correctly predicts and describes the physical reality. But since in this model all the fields are determined by the value of fractal dimension of time and space, they appear as geometrical characteristics of time and space (16-19). Therefore there exists a complete geometrization of all physical fields, based on the idea of time and space with (multi)fractional dimensions, the hypothesis about minimum of functional of fractal dimensions and GFD calculus used in this model. The origin of all physical fields is the result and consequence of the appearing of the fractional dimensions of time and space. One can say that a complete geometrization of all the fields that takes place in our model of fractal time and space is the consequence of the inducing (and describing by GFD) composed structure of multifractal time and space as the multifractal sets of multifractal subsets $`S_t`$ and $`S_{\underset{¯}{\text{r}}}`$ with global and local FD. The fractionality of spatial dimensions $`d_𝐫`$ also leads to a new class of fields and forces (see (22)-(24) with $`\alpha =𝐫`$). For the special case of integer-valued dimensions ($`d_t=1,d_𝐫=3`$) the multifractal sets of time and space $`S_t`$ and $`S_𝐫`$ coincide with the measure carrier $`^4`$. From (19) it follows then that neither particles nor fields exist in such a world. Thus the four-dimensional Minkowski space becomes an ideal physical vacuum (for FD $`d_\alpha >1`$ the exponent of $`R^n`$ has value $`n>4`$). On this vacuum, the multifractal sets of time and space ( $`S_t`$ and $`S_{\underset{¯}{\text{r}}}`$) are defined with their fractional dimensions, and it generates our world with the physical forces and particles.
Now the following question can be asked: what is the reason for the dependence in the considered model of fractal theory of time and space of fractionality of dimensions on Lagrangian densities? One of the simplest hypothesis seems to assume that the appearing of fractional parts in the time and space dimensions with dependence on Lagrangian densities originates from certain deformations or strains in the spatial and time sets of the measure carrier caused by the influence of the real time field on the real space field and vise versa (generating of physical fields caused by deformations of complex manifolds defined in twistor space is well known ). Assuming then that multifractal sets $`S_t`$ and $`S_𝐫`$ are complex manifolds (complex-valued dimensions of time and spatial points can be compacted), deformation, for example, of complex-valued set $`S_t`$ under the influence of the spatial points set $`S_𝐫`$ would result in appearing of spatial energy densities in time dimension, that is generating of physical fields (see ). Fractional dimensions of space appearing (under the influence of set $`S_t`$ deformations) yields new class of fields and forces (or can also not yield). It can be shown also that for small forces (e.g., for gravity - at distances much larger then gravitational radius) generalized fractional derivatives (1)-(2) can be approximated through covariant derivatives in the effective Riemann space and covariant derivatives of the space of the standard model in elementary particles theory (with the corrections taking into account fields generating and characterizing the openness of the world in whole ). All this allows to speak about natural insertion of the offered mathematical tools of GFD, at least for $`\epsilon 1`$, in the structure of all modern physical theories (note here, that the theory of gravitation as the theory of real fields with a spin $`2`$ is invented in ). Note also, that number of problems within the framework of the theory of multifractal time and space (classical mechanics, nonrelativistic and relativistic quantum mechanics) were considered in .
## VII Concluions
In our model we postulate the existence of multifractal space and time and treat vacuum as $`^n`$ space which is the measure carrier for the sets of multifractal time and space. Fractionality of time dimension leads then to appearing of space-time energy densities $`L(𝐫(t),t)`$, that is generating of the known fields and forces, and fractionality of space dimensions gives new time-space energy densities $`L(t(𝐫),𝐫)`$ and a new class of ”temporal” fields. Note, that the roles of $`d_t`$ and $`d_𝐫`$ in distorting accordingly space and time dimensions is relative and can be interchanged. Apparently, one can consider the ”united” dimension $`d_{t,𝐫}`$ \- the dimension of undivided onto time and space multifractal continuum in which time and coordinates are related to each other by relations like those for Minkowski space, not using the approximate relation utilized in this paper $`d_tt,𝐫=d_t+d_𝐫`$. Moreover in some cases it seems to be even impossible to separate space and time variables, and then $`d_t`$ and $`d_𝐫`$ can be chosen to be equal to each other, i.e., there would be only one fractional dimension $`d_t=d_{r_{x,y,z}}=1+\beta _iL_i(𝐫(t),t;t(𝐫),𝐫)`$ describing the whole space-time. In this case one would have to calculate generalized fractional derivatives from the same Lagrangians, and new ”temporal” fields will not be generated.
The considered model of multifractal time and space offers a new look (both in mathematical and philosophical senses) onto the properties of space and time and their description and onto the nature of all the fields they generate. This gives way to many interesting results and conclusions, and detailed discussion of several problems can be found in . Here we restrict ourselves with only brief enumerating of the most important ones.
a) The model does not contradict to the existing physical theories. Moreover, it reduces to them when the potentials and fields are small enough, and gives new predictions (free of divergencies in most cases) for not small fields. Though, the question about applicability of the proposed relation between fractal dimension and Lagrangian densities still remains open.
b) We consider time and space to be material fields which are the basis of our material Universe. In such a Universe there exist absolute frames of reference, and all the conservation laws are only good approximations valid for fields and forces of low energy density since the Universe is an open system, defined on certain measure carrier (the latter probably being the 4-dimensional Minkowski space). Smallness of fractional corrections to the value of time dimension in many cases (e.g., on the Earth’s surface it is about $`d_t110^{12}`$) makes possible to neglect it and use conventional models of the physics of closed systems.
c) The model allows to consider all fields and forces of the real world as a result of the geometrization of time and space (may be more convenient the term ”fractalization” of time and space) in the terms of fractal geometry. It is fractional dimensions of time and space that generate all fields and forces that exist in the world. The model introduces a new class of physical fields (”temporal” fields), which originates from the fractionality of dimensions of space. These fields are analogous to the known physical fields and forces and can arise or not arise depending on certain conditions. Thus the presented model of time and space is the first theory that includes all forces in single theory in the frames of fractal geometry. Repeat once more: the model allows to consider all the fields and forces of the world as the result of geometrization including them in FD of time and space. It is non-integer dimensions of time and space that produce the all observable fields. The new class of fields naturally comes into consideration, originating solely from the fractionality of space dimensions and with the equations similar to those of the usual fields. The presented model of space and time is the first theory that allows to consider all physical fields and forces in terms of a unique approach.
d) Basing on the multifractal model of time and space, one can develop a theory of ”almost inertial” systems which reduces to the special theory of relativity when we neglect the fractional corrections to the time dimension. In such ”almost inertial” frames of reference motion of particles with any velocity becomes possible.
e) On the grounds of the considered fractal theory of time and space very natural but very strong conclusion can be drawn: all the theory of modern physics is valid only for weak fields and forces, i.e. in the domain where fractional dimension is almost integer with fractional corrections being negligibly small.
f) The problem of choosing the proper forms of deformation that would define appearing of fractional dimensions also remains to be solved. So far there is no clear understanding now which type of fractal dimensions we must use, $`d_t`$ and $`d_𝐫`$ or $`d_{t,𝐫}`$. Obviously, solving numerous different problems will depend on this choice as the result of different points of view on the nature of multifractal structures of time and space.
The author hopes that new ideas and mathematical tools presented in this paper will be a good first step on the way of investigations of fractal characteristics of time and space in our Universe. |
warning/0001/cond-mat0001374.html | ar5iv | text | # Magnetic excitations in coupled Haldane spin chains near the quantum critical point
## I Introduction
After two decades of intensive theoretical and experimental studies, 1-dimensional (1D) Heisenberg quantum antiferromagnets (AF) are now rather well understood. A great deal of work has been done on integer-spin systems that have a spin-liquid ground state and a famous Haldane gap in the magnetic excitation spectrum. The focus in quantum magnetism has now shifted towards studies of more complex phenomena, that include inter-chain interactions, spin-lattice coupling, and/or spin-vacancies and substitutions. Of particular current interest is the quantum phase transition between spin-liquid (non-magnetic) and ordered states. This type of transition in gapped 1D systems occurs as 3D magnetic interactions and/or magnetic anisotropy are increased beyond certain threshold values. Their effect is to lower the energy of excitations at certain points in reciprocal space, and ultimately induce to a soft-mode transition to a Néel-like ordered structure. An example of such behavior is found in the extensively studied CsNiCl<sub>3</sub> compound. The corresponding phase diagram has been worked out by several authors, including Sakai and Takahashi (Fig. 1).
The most direct way to observe such a transition experimentally is by looking at a series of isostructural compounds with slightly different inter-chain coupling constants or anisotropy terms. The problem is that most known quasi-1D $`S=1`$ AF materials are otherwise deep inside the spin-liquid area of the phase diagram (good 1-D systems), or obviously in the 3D Néel-like or XY-like ordered phases (Fig. 1). CsNiCl<sub>3</sub> and related compounds are perhaps the only systems close to the phase boundary that have been extensively studied to date. Unfortunately, these compounds order in 3 dimensions at low temperature, and have no isostructural counterpart with a spin-liquid ground state. Only about a year ago the first quasi-1-D integer-spin AF that is still in the spin-liquid state, but is on the verge of a 3D ordering instability, was characterized. This Haldane-gap material, PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is so close to the phase boundary that LRO, absent in the pure compound, can be induced by spin-vacancy substitution. Moreover, an isostructural undoped system, namely SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>does order in 3-dimensions at low temperatures. For the first time two very similar stoicheometric quasi-1D materials with such vastly different ground state properties can be compared in experimental studies.
The magnetic properties of both PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> are due to spin $`S=1`$ octahedrally-coordinated Ni<sup>2+</sup> ions, while the V<sup>5+</sup> sites are presumed to be non-magnetic. The crystal structure, visualized in Fig. 2a, is tetragonal, space group $`I41cd`$, with lattice constants $`a=12.249(3)`$ Å, $`c=8.354(2)`$ Å for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, and $`a=12.1617`$ Å, $`c=8.1617`$ Å for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, respectively. The magnetic Ni-sites are arranged in peculiar spiral-shaped chains that run along the unique crystal axis, as shown in Fig. 2b. Even though all nearest neighbor Ni-Ni bonds are crystallographically equivalent, the spin-spirals have a step-4 periodicity. The dominant magnetic interaction is antiferromagnetic, between nearest-neighbor Ni<sup>2+</sup> spins within each chain. The corresponding exchange constant, $`J8.2`$ meV in both systems, was deduced from the high-temperature part of the experimental $`\chi (T)`$ curves. The ground state of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is a Haldane singlet, and the excitation spectrum has an energy gap, as unambiguously shown in low-temperature $`\chi (T)`$ and $`C(T)`$ measurements. The energy gaps, 1.2 meV and 2.2 meV, for excitations polarized along, and perpendicular to the chain-axis, respectively, were accurately determined in high-field magnetization studies. The observed anisotropy of the spin gap is attributed to single-ion easy-axis magnetic anisotropy on the Ni-sites $`D0.23`$ meV.
Unlike PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> orders magnetically in three dimensions at $`T_\mathrm{N}=7`$ K. The magnetic structure has not been determined to date, but, according to bulk measurements, is of a weak-ferromagnetic type, with a dominant antiferromagnetic component. The ordered staggered moment is along the unique $`c`$ crystallographic axis. The weak-ferromagnet distortion of this Néel (collinear) spin arrangement is attributed to the presence of weak Dzyaloshinskii-Moriya off-diagonal exchange interactions in the non-centric crystal.
Preliminary inelastic neutron scattering studies provided an estimate for the inter-chain interaction strength. Only a limited amount of neutron data are available for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, and, to date, none for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>. The present paper deals with more extensive comparative inelastic neutron scattering studies of both materials. Our results reveal the subtle differences between the two systems, responsible for their vastly distinct ground state properties.
## II Magnetic interactions
Prior to reporting our new experimental findings, we shall briefly discuss the magnetic interactions that may play an important role in the physics of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, and construct a model spin Hamiltonian for these systems. In doing so, we shall introduce the notation used throughout the rest of the paper.
As mentioned in the introduction, all previous studies point to that the dominant magnetic interaction is the antiferromagnetic coupling $`J`$ between nearest-neighbor spins (2.78 Å) within each spiral-shaped chain. It is also clear that taking the expectedly weaker inter-chain interactions into account is crucial to understanding the static and dynamic properties. In our previous work (Ref. ) we assumed that the dominant inter-chain coupling is between pairs of Ni<sup>2+</sup> ions from adjacent chains, such that both interacting spins have the same $`c`$-axis fractional cell coordinate, and that each Ni-site is coupled to 4 other sites. The resulting coupling topology is schematized in Fig. 4a. Considering only these inter-chain bonds, however, can not result in a very realistic model. First of all, they do not correspond to the shortest inter-chain Ni-Ni distance. More importantly, in the crystal structure of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, there seem to be no obvious superexchange pathways corresponding to these links. If one pays close attention to crystal symmetry, bond length, and possible superexchange routes involving the VO<sub>4</sub> tetrahedra, one arrives at a more complex interaction geometry, visualized in Fig. 3. Here one assumes interactions along the shortest inter-chain Ni-Ni bond (5.0 Å). Coupled Ni<sup>2+</sup> ions are offset relative to each other by $`c/4`$ along the chain axis, and are bridged by a VO<sub>4</sub> tetrahedron. All such bonds are crystallographically equivalent. Each Ni-site is linked with only two adjacent chains. The overall coupling topology for this model is as shown in Fig. 4b. We shall denote the corresponding exchange constant as $`J_1`$.
Magnetic anisotropy is clearly manifest in bulk susceptibility and magnetization measurements, and must be explicitly included in the spin Hamiltonian. Because of the strong dispersion along the chain-axis, all the low-energy response of each chain is concentrated at wave vectors close to the 1-D AF zone-center. In this long-wavelength limit, 2-ion anisotropy of in-chain interactions and single-ion anisotropy associated with individual spins can not be distinguished. In our model we shall therefore include only the latter, and write the corresponding term in the Hamiltonian as:
$$\widehat{H}_{\mathrm{single}\mathrm{ion}}=D\underset{i,k}{}(S_{i,k}^z)^2.$$
(1)
The choice of the anisotropy axis along the chain direction is based on previous bulk magnetic studies. Here $`D`$ is the anisotropy constant ($`D<0`$ = easy-axis) and $`𝑺_{i,k}`$ is the spin operator for site $`i`$ in chain $`k`$. At the same time we shall assume in-chain exchange interactions to be isotropic and write this Heisenberg term as:
$$\widehat{H}_{\mathrm{in}\mathrm{chain}}=J\underset{i,k}{}𝑺_{i,k}𝑺_{i+1,k}.$$
(2)
The magnitude of dispersion perpendicular to the chain axis is expected to be rather small, so the entire range of wave vector transfers in the $`(a,b)`$ plabe will influence the low-energy properties. For this reason, two-ion anisotropy of $`J_1`$, unlike that of $`J`$, is a relevant parameter and should be considered. The inter-chain coupling term in the Hamiltonin will thus have the form:
$$\widehat{H}_{\mathrm{inter}\mathrm{chain}}=\underset{i,i^{},k,k^{}}{}\left[J_{1,}S_{i,k}^zS_{i^{},k^{}}^z+J_{1,}\left\{S_{i,k}^xS_{i^{},k^{}}^x+S_{i,k}^yS_{i^{},k^{}}^y\right\}\right].$$
(3)
Here the sum is taken over pairs of next-nearest-neighbor spins. The three terms discussed above constitute the spin Hamiltonian that we propose for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>. Of course, other magnetic interactions can be active in the system as well. As will be discussed below, however, the described model can reproduce the experimental data with a minimal number of parameters.
## III Experimental
Inelastic neutron scattering studies of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> powder samples (about 10 g each) were carried out in two series of experiments. The conventional 3-axis technique was used to measure certain characteristic constant-$`E`$ and constant-$`Q`$ scans. These measurements were performed at the TASP and Druchal spectrometers at the continuous spallation source SINQ at Paul Scherrer Institute. Neutrons of a fixed final energy of 5 meV or 8 meV were used with pyrolytic graphite (PG) monochromator and analyzer, and a typical $`\text{(open)}80^{}80^{}\text{(open)}^{}`$ collimation setup. For some scans a horizontally focusing analyzer was used to increase the useful scattered intensity, in which case no collimators were used after the sample. In most of the measurements a flat analyzer was employed instead. In the 5 meV-final configuration a Be filter was inserted after the sample to suppress higher-order beam contamination. The spectrum of neutrons emerging from the guide is such that no higher-order filter was required for the 8 meV-final measurements.
In the second series of experiments we took advantage of the area-sensitive detector (ASD) 3-axis setup available for the NG-5 “SPINS” spectrometer installed at the National Institute of Standards and Technology. In special cases when a large domain of $`EQ`$ space needs to be surveyed in a spherically symmetric sample (such as powder) this technique can provide an amazing, almost 10-fold increase of data collection rate, as compared to the standard 3-axis geometry, without any penalty in resolution. In our experiments on PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we utilized a PG monochromator and a $`\text{(open)}80^{}80^{}80^{}(\mathrm{radial})`$ array of collimators. The measurements were done in the fixed-final-energy mode, with the central blade of the composite analyzer tuned to $`E_\mathrm{f}=3.125`$ meV. In the experiment scattering events with final neutron energies in the range $`E_\mathrm{f}\pm 0.4`$ meV are registered simultaneously. To suppress higher-order beam contamination we used a Be-O filter after the sample. The data were taken for momentum and energy transfers of up to 2.5 Å<sup>-1</sup> and 7 meV, respectively.
In both series of experiments the sample environment was a standard “ILL-orange” cryostat, and the temperature range $`1.530`$ K was covered. The background was measured by repeating some scans with the analyzer moved away from its elastic position by 10. In the SPINS experiment the background due to air scattering was separately measured with the sample removed from the spectrometer.
## IV Experimental results
### A Results obtained at low temperatures
The inelastic neutron scattering intensities measured in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> at $`T=1.2`$ K with the ASD setup are shown in the false-color plot in Fig. 5 (background subtracted). The resolution of the area-sensitive detector is greater than the actual energy and wave vector resolution of the spectrometer (0.11 meV and 0.023 Å<sup>-1</sup> at $`\mathrm{}\omega =0`$). The customary procedure is to re-bin the mesh data to a coarser grid. Instead, for visualization purposes, the data in Fig. 5 were smeared using a fixed-resolution 0.5 meV$`\times `$0.06 Å<sup>-1</sup> FWHM Gaussian filter. As reported previously, constant-energy scans measured in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> at 4 meV energy transfer have a peculiar and very characteristic shape. Such scans were collected at $`T=2`$ K for both materials, using the 8 meV-final/flat analyzer setup SINQ. Theses scans are shown in Fig. 6. In addition, representative constant-$`Q`$ scans were also measured in this configuration, and are shown in Figs. 7 and 8.
Despite the different ground states, the main features of the powder-averaged dynamic cross section for the two materials are quite similar. This, however, is not particularly surprising. The dynamic structure factor $`S(𝑸,\omega )`$ is severely smeared out by the spherical averaging that occurs in a powder sample. The main features are defined by the form factor of each spiral-shaped chain, and by the steep dispersion along the chain axis, expected to be almost identical in the two systems. The significant difference between the two materials is expected to be in the magnitude of the inter-chain interactions, and possibly single-ion anisotropy. These effects are much more difficult to observe, as they only influence the weak spin wave dispersion in the $`(a,b)`$ crystallographic plane. Our data do in fact contain relevant information on the transverse dispersion of spin excitations, which can be extracted from a quantitative analysis, as described in the following section.
### B Data analysis
#### 1 Model cross section
The general problem in interpreting inelastic neutron scattering data from powder samples is an effective loss of information upon spherical averaging. Indeed, the quantity of interest is the dynamic structure factor $`S(𝑸,\omega )`$ that for each channel of spin polarization is a scalar function in 4-dimensional $`E𝑸`$ space. In a powder sample one measures the spherical average of this function, a scalar function defined in 2-dimensional space:
$$S_{\text{powder}}(Q,\omega )=\frac{1}{4\pi }𝑑\varphi \mathrm{sin}\theta d\theta S([Q\mathrm{sin}\theta \mathrm{cos}\varphi ,Q\mathrm{sin}\theta \mathrm{sin}\varphi ,Q\mathrm{cos}\theta ],\omega )$$
(4)
The transformation $`S(𝑸,\omega )S_{\text{powder}}(Q,\omega )`$ is not reversible. The only way the full correlation function can be extracted from the experiment is by assuming some parameterized model for $`S(𝒒,\omega )`$ and fitting it to the measured $`S_{\text{powder}}(Q,\omega )`$. For PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we used the following analytical form for $`S(𝒒,\omega )`$, as derived below in the theory section:
$`S(𝑸)`$ $`=`$ $`P_{}(𝑸)S_{}(𝑸)+P_{}(𝑸)S_{}(𝑸),`$ (5)
$`2S_{}(𝑸)`$ $`=`$ $`2\mathrm{cos}^2\psi _1\mathrm{cos}^2\psi _2S_{}^{}(h,k,l)`$ (6)
$`+`$ $`\mathrm{cos}^2\psi _1\mathrm{sin}^2\psi _2\left[S_{}^{}(h+1,k,l+1)+S_{}^{}(h+1,k,l+3)\right]`$ (7)
$`+`$ $`\mathrm{sin}^2\psi _1\mathrm{cos}^2\psi _2\left[S_{}^{}(h,k+1,l+1)+S_{}^{}(h,k+1,l+3)\right]`$ (8)
$`+`$ $`2\mathrm{sin}^2\psi _1\mathrm{sin}^2\psi _2S^{}(h+1,k+1,l+2),`$ (9)
$`2S_{}(𝑸)`$ $`=`$ $`2\mathrm{cos}^2\psi _1\mathrm{cos}^2\psi _2S_{}^{}(h,k,l)`$ (10)
$`+`$ $`\mathrm{cos}^2\psi _1\mathrm{sin}^2\psi _2\left[S_{}^{}(h+1,k,l+1)+S_{}^{}(h+1,k,l+3)\right]`$ (11)
$`+`$ $`\mathrm{sin}^2\psi _1\mathrm{cos}^2\psi _2\left[S_{}^{}(h,k+1,l+1)+S_{}^{}(h,k+1,l+3)\right]`$ (12)
$`+`$ $`2\mathrm{sin}^2\psi _1\mathrm{sin}^2\psi _2S^{}(h+1,k+1,l+2),`$ (13)
$`\psi _1`$ $`=`$ $`{\displaystyle \frac{\pi d}{a}}h,`$ (14)
$`\psi _2`$ $`=`$ $`{\displaystyle \frac{\pi d}{a}}k.`$ (15)
In these formulas the argument $`\omega `$ has been dropped. The phases $`\psi _1`$ and $`\psi _2`$ represent the 3-D structure factor of the spiral-shaped spin chains, and $`d=0.08a`$ is the offset of each Ni<sup>2+</sup>-ion relative to the central axis of the corresponding spiral chain, along the $`a`$ or $`b`$ axis. The polarization factors for longitudinal (polarization along the $`c`$-axis) and transverse (polarization in the $`(a,b)`$ plane) spin excitations are defined as:
$`P_{}(𝑸)=\mathrm{sin}^2(\widehat{𝑸,𝒛}),`$ (16)
$`P_{}(𝑸)=1+\mathrm{cos}^2(\widehat{𝑸,𝒛}).`$ (17)
The dynamic structure factors $`S_{}^{}`$ and $`S_{}^{}`$ for straight (as opposed to spiral-shaped) Haldane chains are written in the Single-Mode Approximation (SMA) :
$`S_{}^{}(𝑸,\omega )={\displaystyle \frac{Zv}{\mathrm{}\omega _{}(𝑸)}}\delta (\mathrm{}\omega \mathrm{}\omega _{}(𝑸)),`$ (18)
$`S_{}^{}(𝑸,\omega )={\displaystyle \frac{Zv}{\mathrm{}\omega _{}(𝑸)}}\delta (\mathrm{}\omega \mathrm{}\omega _{}(𝑸)).`$ (19)
Here $`v=2.48J`$ is the spin wave velocity, and $`Z=1.26`$ (Ref. ). Finally, the dispersion relation for weakly coupled chains are given by:
$`\left[\mathrm{}\omega _{}(𝑸)\right]^2`$ $`=`$ $`\mathrm{\Delta }_{}^2+v^2\mathrm{sin}^2(\pi l/2){\displaystyle \frac{1}{2}}ZvJ_{1,}\mathrm{cos}(\pi l/2)\left[\mathrm{cos}(\pi h)+\mathrm{cos}(\pi k)\right](1\mathrm{cos}\pi l/2),`$ (20)
$`\left[\mathrm{}\omega _{}(𝑸)\right]^2`$ $`=`$ $`\mathrm{\Delta }_{}^2+v^2\mathrm{sin}^2(\pi l/2){\displaystyle \frac{1}{2}}ZvJ_{1,}\mathrm{cos}(\pi l/2)\left[\mathrm{cos}(\pi h)+\mathrm{cos}(\pi k)\right](1\mathrm{cos}\pi l/2),`$ (21)
where $`\mathrm{\Delta }_{}`$ and $`\mathrm{\Delta }_{}`$ are the longitudinal and transverse intrinsic Haldane gaps for non-interacting chains, respectively. The splitting of the triplet in our model is caused by single-ion anisotropy:
$`\mathrm{\Delta }_{}=\mathrm{\Delta }0.57D,`$ (22)
$`\mathrm{\Delta }_{}=\mathrm{\Delta }+1.41D.`$ (23)
An important parameter is the mean intrinsic Haldane gap $`\mathrm{\Delta }(\mathrm{\Delta }_{}+2\mathrm{\Delta }_{})/3`$. As can be seen from Eqs. 22 and 23, it is, to the first order, defined by $`J`$ alone: $`\mathrm{\Delta }0.41J`$. It is also useful to define the actual gaps (excitation energies at the 3-D AF zone-center):
$`E_{\text{min},}^2=\mathrm{\Delta }_{}^22Zv|J_{1,}|,`$ (24)
$`E_{\text{min},}^2=\mathrm{\Delta }_{}^22Zv|J_{1,}|.`$ (25)
The spherical average of Eq. 5 was calculated numerically using a Monte-Carlo algorithm that also eliminates the $`\delta `$-functions in Eqs. 18 and 19. The parameters were then refined by a standard least-squares routine to best-fit the data. To accelerate the fitting process the ASD data were binned to a rectangular 0.046 Å$`{}_{}{}^{1}\times `$ 0.1 meV resolution. An additional benefit of this binning is that it allowed us not to worry about resolution effects. The lower 0.5 meV energy transfer range, that contains the elastic-incoherent and possibly phonon scattering, was excluded from the fits.
#### 2 Analysis of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> data.
For The Pb-compound the adjustable parameters were the two intrinsic gap energies $`\mathrm{\Delta }_{}`$ and $`\mathrm{\Delta }_{}`$, and the doublet 3D gap $`E_{\text{min},}`$. The singlet 3-D gap was fixed to $`E_{\text{min},}=1.2`$ meV, as determined in high-field bulk measurements. The only additional parameter was an overall scale factor. The least-squares refinement yields: $`\mathrm{\Delta }_{}=4.0\pm 0.25`$ meV, $`\mathrm{\Delta }_{}=3.1\pm 0.3`$ meV and $`E_{}^{(\mathrm{min})}=2.4\pm 0.2`$ meV ($`J_1<0`$), with $`\chi ^2=2.3`$. The resulting fit is shown in Fig. 9. Substituting the obtained values values into the expression for $`S(𝑸,\omega )`$, performing a powder average and convoluting the result with the spectrometer resolution function, reproduces the measured 3-axis const-$`E`$ and const-$`Q`$ scans rather well, as shown in solid lines in Figs. 6b,7. The relatively large $`\chi ^2`$ of the global fit to the ASD data is to be attributed to systematic error, primarily due to “spurious” scattering. In particular, the data are partially contaminated by spurions of type “$`k_ik_i`$” and “$`k_fk_f`$” originating from the stronger Bragg powder lines from the sample. Areas that are affected by these spurious processes are shown as shaded “streaks” in Fig. 5. Another prominent streak in the data is to the left from, and parallel to, the two shaded areas. The origin of this feature is not known. It is however resolution-limited and temperature-independent, and is thus almost certainly spurious and of non-magnetic origin. Considering that this type of systematic error is unavoidable in powder experiments, the obtained model fit to the data is quite acceptable.
The refined value for $`E_{\text{min},}`$ is in very good agreement with the result of high-field measurements. The interchain coupling constants are obtained from Eqs. 24,25: $`J_{1,}=0.18`$ meV and $`J_{1,}=0.14`$ meV. Note that $`|J_1|`$ is larger by a factor of 1.5–2, compared to our previous estimate ($`J_{}=0.096\pm 0.003`$ meV) in Ref. . This discrepancy should be partly attributed to a difference in the definition of $`J_1`$. In our previous model each spin was coupled to 4 spins in adjacent chains (coordination number 4). In the present model the inter-chain coordination number is 2. This automatically translates into a factor of 2 for $`|J_1|`$. The negative sign of $`J_1`$ indicates that this coupling is actually ferromagnetic. Note, however, that in our model $`J_1`$ couples spins offset by $`c/4`$ along the chain axis, so that the effective mean field coupling between interacting chains is still antiferromagnetic, as in our previous model. The 3-D magnetic zone-center, where transverse dispersion is a minimum, is at $`(1,1,2)`$. From the refined $`\mathrm{\Delta }_{}`$ and $`\mathrm{\Delta }_{}`$ we can also get the in-chain coupling constant: $`J9.0`$ meV. This value is in a better agreement with the high-temperature susceptibility estimate $`J8.2`$ meV in Ref. , than our preliminary neutron result $`J=9.5`$ meV, from the same reference. The anisotropy constant $`D`$ can be estimated from $`\mathrm{\Delta }_{}`$ and $`\mathrm{\Delta }_{}`$ using Eqs. 22,23. For PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we get $`D=0.45`$ meV.
To convince the reader that the $`J_1`$ is in fact a relevant parameter, we have performed simulations of the dynamic structure factors using the same values $`\mathrm{\Delta }_{}=4.0\pm 0.25`$ meV, $`\mathrm{\Delta }_{}=3.1\pm 0.3`$ meV, as determined in the analysis of the ASD data, but with the sign of $`J_1`$ reversed: $`J_{1,}=0.18`$ meV and $`J_{1,}=0.14`$ meV. The resulting simulated powder cross section is visualized in Fig. 10a, and is clearly very different from the inelastic intensity measured experimentally. Similarly, Fig. 10b shows a simulation with $`J_{1,}=J_{1,}=0`$ meV. $`J_1`$ has a particular impact on the shape of the constant-energy powder scans. The dashed lines in Fig. 6 are simulations for uncoupled chains. The well-defined intensity maximum seen in the data at $`|Q|1.2`$ Å<sup>-1</sup>, is replaced with a broad monotonous feature for uncoupled chains. In contrast, the interacting chain model reproduces the peak rather well.
#### 3 Analysis of SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> data
Similar data analysis was performed for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>. $`E_{}^{(\mathrm{min})}`$ in this case was fixed at zero value (see discussion in the theory section below). The best fit to the ASD data is obtained with $`\mathrm{\Delta }_{}=3.9\pm 0.3`$ meV, $`\mathrm{\Delta }_{}=2.8\pm 0.4`$ meV and $`E_{}^{(\mathrm{min})}=2.35\pm 0.3`$ meV with $`\chi ^2=2.0`$, and is shown in Fig. 9b. These values correspond to interchain coupling constants $`J_{1,}=0.18`$ meV and $`J_{1,}=0.15`$ meV, respectively, almost identical to the corresponding PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> values. An important consistency check is that using the gap values to estimate $`J`$ gives almost the same value as for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, as expected: $`J=8.6`$ meV. The easy-axis anisotropy constant in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is larger than in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>: $`D=0.56`$ meV. Simulations for constant-$`E`$ and constant-$`Q`$ scans measured for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> in the standard 3-axis mode are shown in solid lines in Figs. 6a,8.
#### 4 Observation of actual spin gaps at the 3D magnetic zone-center
From the point of view of spin dynamics, the main distinction between the singlet-ground-state PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and the magnetically ordered SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is the presence of a spin gap in the former system, and a gapless excitation spectrum in the latter. The gap is directly accessible experimentally at the 3-D AF zone-center, where the dispersion of magnetic excitations is a global minimum. As mentioned above, for both vanadates the 3D zone-center is at $`𝑸^{(0)}=(1,1,2)`$, which corresponds to a momentum transfer $`|𝑸^{(0)}|1.67`$ Å<sup>-1</sup>. Constant-$`Q`$ scans extracted from our ASD data for this momentum transfer, as well as a standard $`E_f=5`$ meV constant-$`Q`$ scan at this wave vector, are shown in Fig. 11. To better understand these data, we note that near the 3D zone-center the single-mode contribution to the dynamic structure factor, independently of the details of the spin Hamiltonian, should be of the following form:
$`S(𝑸,\omega )`$ $``$ $`{\displaystyle \frac{1}{\omega _𝑸}}\delta (\omega \omega _𝑸),`$ (26)
$`(\mathrm{}\omega _𝑸)^2`$ $`=`$ $`\mathrm{\Delta }^2+v_{}^2(Q_{}Q_{}^{(0)})^2+v_{}^2(Q_{}Q_{}^{(0)})^2.`$ (27)
In this formula $`\mathrm{\Delta }`$ is the gap energy, and $`v_{}`$ and $`v_{}`$ are spin wave velocities along and perpendicular to the chain axis $`c`$, respectively. The data shown in Fig. 11c were collected with a horizontally-focusing analyzer. At energy transfers below 1 meV in this mode we are picking up a great deal of diffuse and phonon scattering ( $`(1,1,2)`$ is an allowed nuclear Bragg peak !). Above this contaminated region though, in a powder sample we effectively observe a $`𝑸`$-integral of the cross-section around the 3D zone-center. The same applies to the scans in Figs. 11a and b, where integration was performed in the range 1.6–1.7 Å<sup>-1</sup>. For the $`Q`$-integrated intensity Eq. 26 gives:
$$S(\omega )\sqrt{\omega ^2\mathrm{\Delta }^2}.$$
(28)
The linear increase of intensity seen for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> in Fig. 11 should thus indeed be interpreted as due to a gapless spin wave ($`\mathrm{\Delta }=0`$). In contrast, for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, the threshold behavior is a clear sign of an energy gap. These effects are also seen in constant-$`Q`$ scans extracted from the ASD data by binning the pixels in a 0.1 Å<sup>-1</sup> $`Q`$-range.
### C SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>: Temperature dependence
To better understand the mechanism of long-range ordering of SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we studied the temperature dependence of inelastic scattering at the momentum transfer $`|𝑸^{(0)}|=1.67`$ Å<sup>-1</sup> in this compound. Typical data are shown in Fig. 12. The solid lines were obtained by fitting our model cross section to the data at each temperature. The lower 0.75 meV energy transfer range was excluded from this analysis, as the difficult-to-estimate phonon contribution in this range is expected to increase dramatically with increasing $`T`$. The set of independent parameters was slightly modified. The anisotropy splitting of the triplet $`D`$, as well as both inter-chain coupling constants $`J_{1,}`$ and $`J_{1,}`$ were fixed at the values determined at $`T=1.5`$ K. The mean gap $`\mathrm{\Delta }`$ and an intensity prefactor were refined to best-fit the scans at each temperature. The obtained temperature dependences are shown in Fig. 13. As previously observed in other Haldane-gap systems, $`\mathrm{\Delta }`$ increases with increasing $`T`$. Even though this change is rather small, according to Eqs. 24,25, it corresponds to an appreciable variation of the gap in the longitudinal mode (Fig. 14). Upon cooling, the longitudinal gap approaches zero at $`T=T_\mathrm{N}`$, which results in a soft-mode transition to a magnetically ordered state. This type of behavior is very similar to that found in CsNiCl<sub>3</sub>.
## V Theory and discussion
### A Derivation of the model cross section
To calculate the effect of weak inter-chain interactions on dynamic spin correlations in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we shall use the Random Phase Approximation (RPA). This approach for directly coupled Haldane spin chains has been successfully applied, for example, to CsNiCl<sub>3</sub> (Ref. ). More recently it was shown to also work well for Haldane chains coupled via classical spins. The only additional difficulties in the present case arise from the complicated 3-D arrangement of magnetic ions in the PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> structure, and from the rather non-trivial geometry of interchain coupling. To somewhat simplify the task we shall break it up into two distinct problems. First, we shall worry only about the topology of magnetic interactions and consider an equivalent Bravais lattice of spins, assuming straight spin chains and a system of inter-chain bonds shown in Fig. 4b. Making use of the general RPA equations for coupled spin chains, summarized in Appendix I, we shall write down the RPA susceptibility and dynamic structure factor $`S^{}(𝑸,\omega )`$ for this model. In a separate step we shall adapt the result to the more complex structure of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, using the formulas of Appendix II.
#### 1 Structure factor for an equivalent Bravais lattice
The first step in the RPA calculation is to write down the bare (non-interacting) susceptibility for an isolated Haldane spin chain, for which the single-mode approximation (SMA) is known to work rather well:
$$\chi _0(𝑸,\omega )=\frac{1\mathrm{cos}(Q_zc/4)}{2}\frac{Zv}{\mathrm{\Delta }^2+v^2\mathrm{sin}^2(Q_zc/4)(\mathrm{}\omega +iϵ)^2}$$
(29)
This expression should be used separately for each channel of spin polarization, with appropriate values of gap energy $`\mathrm{\Delta }`$ for each particular mode.
Expression 29 is to be substituted in the general RPA equations 40. Even for the “straightened out” spin chains in Fig. 4b there are still 16 spins per unit cell, so we end up with 16 couple equations and 16 modes for each polarization! For the low-energy part of the excitation spectrum it is however quite appropriate to use the approximation of Eqs. 41. This reduces the problem to only 4 self-consistent RPA equations (there are 4 spin chains in each crystallographic unit cell). Moreover, at this level of approximation the set of 4 RPA equations is degenerate, and we end up with a single RPA equation for each polarization:
$$\chi _{\mathrm{RPA}}^1(𝑸)=\chi _0^1(𝑸)\left[\chi _0(𝑸)𝒥(𝑸)+1\right],$$
(30)
where
$$𝒥(𝑸)=J_1\mathrm{cos}(Q_zc/4)\left[\mathrm{cos}(Q_xa/2)+\mathrm{cos}(Q_ya/2)\right].$$
(31)
This equation is easily solved analytically. By taking the imaginary part of the thus obtained $`\chi _{\mathrm{RPA}}(𝑸)`$ we arrive at Eqs. 1821.
#### 2 Actual structure
The transformation from the Bravais spin lattice (straight chains) to the spiral-chain structure of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is schematized in Fig. 15. This figure also serves as an illustration of the notation used in Appendix I. Applying Eq. 53 is straightforward, and we will spare the reader the tedious calculations. The final relation between $`S(𝒒,\omega )`$ and $`S^{}(𝒒,\omega )`$ is as in Eqs. 9,15.
### B Temperature effects
#### 1 Disordered phase
The derivation above can be repeated almost verbatim for the case $`T>0`$. Temperature enters RPA calculation indirectly, through an intrinsic temperature dependence of bare susceptibilities of individual chains (Eq. 29). Two effects at $`T>0`$ need be considered: i) the increase of the Haldane gap, compared to its $`T=0`$ value, and ii) damping of Haldane excitations. Both phenomena have been observed in a number of model quasi-1-D systems (see, for example, Refs. ). Unfortunately, it is almost impossible to extract meaningful information regarding excitation lifetimes from powder inelastic data. As most of the observed inelastic signal originates from excitations with energies greater than 2 meV, the temperature dependence of the intrinsic Haldane gap is much easier to observe. For this reason we analyzed the temperature dependence of the inelastic signal measured in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> using the same single-mode approximation (Eq. 29), as at base temperature (see fits in Fig. 12), to obtain the $`T`$-dependences of intrinsic Haldane gap energies (Fig. 13). From the mapping of the Heisenberg Hamiltonian on the quantum non-linear $`\sigma `$-model (NLSM) one expects a rather steep increase of the gap energy with increasing temperature:
$$\mathrm{\Delta }(T)\mathrm{\Delta }(0)+\sqrt{2\pi }\sqrt{T\mathrm{\Delta }(0)}\mathrm{exp}(\frac{\mathrm{\Delta }(0)}{k_\mathrm{B}T}).$$
(32)
The actual increase of the gap energy observed experimentally NENP, NINAZ, and Y<sub>2</sub>BaNiO<sub>5</sub>, was found to be consistently smaller that this NLSM prediction. The same discrepancy, namely a rather slow increase of the gap energy with temperature, is seen in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> as well. Equation 32 is plotted in a dashed line in Fig. 13a. The solid line is an empirical fit in which the prefactor $`\sqrt{2\pi }`$ in Eq. 32 was replaced by an adjustable parameter $`A`$. For SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> we get $`A=0.67(1)\sqrt{2\pi }`$. A self-consistency check for our SMA model is that the refined intensity prefactor is almost $`Tindependent`$, as shown in Fig. 13b. In other words, the decrease of actual excitation intensity is entirely due to the increase of gap energy, and the intensity scales as $`1/\omega `$.
In SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> the ordering temperature $`T_\mathrm{N}=7`$ K is significantly smaller than the intrtinsic Haldane gap energy $`\mathrm{\Delta }4`$ meV. As both the gap energy and excitation width are expected to increase exponentially with $`T`$, such a small $`T_N`$ suggests that inter-chain interactions in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> are barely strong enough to produce a LRO-ground state. Long-range ordering in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, and the absence of such in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, are a very “lucky” coincidence that results from a fine interplay between in-chain interaction strength, magnetic anisotropy and inter-chain coupling.
#### 2 Magnetically ordered state
Our model cross section was derived under the implicit assumption that the system is in a non-magnetic state, which is not applicable to SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> at $`T<T_\mathrm{N}`$. In this regime one expects an increase of the gap energies, at least for the longitudinal mode, due to the presence of a mean static staggered field generated by the ordered staggered moment in the system. To adapt the MF-RPA calculation to this temperature regime one has to know the ordered moment in the system. To date, powder diffraction experiments on SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> failed to detect any magnetic Bragg reflections in this compound at low temperatures, and the magnetic order parameter is thus expected to be very small. This is consistent with our previous observation that the system is “barely” 3-dimensional enough to become ordered at low temperatures. Provided the ordered moment is small, within the accuracy of our powder experiments it seems quite appropriate to use the same cross section for $`T<T_\mathrm{N}`$, as for the paramagnetic phase, and postulate $`E_{\mathrm{min},}=0`$. This assumption is equivalent to assuming $`T_\mathrm{N}=0`$.
#### 3 Placement on the phase diagram
The line of quantum phase transition separating the ordered and spin-liquid states can be derived from Eqs. 22-25. The critical value for $`J_1`$ corresponds to the lower gap energy (in our case $`E_{\mathrm{min},}`$) being equal to zero. The resulting phase diagram is very similar to the direct numerical calculations by Sakai and Takashi, and is shown in Fig. 1. Using the results for PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> and SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> described above, we are able to place the two new materials on the same plot. It has to be emphasized, that for SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> the parameters $`D`$ and $`J_1`$ were extracted from the experimental data using a cross-section for a disordered system, and assuming $`E_{\mathrm{min},}0`$. That SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> lands exactly on the phase boundary is thus an artifact of our data analysis procedure. In reality, SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> must be positioned slightly above the phase boundary. The disordered PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is not influenced by such an artifact, and its positioning can be considered quite reliable.
## VI Summary
Our neutron scattering results help explain the different ground state properties of the two very similar vanadates. The ratio of inter-chain to in-chain coupling in SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is slightly larger than in PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>. SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> is also driven towards LRO by the somewhat larger magnetic anisotropy. The quantitative data analysis enables us to precisely position the two compounds on the Sakai-Takahashi phase diagram, just opposite to each other relative to the ordered-disordered phase boundary. This almost unbelievable coincidence opens many exciting possibilities for future studies. Experiments on aligned powders will enable more accurate measurements of the excitation spectrum near the 3D zone-center. High-pressure studies, as previously attempted for NENP, may lead to the first observation of pressure-induced long-range ordering in a quantum-disordered magnet.
###### Acknowledgements.
We would like to thank Y. Sasago, who has suggested at the early stage of this study that SrNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub> may be a Haldane-gap antiferromagnet. We also thank R. Wickmann for sending a copy of Ref. . Work at the University of Tokyo was supported in part by Grant-in-Aid for COE Research “Phase Control of Spin-Charge-Phonon Coupled Systems” from the Ministry of Education, Science, Sports and Culture of Japan. Work at Brookhaven National Laboratory was carried out under Contract No. DE-AC02-76CH00016, Division of Material Science, U.S. Department of Energy. Studies at NIST were partially supported by the NSF under contract No. DMR-9413101.
## Appendix I: RPA susceptibility for weakly coupled spin chains
The general RPA technique of calculating the dynamic susceptibility and structure factor for weakly interacting systems is well established and documented (see, for example, Ref. ). In particular, the approach has been on many occasions applied to weakly coupled quantum spin chains. In the present section, without any claims of novelty, and for reference only, we shall derive some useful RPA results for weakly coupled spin chains assuming a rather general geometry of inter-chain bonding. Note that the notation used in this Appendix is totally independent from that in the rest of the paper.
We consider a crystal structure composed of identical parallel uniform spin chains that run along the $`𝒂`$ axis. The spacing between spins in each chain is $`a/K`$, where $`K`$ is integer. There are $`M`$ chains in each crystallographic unit cell. The origin of the $`m`$-th chain is at $`𝑹_n+𝝆_m`$. The position of the $`i`$-th spin in chain $`(n,m)`$ can be written as $`𝑹_n+𝝆_m+i𝒂/K`$. A single crystallographic unit cell thus contains $`MK`$ spins. For convenience we shall break up the in-chain spin index $`i`$ into two: $`i=lK+k`$, $`0k<K`$. The spin Hamiltonian for inter-chain interactions will be written as:
$$\widehat{H}_{\mathrm{interchain}}=\underset{n,n^{}}{}\underset{m,m^{}}{}\underset{l,l^{}}{}\underset{k=0}{\overset{K1}{}}\underset{k^{}=0}{\overset{K1}{}}s_{l,k}^{(n,m)}s_{k^{},l^{}}^{(n^{},m^{})}J(n,n^{},m,m^{},Kl+k,Kl^{}+k^{})$$
(33)
We assume all exchange interactions to be simultaneously diagonal in spin projection indexes. All the equations in this section thus apply to a particular channel of spin polarization. In our model the exchange constant satisfies certain translational-symmetry relations:
$$J(n,n^{},m,m^{},Kl+k,Kl^{}+k^{})=J((𝑹_n^{}𝑹_n),m,m^{},(l^{}l),k,k^{}).$$
(34)
In order to write down self-consistent RPA equations, we have to answer the following question. Suppose we have artificially induced a spin density $`s^{(m^{})}(𝒓)`$ (or, equivalently, $`s^{(m^{})}(𝒒)`$)in the chain-sublattice $`m^{}`$. What exchange field will this spin density project on the chains in sublattice $`m`$? By definition, the exchange field acting on spin $`s_{l,k}^{(n,m)}`$ is given by:
$$h_{l,k}^{(n,m)}=\underset{n^{}}{}\underset{m^{}=0}{\overset{M1}{}}\underset{l^{}}{}\underset{k^{}=0}{\overset{K1}{}}s_{l^{},k^{}}^{(n^{},m^{})}J((𝑹_n^{}𝑹_n),m,m^{}(l^{}l),k,k^{})$$
(35)
We, of course, will be interested in the Fourier transform of the exchange field acting on sublattice $`m`$:
$$h^{(m)}(𝒒)\underset{n}{}\underset{l}{}\underset{k=0}{\overset{K1}{}}h_{l,k}^{(n,m)}\mathrm{exp}\left[\mathrm{i}𝒒(𝑹_n+𝝆_m+𝒂(k/K+l))\right].$$
(36)
Having introduced the definition
$$J_{k,m,m^{}}(𝒒)\underset{n,k^{},l}{}\mathrm{exp}\left[\mathrm{i}𝒒(𝑹_n+𝝆_m^{}𝝆_m+𝒂l+𝒂(k^{}k)/K)\right]J(𝑹_n,m^{},m,l,k^{},k),$$
(37)
it is straightforward to verify that:
$$h^{(m)}(𝒒)=\underset{m^{}=0}{\overset{M1}{}}\underset{\kappa =0}{\overset{K1}{}}𝒥_{\kappa ,m,m^{}}(𝒒)s^{(m^{})}(𝒒+\kappa 𝒂^{}),$$
(38)
where
$$𝒥_{\kappa ,m,m^{}}(𝒒)\frac{1}{K}\underset{k}{}J_{k,m,m^{}}(𝒒)\mathrm{exp}\left[2\pi \mathrm{i}\kappa k/K\right].$$
(39)
Equation 38 enables us to immediately write down the self-consistent RPA equations:
$$\chi ^{(m)}(𝒒,\omega )=\left[\underset{m^{}=0}{\overset{M1}{}}\underset{\kappa =0}{\overset{K1}{}}\chi _0^{(m^{})}(𝒒+\kappa 𝒂^{},\omega )𝒥_{\kappa ,m,m^{}}(𝒒)+1\right]^1\chi _0^{(m)}(𝒒,\omega ).$$
(40)
Here $`\chi _0^{(m)}(𝒒,\omega )`$ and $`\chi ^{(m)}(𝒒,\omega )`$ are the bare (non-interacting) and RPA-corrected wave vector dependent dynamic susceptibilities of spin chains in the $`m`$-th sublattice, respectively. For any wave vector $`𝒒`$ one obtains $`KM`$ equations: susceptibilities at $`K`$ wave vectors become coupled to each other and there are $`M`$ independent chain-sublattices.
For a Heisenberg AF chain, at energy transfers much smaller than the in-chain exchange constant, $`\chi _0^{(m)}(𝒒)`$ is typically very small except when momentum transfer along the chain axis is close to $`q_0`$, defined by $`q_0a=\pi `$. To a good approximation we can thus replace Eq. 40 with:
$`\chi ^{(m)}(𝒒𝒂\pi ,\omega )`$ $``$ $`\left[{\displaystyle \underset{m^{}=1}{\overset{M}{}}}\chi _0^{(m^{})}(𝒒,\omega )J_{m,m^{}}(𝒒)+1\right]^1\chi _0^{(m)}(𝒒,\omega ),`$ (41)
$`\chi ^{(m)}(𝒒𝒂\pi ,\omega )`$ $``$ $`0,`$ (42)
where
$$J_{m,m^{}}(𝒒)\frac{1}{K}\underset{n}{}\underset{l}{}\underset{k=1}{\overset{K}{}}\underset{k^{}=1}{\overset{K}{}}J(𝑹_n,m,m^{},l,k,k^{})\mathrm{exp}\left[\mathrm{i}𝒒(𝑹_n+𝒂l+𝒂(k^{}k)/K+(𝝆_m^{}𝝆_m))\right]$$
(43)
is simply the Fourier transform of exchange interactions between sublattices $`m`$ and $`m^{}`$. Note that there are only $`M`$ equations in the system 41.
## Appendix II: Equivalent Bravais lattice
In this section we shall derive a useful relation between magnetic dynamic structure factors of two spin lattices described by the same Hamiltonian, but featuring different 3D arrangements of magnetic sites. Lattice I is assumed to be a simple Bravais lattice of magnetic sites and consists of $`N`$ unit cells of volume $`v`$. The origin of the $`n`$-th unit cell $`𝑹_n^{(I)}`$ coincides with the position of the $`n`$-th spin: $`𝒓_n^{(I)}𝑹_n^{(I)}`$. Unlike lattice I, lattice II is non-Bravais, having $`M`$ spins per unit cell, and a unit cell volume $`V`$=$`Mv`$. There are $`K=N/M`$ unit cells with origins at $`𝑹_k^{(II)}`$, respectively. Lattice II is obtained from lattice I by shifting each spin $`n`$ by $`𝝆_n`$. In lattice II the spins are thus positioned at $`𝒓_n^{(II)}𝒓_n^{(I)}+𝝆_n`$. The relation between the two lattices allows us to write $`𝒓_n^{(II)}=𝑹_k^{(II)}+𝒓_m^{(II)}=𝑹_k^{(II)}+𝒓_m^{(I)}+𝝆_m`$ ( $`n=kM+m`$, $`0m<M`$). We assume the spin systems to be equivalent in the sense that the spin Hamiltonian of system II written in terms of site-spin operators $`\widehat{𝒔}_n`$ is identical to that of system I. The notation introduced above is independent of that used in the rest of the paper or in Appendix I, and is illustrated, for the particular case of PbNi<sub>2</sub>V<sub>2</sub>O<sub>8</sub>, in Fig. 15.
The Fourier transform of the total spin operator for lattice I is defined as:
$$s^{(I)}(𝒒)=\underset{n=0}{\overset{N1}{}}s_n\mathrm{exp}[\mathrm{i}𝒒𝑹_n^{(I)}]=\underset{k=0}{\overset{K1}{}}\underset{m=0}{\overset{M1}{}}s_{kM+m}\mathrm{exp}[\mathrm{i}𝒒(𝑹_{kM}^{(I)}+𝒓_m^{(I)})].$$
(44)
Again, for simplicity, we have omitted the spin projection indexes. The reverse relation is given by:
$$s_{kM+m}=\frac{v}{(2\pi )^3}𝑑𝒒s^{(I)}(𝒒)\mathrm{exp}[\mathrm{i}𝒒(𝑹_{kM}^{(I)}+𝒓_m^{(I)})],$$
(45)
where the integral is taken over the Brillouin zone of the Bravais lattice ( a reciprocal-space volume $`K`$ times as large as the Brillouin zone of the lattice II). Similarly, for lattice II, by definition:
$$s^{(II)}(𝒒)=\underset{k=0}{\overset{K1}{}}\underset{m=0}{\overset{M1}{}}s_{kM+m}\mathrm{exp}[\mathrm{i}𝒒(𝑹_k^{(II)}+𝒓_m^{(II)})].$$
(46)
The equivalence of the two systems allows us to directly combine the last two equations:
$$s^{(II)}(𝒒)=\frac{v}{(2\pi )^3}\underset{k=0}{\overset{K1}{}}\underset{m=0}{\overset{M1}{}}\mathrm{exp}[\mathrm{i}𝒒(𝑹_k^{(II)}+𝒓_m^{(II)})]𝑑𝒑\mathrm{exp}[\mathrm{i}𝒑(𝑹_{kM}^{(I)}+𝒓_m^{(I)})]s^{(I)}(𝒑).$$
(47)
The expression can be simplified by noting that
$$\underset{k=0}{\overset{K1}{}}\mathrm{exp}[\mathrm{i}(𝒒𝒑)𝑹_k^{(II)}]=\frac{(2\pi )^3}{V}\underset{𝑻}{}\delta (𝒑𝒒𝑻).$$
(48)
Here the inner sum is taken over reciprocal-lattice points $`𝑻_k`$, $`k=0\mathrm{}K`$, for lattice II. Plugging this result into the expression into Eq. 47 gives:
$`s^{(II)}(𝒒)`$ $`=`$ $`{\displaystyle \frac{v}{V}}{\displaystyle 𝑑𝒑s^{(I)}(𝒑)\underset{𝑻}{}\delta (𝒑𝒒𝑻)\underset{m=1}{\overset{M1}{}}\mathrm{exp}[\mathrm{i}(𝒒𝒓_m^{(II)}𝒑𝒓_m^{(I)})]}`$ (49)
$`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{m=0}{\overset{M1}{}}}\mathrm{exp}(\mathrm{i}𝒒𝝆_m){\displaystyle \underset{\mu =1}{\overset{M1}{}}}s^{(I)}(𝒒+𝑻_\mu )\mathrm{exp}(\mathrm{i}𝑻_\mu 𝒓_m^{(I)})`$ (50)
In this the last expression the first sum is taken over all reciprocal-lattice vectors $`𝑻_\mu `$ of lattice II within the first Brillouin zone of lattice I.
Equation 50 for the Fourier transform of spin operators can now be directly plugged into the definition of dynamic structure factors for the two lattices:
$`S^{(I)}(𝒒,\omega ){\displaystyle \frac{1}{2\pi \mathrm{}}}{\displaystyle }dt\mathrm{exp}(\mathrm{i}\omega t)s^{(I)}(𝒒,0)s^{(I)}(𝒒,t))`$ (51)
$`S^{(II)}(𝒒,\omega ){\displaystyle \frac{1}{2\pi \mathrm{}}}{\displaystyle }dt\mathrm{exp}(\mathrm{i}\omega t)s^{(II)}(𝒒,0)s^{(II)}(𝒒,t)).`$ (52)
A straightforward, though somewhat tedious calculation gives the result that we were after:
$$S^{(II)}(𝒒,\omega )=\underset{\mu =0}{\overset{M1}{}}\left|\underset{m=0}{\overset{M1}{}}\mathrm{exp}(\mathrm{i}𝒒𝝆_m)\mathrm{exp}(\mathrm{i}𝑻_\mu 𝒓_m^{(I)})\right|^2S^{(I)}(𝒒+𝑻_\mu ,\omega )$$
(53) |
warning/0001/astro-ph0001410.html | ar5iv | text | # NEW EXTREME SYNCHROTRON BL LAC OBJECTS
## 1. Introduction
BL Lac objects are usually divided in two main classes, on the basis of their overall Spectral Energy Distribution (SED): LBL or HBL (low or high energy peaked BL Lacs), according as the peak of the synchrotron emission (in a $`\nu F_\nu `$ representation) is in the IR–optical or EUV–soft-X band, respectiveley. In the X-ray band this usually translates in a spectral index which is steep ($`\alpha _x>1`$) for HBLs (corresponding to the tail of the synchrotron emission) and flat ($`\alpha _x<1`$) for LBLs (corresponding to the upcoming of the Inverse Compton emission). In 1997, the BeppoSAX observations of Mkn 501 (Pian et al. 1997) and 1ES 2344+514 (Giommi et al. 1997) revealed that, at least in a flaring state, the peak of the synchrotron emission can actually reach very high energies, around 100 keV, with a consequently flat synchrotron X–ray spectral index. In order to find and study other sources with such “extreme” properties, we have selected several candidates from the Einstein Slew Survey and the Rosat All Sky Survey Bright Sources Catalogue (RASSBSC).
The selection criteria were based on properties suggesting a high $`\nu _{peak}`$:
a) very high $`F_x/F_{radio}`$ ratio ($`>3\times 10^{10}`$ erg cm<sup>-2</sup> s<sup>-1</sup> / Jy, at \[0.1–2.4\] keV and 5 GHz respectively); b) flat X-ray spectrum (when available), connecting smoothly with the flux at lower frequencies; c) appropriate values of $`\alpha _{ro}`$, $`\alpha _{ox}`$ and $`\alpha _{rx}`$ (Padovani & Giommi 1995). A high X–ray flux ($`>10^{11}`$ erg cm<sup>-2</sup> s<sup>-1</sup>) in the 2–10 keV band was also requested, to achieve a good detection in the PDS instrument.
We used the BeppoSAX satellite, whose wide X–ray energy range (0.1–200 keV) is ideal to constrain the synchrotron peak. Four objects have been observed, between June 1998 and April 1999: 1ES 0120+340, PKS 0548–322, 1ES 1426+428 and H 2356–309. In Fig. 1 and Table 1 only the best fit results are reported (for a complete discussion, see Costamante et al., in preparation). LECS, MECS and PDS data have been reduced and analysed according to the SDC Cookbook instructions, using the latest calibration matrices available. Standard extraction radii of $`4^{}`$ and $`8^{}`$ for MECS and LECS were used, except for the PKS 0548–322 observation of 20/2/99: in this case a $`6^{}`$ radius for the LECS has been used, due to the presence of a contaminating source in the field of view (identified as the star `GSC_07061_01558` in the Guide Star Catalog, probably flaring).
The PDS instrument doesn’t have imaging capabilities, and its f.o.v. (radius $`45^{}`$) is larger than LECS and MECS ($`28^{}`$ for the MECS). Therefore there is the possibility for PDS spectra to be contaminated by hard serendipitous sources in the f.o.v, not visible in the MECS images. Analyzing PDS data, we have taken this into account, also checking in the NED database for potentially contaminating sources.
## 2. 1ES 1426+428
At $`41^{}`$ from this source, thus in the PDS f.o.v., there is the quasar GB 1428+422 (Fabian et al 1998). To account for its contribution, we have added a component to the PDS model, based on the GB 1428+422 data from the BeppoSAX observation of 4/2/99 ($`\alpha =0.42`$, F$`{}_{1keV}{}^{}=0.30\mu `$Jy; Celotti & Iwasawa, priv. comm.). We have also checked in the NED and WGACAT databases for other potentially contaminating objects: we added the contributions of the two most important objects (WGA J1426.1+4247 and CRSS 1429.7+4240, Fig. 2 right panel), according to the fluxes and spectral indices extrapolated from the ROSAT band (when data were not available in literature, we used galactic N<sub>H</sub> and a HR–$`\alpha `$ conversion by Giommi, priv. comm.). Summing all components, the different off–axis response of the instrument has been taken into account. The PDS/MECS normalization has been fixed at 0.9.
With this model (Fig. 2 right panel), adding the PDS data to the LECS+MECS fit yields a $`\chi _r^2=1.06`$, with the PDS points still slightly above the model (Fig.2 left upper panel). A better $`\chi _r^2`$ (0.95) is obtained with GB1428+422 flux as a free parameter: in this case the resulting flux is F$`{}_{1keV}{}^{}=1.44\pm 0.54`$, a factor more than 4 higher during this observation than the week before. Anyway, in both cases, the spectrum of 1ES 1426+428 remains flatter than unity up to 100 keV.
## 3. Results
The main results for all sources are presented in Table 1. All have been detected in the PDS band. For three of them the spectrum is best fitted with a convex broken power-law: this locates the peak of the synchrotron emission in the X–ray band, around 1-4 keV, thus confirming the “extreme” nature of these sources. The spectrum of 1ES 1426+428 is instead well fitted by a single powerlaw, with a flat spectral index ($`\alpha =0.92`$) up to 100 keV. This constrains the synchrotron peak to lie near or above 100 keV. Such high values of the synchrotron peak frequencies, flagging the presence of high relativistic electrons, make these sources good candidates for TeV emission through the Inverse Compton mechanism.
## ACKNOWLEDGEMENTS
This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with the National Aeronautics and Space Administration. We thank the BeppoSAX Science Data Center for their support in the data analysis. This research is financially supported by the Italian Space Agency.
L.C. thanks the Cariplo Foundation and the Italian MURST for support.
## REFERENCES
Fabian A.C., Iwasawa K. et al., 1998, MNRAS 295L. 25F
Giommi P., Padovani P., Perlman E., Nucl. Physics B (Proc. Suppl.), vol. 69, p.407
Padovani P. & Giommi P., 1995, ApJ 444, 567
Pian E. et al. 1998, ApJ, 492, L17 |
warning/0001/physics0001044.html | ar5iv | text | # Finesse and mirror speed measurement for a suspended Fabry-Perot cavity using the ringing effect
## 1 Introduction
Interferometric gravitational wave detectors, like VIRGO , LIGO , GEO and TAMA , make use of suspended Fabry-Perot cavities for their properties: spatial filtering of the laser beam, optical path amplification and power recycling to name a few. Due to the suspension system, the motion of each mirror is dominated by oscillations at the pendulum’s fundamental resonance, typically below $`1Hz`$ and with amplitudes as large as tens of a laser wavelength.
A VIRGO mode-cleaner prototype MC30 , which operated in Orsay for several years, consists of a $`30m`$ triangular Fabry-Perot cavity with two possible finesse values, 100 or 1600, depending on the incoming laser beam polarization state. The two-stage suspension system, together with a local damping control system, results in a residual RMS displacement value, for each mirror, of $`0.8\mu m`$.
Prior to the lock acquisition, the cavity length sweeps the optical resonance at different rates of expansion. If the relative velocity between the mirrors is constant, the DC transmitted power delineates the Airy peak as a function of time, easily observed for the optical system with $`=100`$. However, in the case of $`=1600`$, a deformation of the Airy peak, similar to a ringing, was observed (see fig.(1)).
Both , and references therein, discuss this phenomenon. Briefly, this effect arises once the cavity sweeps the optical resonance in a time $`\tau _{sw}`$ of the order of or less than the cavity storage time $`\tau _{st}=2L_0/c\pi `$, where $`L_0=30m`$ is the cavity length and $`c`$ the speed of light.
This effect is observed when the rate of expansion is so high that, as resonance is approached, the cavity doesn’t have enough time to completely fill itself. It is the beating between the incoming laser field and the evolving stored field that gives rise to this oscillatory behavior.
The goal of this letter is to present the informations which can be extracted from such deformation, in particular, the finesse of the cavity and the relative mirror speed.
## 2 The ringing effect in Fabry-Perot cavities
The model used for this study is shown in fig.(2). Assuming a negligible mirror displacement for times of the order of the round trip time of light $`\tau =2L_0/c=0.2\mu s`$, the stored field $`\mathrm{\Psi }_1(t)`$ at time $`t`$ can be written as
$$\mathrm{\Psi }_1(t)=t_1\mathrm{\Psi }_{in}+r_1^2\mathrm{exp}(\mathrm{\hspace{0.17em}2}ikL)\mathrm{\Psi }_1(t\tau )$$
(1)
where $`r_1`$ and $`t_1`$ denote the amplitude reflectivity and transmittivity of each mirror, $`\mathrm{\Psi }_{in}`$ is the incoming laser field, and $`L`$ is the cavity length. Assuming that the cavity expands at a constant rate $`v`$, we can write $`L=L_0+vt`$ and solve eq.(1) iteratively, for different velocities $`v`$ and finesse $``$.
Fig.(3) shows the Airy peak for three velocities: $`v=0`$ (static approximation), $`v=1\lambda /s`$, and $`v=2.6\lambda /s`$, with $`=4000`$. The curve labeled static, corresponding to $`v=0`$, was generated by neglecting the travel time of light, i.e the cavity has reached its equilibrium point at each step. The two other curves, on the other hand, were simulated according to the dynamical model here presented. Notice how the main peak height decreases, its width increases and its position shifts ahead of the resonance. These changes are greater for larger velocities.
## 3 The speed measurement
We would now like to discuss a property of the ringing effect observed from the simulation runs. Fig.(3) graphs the stored power as a function of cavity length for a given finesse and for different values of velocity. We can now plot the stored power as a function of time, setting the velocity to a fixed value, but varying the finesse. One example is given in fig.(4). The top graph of this figure shows the stored power as a function of time, for an expansion rate set to $`v=10\lambda /s`$, for three different finesse values: $`=1000`$, $`2000`$, and $`3000`$. The bottom graph is the curves’ time derivative. From these plots, we remark a particular characteristic of the phenomenon: the position of the minima and maxima, with the exception of the main peak, are almost independent from the finesse value.
Furthermore, going back to fig.(3), we can now note that the derivative zeros depend only on the relative mirror velocity. The simulation output shown in fig.(4) not only shows how the derivative zeros are independent, at least to first approximation, from the finesse, but it also shows a particular regularity in the spacing between the minima and maxima.
The upper graph of fig.(5) shows the simulated stored power of a cavity with $`=3500`$, expanding at a rate $`10\lambda /s`$. Let the position of the curve’s derivative zeros, $`t_{zero}`$, be labeled by the index $`n`$, so that, for the first zero, positioned at $`t_{zero}5.07ms`$, $`n=0`$, for the second zero, located at $`t_{zero}5.103ms`$, $`n=1`$ and so on. Then, the bottom graph of fig.(5) shows the plot of index $`n`$ as a function of time. We remark that the $`n`$-th zero of the derivative is a quadratic function of the zero crossing time $`t_{zero}`$: $`nt_{zero}^2`$. By fitting the simulation outputs to the expression $`n_{zero}=p_1+p_2t_{zero}+p_3t_{zero}^2`$ where $`p_{1,2,3}`$ are fitting parameters, we empirically found that the coefficient $`p_3`$ can be written as $`p_3=cv/\lambda L`$ where $`L`$ is the cavity length and $`v`$ is the cavity expansion rate (an example is shown in the bottom graph of fig.(5)). Therefore, an estimate of coefficient $`p_3`$ would also give us an estimate of the relative velocity $`v`$.
Fig.(6) shows the results of a fit on a measured event. Notice how the parabolic behavior is in agreement with the experimental points, resulting in a measured speed of $`12.8\pm 10^2\mu m/s`$.
## 4 The Finesse measurement
Once the speed is extracted, it is possible to fit the measurements with the simulation’s output to find the remaining parameter: the finesse. The fit results of a set of measurements, six of which are shown in fig.(7), led to a mean finesse of $`\overline{}=1554\pm 160`$, a value later confirmed by a measurement of the cavity pole. The ten percent precision on the finesse measurement is most probably due to the alignment state of the cavity, as suggested by simulation studies.
## 5 Conclusion
An analysis of the optical ringing effect, observed on the VIRGO mode-cleaner prototype in Orsay, was here presented. We investigated a method to extract, from the oscillatory behavior, both the relative mirror speed and the finesse of the system.
The numerical results showed how the position of the oscillations’ minima and maxima, when plotted as a function of time, weakly depend on the finesse value and are completely determined by the cavity expansion rate as the resonance is being crossed. In particular, we showed how a simple empirical formula can determine the cavity expansion rate by observing these minima and maxima.
Once the speed was reconstructed, it was possible to fit the measurements with the simulation’s output and estimate the cavity’s finesse to $`\overline{}=1554\pm 160`$.
The present letter gives an alternative method to the finesse measurement for a suspended cavity. The simplicity in the velocity reconstruction algorithm may be useful for a future control system capable of guiding the cavity into lock.
Figure Captions
Fig.1: The observed ringing effect on the transmitted DC power of the MC30 prototype. The transmitted power is shown as a function of time as the cavity length sweeps the optical resonance at an unknown rate.
Fig.2: The model used for the study of the MC30 ringing effect.
Fig.3: The calculated Fabry-Perot transmitted power, with $`=4000`$, as a function of cavity length $`\mathrm{\Delta }L`$ as the resonance is swept at $`v=0`$ (static approximation), $`v=1\lambda /s`$, and $`v=2.6\lambda /s`$. In the figure, $`\mathrm{\Delta }L=0`$ corresponds to resonance.
Fig.4: The calculated stored power as a function of time, with a fixed expansion rate set to $`v=10\lambda /s`$, for different finesse values: $`=1000`$, $`2000`$, and $`3000`$. Top graph: the stored power. Bottom graph: the stored power time derivative.
Fig.5: The simulated stored power of a Fabry-Perot, expanding at a constant rate $`v=10\lambda /s`$, with $`=3500`$. Top graph: the stored power as a function of time. Bottom graph: the index $`n`$, corresponding to the $`n`$-th derivative zero, as a function of time. The curve is fit to the expression $`n=p_1+p_2t+p_3t^2`$. Notice that $`p_3=cv/\lambda L=100[1/ms^2]`$.
Fig.6: Fit results for the mirror relative velocity reconstruction.On the left: The measured DC transmitted power. On the right: the plot of $`t_{zero}`$ as a function of index $`n`$. The error bars correspond to half of the oscilloscope’s sampling time. The reconstructed speed is $`12.8\pm 10^2\mu m/s`$.
Fig.7: The observed ringing effect for the MC30 prototype: measurements and fits. The finesse and velocity values are shown for each graph. $`T_{DC}`$ in arbitrary units. |
warning/0001/hep-th0001172.html | ar5iv | text | # Contents
## Chapter 1 SOME GENERAL FACTS
After a short introduction, the most important known facts about sine-Gordon and massive Thirring models are exposed. It is also explained their connection with the $`c=1`$ free boson.
### 1.1 Introduction
As a very large number of papers in literature, this thesis principally deals with the sine-Gordon model, which has been well known at the classical level for the late fifty years and plays also an important role in quantum theory, thanks to its particular properties of non-linearity and integrability. It has been successfully applied in very different sectors of Mathematics and Physics, from partial differential equation theory to particle physics or solid state physics. Recent applications of the classical model are related to nonlinear optics (resonant dielectric media) and optical fibers, magnetic properties of polymers, propagation of waves in crystals and so on. Interesting applications of the quantum model are related to Kondo effect and to the thermodynamics of some chemical compound, as can be found in (see also section 4.10). At the same time, the quantum theory shows a phenomenology that is similar to the Skyrme model used before QCD era to describe barions and strong interactions.
The most relevant properties of the model are
* at a classical level, all the solutions of the equations of motion are known (exact integrability via inverse scattering method)
* the classical solutions describe solitons, antisolitons and bound states (breathers)<sup>1</sup><sup>1</sup>1 The so called mesonic solutions are excluded in this analysis. ; in a scattering process this solutions are transparent (it is the mathematical meaning of “soliton”)
* it admits, both at a classical and at the quantum level, a countable infinite set of conserved charges
* the quantization of the theory describes an interacting particle with its antiparticle and, in a certain (attractive) regime, bound states
* the S matrix has been exactly determined; only elastic scattering processes can take place (i.e. no particle production), that is the quantum analog of the classical transparency of solitons
#### Finite size effects
Finite size effects are widely recognized to play an important role in modern statistical mechanics and quantum field theory. From a statistical point of view, it is known that no phase transitions take place in a finite volume system. For example, specific heat $`c(T)`$, that is divergent at the critical point, if the system has finite size looses it divergence; one observes only a rounded peak, in the plot $`c(T)`$ versus $`T`$. Moreover, there is only an interval around critical temperature $`T_c`$ where the finite size effects are relevant. Out of this interval, they are negligible (because only near $`T_c`$ the correlation length can be comparable with the size of the system). The interesting fact is that specific heat (and other critical quantities) have a scaling behaviour (i.e. varying the size $`L`$) that is fixed by the (infinite size) critical exponents (see ). This is a general fact: as argued in , the UV behaviour of the scaling functions (see later) is fixed by the conformal dimensions of the operators that belong to the universality class of the critical point (i.e. the CFT describing the critical point of the statistical system).
Also in quantum field theory interesting phenomena appear. If the space-time geometry is a cylinder of circumference $`L`$, Casimir effects change the energy of a two body interaction, because particles interact in the two possible directions, as shown in figure 1.1 (looking forward one can see his own back). Also new radiative corrections to a propagating particle may appear because of the closed geometry.
A system on a finite volume has discrete energy and momentum spectra. Lüscher has show that the corrections to the free system eigenvalues of $`E,P`$ depends on the scattering amplitudes (i.e. on the $`S`$ matrix). This general fact can be used to extract informations on an infinite space QFT from a numerical simulation, that obviously is affected by finite size effects.
Notice that a cylinder geometry can also be interpreted as a finite temperature field theory (temperature is $`L^1`$). Such a theory is related to the phase transition between confined and deconfined phase in QCD.
In general, a physical quantity (for example the free energy) that depends on the renormalization scale $`L`$ (it parameterizes the renormalization flow) is a *scaling function*. The energy and momentum computed in (3.58, 3.59) are scaling functions, because in the NLIE the scale $`L`$ appears.
There are many methods of investigation of the finite size effects. One of them, that has been very fruitful even for theories which are not integrable, is the *truncated conformal space approach* (see and section 1.4), which is an intrinsically non-perturbative approximation method. It has problems of principal nature, coming from the fact that one does not have an analytic control of the spectrum, and of practical nature, because to reach a certain precision in the resulting energy levels one has sometimes to resort to very high truncation levels and introduce enormous matrices to diagonalize.
For integrable QFT, there also exist *exact analytic* methods to compute the finite size effects, like e.g. the *Thermodynamic Bethe Ansatz* (TBA), which was used to calculate the vacuum (Casimir) energy . The method was later extended to include ground states of charged sectors . More recently, using analytic properties of the TBA equations extended for complex values of the volume parameter, an approach to get excited states was proposed in . Their method to get excited states sheds light on the analytic structure of the dependence of scaling functions on the spatial volume and up to now was the only method developed to deal with excited states in perturbations of minimal models. Its main drawback is that to obtain the equation for a given excited state one has to do analytic continuation for each case separately, and a major part of this continuation can only be carried out numerically. Because of the complications of the analytic continuation, this method is limited at present to simple cases of integrable perturbations of Virasoro minimal models and some other perturbed conformal field theories. Similar results were obtained in .
This thesis reports on a novel approach to the excited states of IQFTs in finite volume, based on the *nonlinear integral equation* (NLIE) method, which has its origin in the so-called *light-cone lattice Bethe Ansatz* approach to regularize integrable QFTs. It was argued in that sine-Gordon theory can be regularized using an inhomogeneous $`6`$-vertex model (or equivalently, an inhomogeneous $`XXZ`$ chain). The NLIE was originally developed in this framework to describe the ground state scaling function (Casimir energy) in sine-Gordon theory in and it was shown that in the ultraviolet limit it reproduces the correct value of the central charge $`c=1`$. Similar methods were independently introduced in Condensed Matter Physics by other authors .
The NLIE was first extended to excited states in where the spectrum of states containing only solitons (and no antisolitons/breathers) has been described. Using an idea by Zamolodchikov they also showed that a twisted version of the equation was able to describe ground states of unitary Virasoro minimal models perturbed by the operator $`\mathrm{\Phi }_{(1,3)}`$. A framework for generic excited states of even topological charge in sine-Gordon theory was outlined by Destri and de Vega in . However, there has been a contradiction between the results of the two papers, which was resolved in where it was shown that it was related to the locality and the operator content of limiting ultraviolet conformal field theory (CFT). Besides that, strong evidence for the correctness of the predicted spectrum was given by comparing it to predictions coming from the truncated conformal space (TCS) method, pioneered by Yurov and Zamolodchikov in and extended to $`c=1`$ theories in . Later a modification of the NLIE to describe the states of sine-Gordon/massive Thirring theory with odd topological charge was conjectured .
The NLIE for sine-Gordon theory was generalized to models built on general simply-laced algebras of $`ADE`$ type in for the case of the vacuum. More recently, in P. Zinn-Justin extended the method to the spectrum of excited states for these models and he also made a first attempt to describe perturbations of minimal models of CFT. A general framework for describing general excited states of minimal models perturbed by $`\mathrm{\Phi }_{(1,3)}`$ (only massive case) can be found in , where is stated the correct form of NLIE equation to deal with this case, and also the simplest examples of the resulting excited states are checked (see chapter 4).
In the following chapters, the general setup of NLIE will be presented, with the most relevant examples.
Chapter 1 is devoted to summarize some well known facts that will be used in the following.
In Chapter 2 the light-cone lattice is introduced and the Bethe equations for the 6 vertex model are written.
Chapter 3 is devoted to obtain an integral equation equivalent to Bethe Ansatz, but that allows a continuum limit procedure. The analysis of the so obtained continuum theory is in Chapter 4.
### 1.2 c=1 CFT: free boson
To fix some conventions and to define certain objects which are used later, a brief summary is given of the $`c=1`$ free boson with a target space of a circle of radius $`R`$. The Lagrangian of this CFT is taken to be
$$=\frac{1}{8\pi }_0^L_\mu \phi (x,t)^\mu \phi (x,t)dx,x[0,L],$$
(1.1)
where $`L`$ is the spatial volume (i.e. the theory is defined on a cylindrical spacetime with circumference $`L`$). In the sequel often the complex Euclidean coordinates will be used $`z=e^{2\pi (tix)/L},\overline{z}=e^{2\pi (t+ix)/L}`$. The superselection sectors are classified by the $`\widehat{U(1)}_L\times \widehat{U(1)}_R`$ Kac-Moody symmetry algebra, generated by the currents
$$J(z)=i_z\phi ,\overline{J}(\overline{z})=i_{\overline{z}}\phi .$$
The left/right moving energy-momentum tensor is given by
$$T(z)=\frac{1}{8\pi }_z\phi _z\phi =\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}L_kz^{k2},\overline{T}(\overline{z})=\frac{1}{8\pi }_{\overline{z}}\phi _{\overline{z}}\phi =\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\overline{L}_k\overline{z}^{k2}$$
The coefficients $`L_n`$ and $`\overline{L}_n`$ of the Laurent expansion of these fields generate two mutually commuting Virasoro algebras. If the (quasi)periodic boundary conditions are required
$$\phi (x+L,t)=\phi (x,t)+2\pi mR,m,$$
then the sectors are labelled by a pair of numbers $`(n,m)`$, where $`\frac{n}{R}`$ ($`n`$ is half integer because of the locality, see later) is the eigenvalue of the total field momentum $`\pi _0`$
$$\pi _0=_0^L\pi (x,t)𝑑x,\pi (x,t)=\frac{1}{4\pi }_t\phi (x,t),$$
and $`m`$ is the winding number, i.e. the eigenvalue of the topological charge $`Q`$ defined by
$$Q=\frac{1}{2\pi R}_0^L_x\phi (x,t)dx.$$
In the sector with quantum numbers $`(n,m)`$, the scalar field is expanded in modes as follows:
$$\begin{array}{cc}\hfill \phi (x,t)=& \varphi (z)+\overline{\varphi }(\overline{z}),\hfill \\ \hfill \varphi (z)=& \frac{1}{2}\phi _0ip_+\mathrm{log}z+i\underset{k0}{}\frac{1}{k}a_kz^k,\hfill \\ \hfill \overline{\varphi }(\overline{z})=& \frac{1}{2}\phi _0ip_{}\mathrm{log}\overline{z}+i\underset{k0}{}\frac{1}{k}\overline{a}_k\overline{z}^k,\hfill \end{array}$$
where the left and right moving field momenta $`p_\pm `$ (which are in fact the two $`U(1)`$ Kac-Moody charges) are given by
$$p_\pm =\frac{n}{R}\pm \frac{1}{2}mR.$$
(1.2)
The Virasoro generators take the form
$$L_n=\frac{1}{2}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}:a_{nk}a_k:,\overline{L}_n=\frac{1}{2}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}:\overline{a}_{nk}\overline{a}_k:,$$
where the colons denote the usual normal ordering, according to which the oscillator with the larger index is put to the right.
The ground states of the different sectors $`(n,m)`$ are created from the vacuum by the (Kac-Moody) primary fields, which are vertex operators of the form
$$V_{(n,m)}(z,\overline{z})=:\mathrm{exp}i(p_+\varphi (z)+p_{}\overline{\varphi }(\overline{z})):.$$
(1.3)
The left and right conformal weights of the field $`V_{(n,m)}`$ (i.e. the eigenvalues of $`L_0`$ and $`\overline{L}_0`$) are given by the formulae
$$\mathrm{\Delta }^\pm =\frac{p_\pm ^2}{2}.$$
(1.4)
The Hilbert space of the theory is given by the direct sum of the Fock modules built over the states
$$|n,m=V_{(n,m)}(0,0)|vac,$$
(1.5)
with the help of the creation operators $`a_k,\overline{a}_kk>0`$:
$$=\underset{(n,m)}{}\{a_{k_1}\mathrm{}a_{k_p}\overline{a}_{l_1}\mathrm{}\overline{a}_{l_q}|n,m,k_1,\mathrm{}k_p,l_1,\mathrm{}l_q_+\}$$
The boson Hamiltonian on the cylinder is expressed in terms of the Virasoro operators as
$$H_{CFT}=\frac{2\pi }{L}\left(L_0+\overline{L}_0\frac{c}{12}\right),$$
(1.6)
where the central charge is $`c=1`$. The generator of spatial translations is given by
$$P=\frac{2\pi }{L}\left(L_0\overline{L}_0\right).$$
(1.7)
The operator $`L_0\overline{L}_0`$ is the conformal spin which has eigenvalue $`nm`$ on the primary field $`V_{(n,m)}`$.
One can also introduce twisted sectors using the operator $`𝒯`$ that performs spatial translations by $`L`$: $`xx+L`$. The primary fields $`V_{(n,m)}`$ as defined above satisfy the periodicity condition $`𝒯V_{(n,m)}=V_{(n,m)}.`$ If the more general twisted boundary condition labelled by a real parameter $`\nu `$ is required
$$𝒯V_{(n,m)}=\mathrm{exp}\left(i\nu Q\right)V_{(n,m)},$$
then it is possible to generate superselection sectors for which $`n+\frac{\nu }{2\pi }`$.
It is important to stress that a particular $`c=1`$ CFT is specified by giving the spectrum of the quantum numbers $`(n,m)`$ (and the compactification radius $`R`$) such that the corresponding set of vertex operators (and their descendants) forms a *closed and local* operator algebra. The locality requirement is equivalent to the fact that the operator product expansions of any two such local operators is single valued in the complex plane of $`z`$. This condition, which is weaker than the modular invariance of the CFT, is the adequate one since the theory is considered on a space-time cylinder and do not wish to define it on higher genus surfaces.
By this requirement of locality, it was proved in that there are only two maximal local subalgebras of vertex operators: $`𝒜_b`$ generated by the vertex operators
$$\{V_{(n,m)}:n,m\},$$
and $`𝒜_f`$ generated by
$$\{V_{(n,m)}:n,m2orn+\frac{1}{2},m2+1\}.$$
Other sets of vertex operators can be built, but the product of two of them gives a nonlocal expression.
### 1.3 Sine-Gordon/massive Thirring field theory
The minkowskian lagrangian of sine-Gordon theory is given by<sup>2</sup><sup>2</sup>2 Integration sign without any specification means integration on the whole real axis.
$$_{sG}=(\frac{1}{2}_\nu \mathrm{\Phi }^\nu \mathrm{\Phi }+\frac{\mu ^2}{\beta ^2}:\mathrm{cos}\left(\beta \mathrm{\Phi }\right):)dx,$$
(1.8)
where $`\mathrm{\Phi }`$ denotes a real scalar field, while that of the massive Thirring theory is of the following form:
$$_{mTh}=\left(\overline{\mathrm{\Psi }}(i\gamma _\nu ^\nu +m_0)\mathrm{\Psi }\frac{g}{2}\overline{\mathrm{\Psi }}\gamma ^\nu \mathrm{\Psi }\overline{\mathrm{\Psi }}\gamma _\nu \mathrm{\Psi }\right)𝑑x,$$
(1.9)
describing a current-current selfinteraction of a Dirac fermion $`\mathrm{\Psi }`$. It is known that the two theories are deeply related provided their coupling constants satisfy
$$\frac{\beta ^2}{4\pi }=\frac{1}{1+g/\pi }.$$
For comparison with the Destri-de Vega nonlinear integral equation, it is important to deal with cylindrical sine-Gordon and massive Thirring, i.e. the integrals in (1.8, 1.9) must be taken in the interval $`[0,L]`$.
The $`\mathrm{cos}`$ term in (1.8) can be considered as a perturbation of the $`c=1`$ free boson compactified on a cylinder, as described in section 1.2. Similarly the massless ($`m_0=0`$) Thirring model is a $`c=1`$ conformal field theory and the mass term plays the role of a perturbation. From Coleman’s paper it is known that correlation functions of the perturbing fields $`\overline{\mathrm{\Psi }}\mathrm{\Psi }`$ and $`:\mathrm{cos}\beta \mathrm{\Phi }:`$ are identical, then both models can be considered as the perturbations of a $`c=1`$ compactified boson by a potential $`V`$ :
$$H_{sG/mTh}=H_{CFT}+V,V=\lambda _0^L(V_{(1,0)}(z,\overline{z})+V_{(1,0)}(z,\overline{z}))dx,$$
(1.10)
which is related to the bosonic lagrangian (1.8) by the following redefinitions of the field and the parameters:
$$\phi =\sqrt{4\pi }\mathrm{\Phi },R=\frac{\sqrt{4\pi }}{\beta },\lambda =\frac{\mu ^2}{2\beta ^2}.$$
(1.11)
For later convenience, a new parameter $`p`$ can be defined by
$$p=\frac{\beta ^2}{8\pi \beta ^2}=\frac{1}{2R^21}.$$
(1.12)
The point $`p=1`$ (i.e. $`g=0`$) is the free fermion point, corresponding to a massive Dirac free fermion. The particle spectrum of sG for $`p>1`$ is composed by the soliton ($`s`$) and its antiparticle, the antisoliton ($`\overline{s}`$). It is known as repulsive regime because no bound states can take place. $`p<1`$ is the attractive regime, because $`s`$ and $`\overline{s}`$ can form bound states that are known as breathers. The values $`p={\displaystyle \frac{1}{k}},k=1,2,\mathrm{}`$ are the thresholds where a new bound state appears. The potential term becomes marginal when $`\beta ^2=8\pi `$ which corresponds to $`p=\mathrm{}`$. The perturbation conserves the topological charge $`Q`$, which can be identified with the usual topological charge of the sG theory and with the fermion number of the mTh model.
Mandelstam showed that a fermion operator satisfying the massive Thirring equation of motion can be constructed as a nonlocal functional of a pseudoscalar field (boson) satisfying the sine-Gordon equation. But the fermion and the boson are not relatively local and then do not create the same particle (the two theories are not equivalent).
The difference between them is that they correspond to the perturbation by the same operator of the two *different local* c=1 CFTs $`𝒜_b`$ and $`𝒜_f`$ as in 1.2. The short distance behaviour of the sG theory is described by the local operator algebra $`𝒜_b`$, while the primary fields of the UV limit of mTh theory are $`𝒜_f`$.
Note that the two algebras share a common subspace with even values of the topological charge, generated by $`\{V_{(n,m)}:n,m2\}`$, where the massive theories described by the lagrangians (1.8) and (1.9) are identical. Exactly in this subspace holds the well known proof by Coleman about the equivalence of the two theories . The figure 1.2 shows the four sectors where all the vertex operators live.
### 1.4 Truncated Conformal space at $`c=1`$
In this section is given a brief description of Truncated Conformal Space (TCS) method for $`c=1`$ theories. It will be usefull, in chapter 4 to have a comparison of the data obtained from the non-linear integral equation that will be introduced.
The TCS method was originally created to describe perturbations of Virasoro minimal models in finite spatial volume . Here is presented an extension of the method to study perturbations of a $`c=1`$ compactified boson, more closely the perturbation corresponding to sine-Gordon theory.
The Hilbert space of the $`c=1`$ theory (on a cylinder) can be split into sectors labelled by the values of $`P`$ and $`Q`$, which are quantised by integers. The numerical computations shall be restricted to the $`P=0`$ sector (it is not expected that any relevant new information would come from considering $`P0`$). The TCS method consists of retaining only those states in such a sector for which the eigenvalue of $`H_{CFT}`$ is less than a certain upper value $`E_{cut}`$, so the truncated space is defined as
$$_{TCS}(s,m,E_{cut})=\left\{\right|\mathrm{\Psi }:P|\mathrm{\Psi }=s|\mathrm{\Psi },Q|\mathrm{\Psi }=m|\mathrm{\Psi },H_{CFT}|\mathrm{\Psi }E_{cut}|\mathrm{\Psi }\}.$$
(1.13)
For a given value of $`s`$, $`m`$ and $`E_{cut}`$ this space is always finite dimensional. In this space, the Hamiltonian, represented on the basis of the eigenstates of the unperturbed $`c=1`$ Hamiltonian (i.e. eigenstates of $`\widehat{L_0}+\widehat{\overline{L}_0}`$)<sup>3</sup><sup>3</sup>3 In this section, an hat on a letter means matrix , is a finite size matrix whose eigenvalues can be computed using a numerical diagonalization method. The explicit form of this matrix is the following:
$$\widehat{H}=\frac{2\pi }{L}\left(\widehat{L_0}+\widehat{\overline{L}_0}\frac{c}{12}\widehat{I}+\lambda \frac{L^{2h}}{\left(2\pi \right)^{1h}}\widehat{B}\right),$$
(1.14)
where $`\widehat{L_0}`$ and $`\widehat{\overline{L_0}}`$ are diagonal matrices with their diagonal elements being the left and right conformal weights, $`\widehat{I}`$ is the identity matrix,
$$h=\mathrm{\Delta }_V^++\mathrm{\Delta }_V^{}=\frac{\beta ^2}{4\pi }=\frac{2p}{p+1}$$
(1.15)
is the scaling dimension of the perturbing potential $`V`$ and the matrix elements of $`\widehat{B}`$ are
$$\widehat{B}_{\mathrm{\Phi },\mathrm{\Psi }}=\frac{1}{2}\mathrm{\Phi }\left|V_{(1,0)}(1,1)+V_{(0,1)}(1,1)\right|\mathrm{\Psi }.$$
(1.16)
The units are chosen in terms of the soliton mass $``$ which is related to the coupling constant $`\lambda `$ by the mass gap formula<sup>4</sup><sup>4</sup>4 Notice the analogy with both s-G/mTh that are obtained perturbing the free boson theory by a “mass term”. obtained from TBA in :
$$\lambda =\kappa ^{2h},$$
(1.17)
where
$$\kappa =\frac{2\mathrm{\Gamma }(h/2)}{\pi \mathrm{\Gamma }(1h/2)}\left(\frac{\sqrt{\pi }\mathrm{\Gamma }\left({\displaystyle \frac{1}{2h}}\right)}{2\mathrm{\Gamma }\left({\displaystyle \frac{h}{42h}}\right)}\right)^{2h}.$$
(1.18)
In what follows the energy scale is normalized by taking $`=1`$; the dimensionless volume $`L`$ is denoted by $`l`$. For numerical computations, the dimensionless Hamiltonian
$$\widehat{h}=\frac{\widehat{H}}{}=\frac{2\pi }{l}\left(\widehat{L_0}+\widehat{\overline{L}_0}\frac{c}{12}\widehat{I}+\kappa \frac{l^{2h}}{\left(2\pi \right)^{1h}}\widehat{B}\right)$$
(1.19)
will be used. The usefulness of the TCS method lies in the fact that it provides a nonperturbative method of numerically obtaining the spectrum (the mass gap, the mass ratios and the scattering amplitudes) of the theory. Therefore it can serve as a tool to check the exact results obtained for integrable field theories and get a picture of the physical behaviour even for the nonintegrable case. The systematic error introduced by the truncation procedure is called the *truncation error*; it increases with the volume $`L`$ and can be made smaller by increasing the truncation level (at the price of increasing the size of the matrices, which is bound from above by machine memory and computation time).
Let us make some general remarks on how the TCS method applies to $`c=1`$ theories. First note that the Hilbert space (even after specifying the sector by the eigenvalues of $`P`$ and $`Q`$) consists of infinitely many Verma modules labelled by the quantum number $`n`$. At any finite value of $`E_{cut}`$ only finitely many of such Verma modules contribute, but their number increases with $`E_{cut}`$. As a result one has to deal with many more states than in the case of minimal models. The results of TCS are supposed to approach the exact results in the limit $`E_{cut}\mathrm{}`$. The convergence can be very slow, while the number of states rises faster than exponentially with the truncation level. The perturbing operator has scaling dimension which ranges between $`0`$ and $`2`$, becoming more relevant in the attractive regime, while the number of states corresponding to a given value of $`E_{cut}`$ becomes larger as moving towards $`p=0`$, which affects the convergence just the other way around.
Generally, the energy of any state goes with the volume $`L`$ as
$$\frac{E_\mathrm{\Psi }(L)}{}=\frac{\pi \left(c12\left(\mathrm{\Delta }_\mathrm{\Psi }+\overline{\mathrm{\Delta }}_\mathrm{\Psi }\right)\right)}{6l}+Bl+\underset{k=1}{\overset{\mathrm{}}{}}C_k\left(\mathrm{\Psi }\right)l^{k(2h)},$$
(1.20)
where $`\mathrm{\Delta }_\mathrm{\Psi }`$ ($`\overline{\mathrm{\Delta }}_\mathrm{\Psi }`$) are the left (right) conformal dimensions of the state in the ultraviolet limit, $`B`$ is the universal bulk energy constant (the vacuum energy density) and the infinite sum represents the perturbative contributions from the potential $`V`$.
The bulk energy constant has also been predicted from TBA and reads
$$B=\frac{1}{4}\mathrm{tan}\left(\frac{p\pi }{2}\right)$$
(1.21)
(the same result was obtained from the NLIE approach in ). This is a highly nonanalytic function of $`p`$ and it becomes infinite at the points where $`p`$ is an odd integer. In fact, at these points there is a value of $`k`$ for which $`k(2h)=1`$, and $`C_k\left(\mathrm{\Psi }\right)\mathrm{}`$. The infinite parts of $`B`$ and $`C_k\left(\mathrm{\Psi }\right)`$ exactly cancel, leaving a logarithmic (proportional to $`l\mathrm{log}l`$) and a finite linear contribution to the energy, by a sort of a resonance mechanism. All of these “logarithmic points” are in the repulsive regime. However, due to UV problems in the repulsive regime we are not able to check numerically the logarithmic corrections to the bulk energy.
The origin of UV divergences can be understood from *conformal perturbation theory* (CPT). It is known that when the scaling dimension $`h`$ of the perturbing potential exceeds $`1`$, CPT suffers from ultraviolet divergences which should be removed by some renormalization procedure. The TCS method is something very similar to CPT: it operates in the basis of the UV wave functions as well, but computes the energy levels using the variational approach and therefore could be called “*conformal variation theory*(CVT). As a result, it is expected that there could be UV divergences for the range of couplings where $`h>1`$ which is exactly the repulsive regime $`p>1`$ . The numerical analysis has in fact shown that in the repulsive regime the TCS energy eigenvalues did not converge at all when increasing the truncation level.
Fortunately, there exists a way out: since it is expected to find a sensible quantum field theory when the UV cutoff is removed, it should be the case that the *relative energy levels* $`_\mathrm{\Psi }(L)=E_\mathrm{\Psi }(L)E_{vac}(L)`$ converge to some limit. This is exactly the behaviour that has been observed. Consequently, in the repulsive regime one can only trust the relative scaling functions produced by the TCS method, while in the attractive regime even the absolute energy values can be obtained (including the predicted bulk energy constant (1.21), which is completely analytic for $`p<1`$ and thus logarithmic corrections are absent as well).
Many numerical results show that the smaller the value of $`p`$ is the faster the convergence is (with the understanding that in the repulsive regime by convergence of TCS we mean the convergence of the energies relative to the vacuum). On the other hand, even in the attractive regime the convergence is so slow that to get reliable results (which means errors of order $`10^310^2`$ for the volume $`l`$ ranging from $`0`$ to somewhere between $`5`$ and $`10`$) requires to work with matrix of dimensions around $`4000`$. This means that the TCS for $`c=1`$ theories is far less convergent than the one for minimal models (in the original Lee-Yang example the authors of took a $`17`$ dimensional Hilbert space (!) and arrived to very accurate results).
## Chapter 2 LIGHT-CONE LATTICE QFT
In this chapter, the most important tools to deal with Destri-de Vega equation are introduced, with reference to the original papers . Particular attention will be taken to indicate the path that has been done from the light-cone lattice until the definition of a particular system (the 6 vertex model, alias XXZ chain), whose Hilbert space is quite completely known.
### 2.1 Kinematics on light-cone
It is a usual way to regularize quantum field theories by defining them on a space-like “hamiltonian” lattice (where time is continuous and space discrete) or space and time-like “euclidean” lattice (when both space and time are discrete). In statistical mechanics this is not just a regularization method but can be a right microscopic way to describe physical systems. In two dimensions, the most known approach is to define a rectangular lattice with axis corresponding to space and time directions and associate to each site an interaction depending only on the nearest neighbouring sites. In this case the partition function can be expressed in terms of a transfer matrix.
In what follows, a different approach is adopted: Minkowski and Euclidean space-time can, in fact, be discretized along light-cone directions. Light-cone coordinates are:
$$x_\pm =x\pm t$$
and the choice
$$=\{x_\pm =\frac{a}{\sqrt{2}}n_\pm ,n_\pm \}$$
defines a light-cone lattice of “events” as in figure 2.1 (a) . They are spaced by $`a`$ in the space and time directions and by $`a/\sqrt{2}`$ in light-cone directions. At every event $`P`$ there is associated a double light-cone (in the past and in the future) and only events within this light-cone can be causally connected (see fig. 2.1 (b)). Then, any rational and not greater than $`1`$ speed is permitted for particles, in an infinite lattice. The shortest displacement of the particle (one lattice spacing) is realized at speed $`\pm 1`$ and corresponds, from the statistical point of view, to nearest neighbours interactions. Smaller speeds can be obtained with displacements longer than the fundamental plaquette, and correspond to high order neighbours interactions. In quantum field theory, these are nonlocal interactions.
In the following, only the local case (nearest neighbours) is treated. The nearest neighbours of the event $`P`$ are the four nearest points in the light-cone directions. This implies that particles can have only the speed of light $`\pm 1`$ and are massless. They are called right-movers (R) and left-movers (L).
In this case, it is possible to introduce a useful language for connection with statistical mechanics associating a particle to a link. Consider the two links in the future that come out from a event $`P`$. Particles R and L in $`P`$, by definition, are respectively associated to the right-oriented link and to the left-oriented link. In this way, the state of a link is defined to be the state of the point where it begins (also the opposite choice, of connecting a link with the site where it ends, can be done; it is simply a matter of convention). For example, if in a point $`O`$ there is a particle R, one tells that the “right-oriented” link outcoming from it is occupied by R. This correspondence of points and links is possible because only local interactions are assumed, and it is useful because the counting of states is simpler. But the physically correct interpretation is that particles live on events, not on links. Links are the possible world lines for particles.
In the following, the lattice is assumed of a finite extent $`L=aN`$ in space direction ($`N`$ is the number of sites, counted in the space direction), with periodic boundary conditions, but infinite in time direction, as shown in figure 2.2. In this way a cylinder topology is defined for space time. The Hilbert space of states in an event $`P`$ is the tensor product
$$=_L_R$$
of R and L space of states. The fact that particles can be classified in left and right does not mean, in general, that the two dynamics are independent, as happens for example in conformal field theory. In what follows, exactly the interacting case will be treated.
Call $`|\alpha _{Li},\alpha _{Ri}`$ the generic vector of a basis of $`_i`$ where $`i=1,\mathrm{},N`$ labels the sites. The notation
$$|\alpha _{2i1},\alpha _{2i}=|\alpha _{Li},\alpha _{Ri}$$
is useful and not ambiguous (even number refers to right, odd number refers to left). The total Hilbert space is:
$$_N=\underset{i=1}{\overset{N}{}}_i$$
and a basic vector can be represented by
$$|\alpha _1,\alpha _2\mathrm{}|\alpha _{2N1},\alpha _{2N}=|\alpha _1,\alpha _2,\mathrm{},\alpha _{2N}_N.$$
Note that in a $`N`$ sites lattice, due to light-cone, $`2N`$ labels are required.
If at a given time $`t`$ there is a line of sites, the particular characteristic of the light-cone is that at time $`t+a/2`$ there is another line of sites, but not equivalent to the previous one, because it is shifted. Only at $`t+a`$ there is an equivalent line (see figure 2.2). Then two different evolution operators can be defined, depending on the initial state:
$$\begin{array}{c}U_+|\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t=|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t+a/2\\ \\ U_{}|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t+a/2=|\alpha _1^{\prime \prime },\alpha _2^{\prime \prime },\mathrm{},\alpha _{2N}^{\prime \prime },t+a\end{array}$$
(2.1)
where the initial states are chosen as in figure 2.3.
Schrödinger form of equations of motion is used, for a state in Hilbert space.<sup>1</sup><sup>1</sup>1 In Schrodinger form, if $`|\alpha `$ is a state vector, its time evolution is given by $`|\alpha ,t=U|\alpha ,0`$ where $`U`$ satisfies motion’s equations: $`U=Te^{i{\displaystyle 𝑑tH}}`$ (Dyson’s series) The global time operator can be chosen as
$$U=U_+U_\text{ }\text{ or }U^{}=U_{}U_+$$
depending on the initial state. For a consistent quantum theory, both this operators must be unitary. This can be guaranteed if the assumption $`U_+^{}U_+=U_{}^{}U_{}=1`$ is made, that is the elementar operators themselves must be unitary.
Another operator plays an important role and is defined as follows (the states are at a certain fixed time):
$$V|\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t=|\alpha _{2N},\alpha _1,\mathrm{},\alpha _{2N1},t$$
(2.2)
it corresponds to an half-space shift in the space direction, with exchange of right and left states(see the figure 2.4).
Two applications of $`V`$ give a shift by an entire lattice spacing, then $`V^2`$ is the lattice space evolution operator. It is possible to write an expression for $`V`$ in terms of the permutation operator:
$$P_{nm}\left(|\alpha _n|\alpha _m\right)=|\alpha _m|\alpha _n,|\alpha _n|\alpha _m_n_m.$$
that looks like:
$$V=P_{12}P_{23}\mathrm{}P_{2N1,\mathrm{\hspace{0.17em}2}N}$$
(2.3)
The verification is direct:
$$\begin{array}{c}P_{12}P_{23}\mathrm{}P_{2N1,\mathrm{\hspace{0.17em}2}N}|\alpha _1,\alpha _2,\mathrm{},\alpha _{2N}=P_{12}P_{23}\mathrm{}P_{2N2,\mathrm{\hspace{0.17em}2}N1}|\alpha _1,\mathrm{},\alpha _{2N2},\alpha _{2N},\alpha _{2N1}=\\ \\ =P_{12}P_{23}\mathrm{}P_{2N3,\mathrm{\hspace{0.17em}2}N2}|\alpha _1,\mathrm{},\alpha _{2N},\alpha _{2N2},\alpha _{2N1}=\mathrm{}=|\alpha _{2N},\alpha _1,\mathrm{},\alpha _{2N1}.\end{array}$$
Also $`V`$ is a unitary operator. To show this, one can use the fact that the permutation is self-adjoint. A more complete list of known properties of $`P`$ is:
$$P=P^1=P^{},P^2=1$$
(2.4)
Then, the adjoint of $`V`$ is:
$$V^{}=P_{2N1,2N}\mathrm{}P_{23}P_{12}$$
and $`VV^{}=1`$ simply follows. The operators defined up to now have the following commutation rules:
$$\begin{array}{cc}[V^2,U_\pm ]=0;U_\pm =VU_{}V^{}.& \end{array}$$
(2.5)
In figure 2.5, the first case is proved, by showing the equivalence of the two paths $`V^2U_+`$ and $`U_+V^2`$. “Mutatis mutandis”, all the other cases can be simply obtained.
Consequently, the two principal evolution operators, $`V^2`$ and $`U`$, are commuting:
$$[V^2,U]=[U,V^2]=0.$$
As previously shown, they are also unitary. Thanks to all these properties, it is very natural to identify them as the exponential of the hamiltonian operator, and the exponential of the linear momentum:
$$\begin{array}{c}U=e^{iaH}\\ V^2=e^{iaP}.\end{array}$$
(2.6)
There are other two important operators, defined as:
$$\begin{array}{c}U_R=U_+V\\ U_L=U_+V^{}\end{array}$$
(2.7)
As clearly shown in figure 2.6 for one of them, they correspond to one step evolution in light-cone directions.
They are commuting and give the expressions:
$$\begin{array}{cc}U=U_RU_L& V^2=U_RU_L^{}\\ [U_R,U_L]=0& U_R^{}U_R=U_L^{}U_L=1\end{array}$$
then, using also (2.6) yields:
$$U_R=e^{i{\displaystyle \frac{a}{2}}\left(H+P\right)},U_L=e^{i{\displaystyle \frac{a}{2}}\left(HP\right)}.$$
(2.8)
### 2.2 Dynamics on light-cone
A dynamics can be defined by giving all the amplitudes of the different processes that can take place. The fundamental assumption is that in every site a whole process can happen, in the sense that if $`|\alpha _L,\alpha _R_{in}`$ and $`|\beta _L,\beta _R_{out}`$ are the incoming and outgoing states in a certain site, they can be considered asymptotic states and the transition amplitude is an S-matrix element:
$${}_{out}{}^{}\beta _R,\beta _L|\alpha _R,\alpha _L_{in}^{}=S_{(\alpha _R,\alpha _L)_{in},(\beta _R,\beta _L)_{out}}$$
(2.9)
This, in general, is an $`mn`$ scattering. The system is then defined via its microscopic amplitudes.
From the definition (2.1) of evolution operators, the following expression holds (the “in” state is at time $`t`$, the “out” state at $`t+a/2`$, with reference at the figure 2.3:
$$\begin{array}{c}\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t\left|U_+\right|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t=_{in}\alpha _1,\alpha _2,\mathrm{},\alpha _{2N}|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{}_{out}=\\ \\ =_{in}\alpha _1,\alpha _2|\alpha _1^{},\alpha _2^{}_{out}\mathrm{}_{in}\alpha _{2N1},\alpha _{2N}|\alpha _{2N1}^{},\alpha _{2N}^{}_{out}=\\ \\ =_{i=1,\mathrm{},N}S_{(\alpha _{2i1},\alpha _{2i})_{in},(\alpha _{2i1}^{},\alpha _{2i}^{})_{out}}^{}\end{array}$$
(2.10)
The operator $`U_{}`$ can be obtained from (2.5). In the last line, the product is on all the sites at a given time.
It is interesting to observe that there is a powerful connection with statistical mechanics, at this point. Assume, for simplicity, that the lattice is euclidean and has a very large (thermodynamics) but finite extension in the “time” direction. In this way there is an initial time (with an initial state $`|\alpha _{in}`$) and a final time (with a final state $`|\beta _{out}`$). The total transition amplitude is given by the product on all sites of the amplitudes, and the sum on all the internal states ($`\gamma ^{},\gamma ^{\prime \prime }`$):
$${}_{out}{}^{}\beta |\alpha _{in}^{}=\underset{intstates}{}\underset{sites}{}S_{\gamma ^{},\gamma ^{\prime \prime }}Z$$
The important fact is that if $`S_{\gamma ^{},\gamma ^{\prime \prime }}`$ is real and positive, it plays the role of a Boltzmann weight for a statistical system whose partition function is given by the previous expression, that was conveniently called $`Z`$. It depends only on the external states. The first appearance of this correspondence is in .
In both the statistical and the particle interpretation, the integrability (that will be always assumed in what follows) can be obtained by requiring that the scattering amplitude $`S_{\gamma ,\gamma ^{}}`$ or the Boltzmann weight satisfy the Yang-Baxter equation (see ). In this case a general $`mn`$ particles scattering can be factorised in the product of “elementary” $`2\mathrm{\hspace{0.17em}2}`$ particles scattering. Then at every site there will be associated one of this “elementary” amplitudes:
$$_{out}\beta _L,\beta _R|\alpha _R,\alpha _L_{in}=S_{\alpha _R,\alpha _L}^{\beta _R,\beta _L}$$
(2.11)
and “in” and “out” states are two particle states (as in the previous equation). In statistical language, usually the Boltzmann weight is indicated with $`R`$ and there is a little change of notation with respect to scattering case:
$$R_{\alpha _L,\alpha _R}^{\beta _R,\beta _L}=S_{\alpha _R,\alpha _L}^{\beta _R,\beta _L}$$
(2.12)
(the lower indices are interchanged). A pictorial interpretation of that is in figure 2.7.
The Yang-Baxter equation takes the form:
$$S_{12}S_{13}S_{23}=S_{23}S_{13}S_{12}$$
(2.13)
(with the indicated change it holds clearly also for R-matrix).
Apparently, this “phenomenological” approach is quite unusual, because in traditional lattice quantum field theory at every site is associated one interaction potential (ex: $`\varphi ^4(i)`$ is the potential on the site $`i`$), not a whole scattering process. This can appear as a sort of “macroscopic” approach, not based on fundamental interactions. But the properties of factorisable scattering must be taken into account. Factorization of generic amplitudes in $`2\mathrm{\hspace{0.17em}2}`$ particles amplitudes is a sort of quantum superposition principle and the remarkable fact is that between one scattering and the other, the particles are asymptotic ones, that means that they are free. For example in figure 2.8 (1+1 QFT on the continuum) between the point 1 and the point 2 the motion is free.
Every point contain all the interaction. This is what was assumed in the definition of the lattice. Then it is perfectly justified that every site is connected with a whole $`2\mathrm{\hspace{0.17em}2}`$ particles scattering process.<sup>2</sup><sup>2</sup>2 In the next paragraphs, it will be explained that, for the particular case of the 6 vertex R matrix, a more traditional lattice QFT approach can be formulated, in terms of a fermionic field. The previous equation (2.10) can be written now in a more specific form:
$$\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t\left|U_+\right|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t=R_{\alpha _1\alpha _2}^{\alpha _1^{}\alpha _2^{}}R_{\alpha _3\alpha _4}^{\alpha _3^{}\alpha _4^{}}\mathrm{}R_{\alpha _{2N1}\alpha _{2N}}^{\alpha _{2N1}^{}\alpha _{2N}^{}}$$
(2.14)
The physical system dynamics is therefore defined by the assignment of an R matrix, and integrability is guaranteed. The consistency on quantum mechanical interpretation is obtained by requiring unitarity, hermitian analyticity and crossing symmetry for the S matrix obtained by (2.12).
It is important to obtain an explicit expression for $`U_L`$ and $`U_R:`$
$$\begin{array}{c}\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t\left|U_R\right|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t=\underset{\stackrel{}{\gamma }}{}\stackrel{}{\alpha },t\left|U_+\right|\stackrel{}{\gamma }\stackrel{}{\gamma }\left|V\right|\stackrel{}{\alpha ^{}}t=\\ \\ =\underset{\stackrel{}{\gamma }}{}R_{\alpha _1\alpha _2}^{\gamma _1\gamma _2}R_{\alpha _3\alpha _4}^{\gamma _3\gamma _4}\mathrm{}\gamma _1,\mathrm{},\gamma _{2N}|\alpha _{2N}^{},\alpha _1^{},\mathrm{},\alpha _{2N1}^{},t=\\ \\ =R_{\alpha _1\alpha _2}^{\alpha _{2N}^{}\alpha _1^{}}R_{\alpha _3\alpha _4}^{\alpha _2^{}\alpha _3^{}}\mathrm{}R_{\alpha _{2N1}\alpha _{2N}}^{\alpha _{2N2}^{}\alpha _{2N1}^{}}\end{array}$$
(2.15)
and, in a similar way,
$$\alpha _1,\alpha _2,\mathrm{},\alpha _{2N},t\left|U_L\right|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{},t=R_{\alpha _1\alpha _2}^{\alpha _2^{}\alpha _3^{}}R_{\alpha _3\alpha _4}^{\alpha _4^{}\alpha _5^{}}\mathrm{}R_{\alpha _{2N1}\alpha _{2N}}^{\alpha _{2N}^{}\alpha _1^{}}$$
(2.16)
This expression is consistent with the figure 2.6. The important fact is that this operators can be expressed in terms of the transfer matrix of an inhomogeneous lattice model, that will be defined in the next section.
### 2.3 Euclidean transfer matrix
Consider a two dimensional euclidean square lattice, with periodic boundary conditions, in both the directions. The links are the physical objects of the system. They can be in different states belonging to the vector spaces $`𝒜`$ and $`𝒱`$. In principle, this vector states can be different, but in the following they will be identified. At every site can be associated a Boltzmann weight depending on the four links crossing at this site and on a spectral parameter $`\lambda `$ (see the figure 2.9).
The simplest case is to take the same Boltzmann weights in all the sites, but more general configurations are possible. In the following an inhomogeneity $`\lambda _i`$ will be assumed, where $`i`$ is the column index (all the sites on a column have the same inhomogeneity), in the sense that the Boltzmann weight in the column $`i`$ is taken to be:
$$w(ab|\alpha \beta ,\lambda \lambda _i).$$
Also for the boundary conditions it is possible to assume more general configurations than the simplest one (i.e. toroidal b.c.). Assume that between the column N and the N+1 (that is 1) there is a nontrivial seam line, in such a way that the Boltzmann weights on the column N (with respect to the normal ones) are given by:
$$w^{(N)}(ab|\alpha \beta ,\lambda )=e^{i\omega b}w(ab|\alpha \beta ,\lambda ).$$
(2.17)
This choice is made because only the link $`b`$ cross the seam line. This are called *twisted boundary conditions* (obviously the choice $`\omega =0`$ reproduces the periodic case).
The partition function for such a model is given by (see the figure 2.9):
$$Z=\underset{linkstates}{}\underset{rows}{}w(ab|\alpha \beta ,\lambda \lambda _1)\mathrm{}w(ab|\alpha \beta ,\lambda \lambda _{N1})w^{(N)}(ab|\alpha \beta ,\lambda \lambda _N)$$
(2.18)
(observe the twisted Boltzmann weight in the last column). It is simple to show that it can be expressed, thanks to periodicity, in terms of a row-to-row transfer matrix. To obtain that, it is convenient to define an operator $`t`$ that acts on both the horizontal and the vertical spaces:
$$\begin{array}{c}t(\lambda ):𝒜𝒜\\ t_{ab}(\lambda )=a\left|t(\lambda )\right|b\end{array}$$
(2.19)
defines the horizontal action; $`t_{ab}`$ depends on the horizontal links and acts on the vertical space as:
$$\begin{array}{c}t_{ab}(\lambda ):𝒱𝒱\\ \alpha \left|t_{ab}(\lambda )\right|\beta =\left[t_{ab}(\lambda )\right]_{\alpha \beta }=w(ab|\alpha \beta ,\lambda ).\end{array}$$
(2.20)
The twist can be obtained with an appropriate operator that changes the horizontal states by a phase after the seam line:
$$|a_{N+n}=e^{i\omega A}|a_n,|a_i𝒜_i$$
(2.21)
where $`\omega `$ is a real parameter (twist), $`A`$ is a self-adjoint operator and $`n`$ is an horizontal index (in the vertical direction there is no twist). $`A`$ must be chosen diagonal on the whole horizontal space $`𝒜,`$ therefore it acts as a phase. Observe that the lattice euclidean structure is a normal periodic one, because the twist acts only on the quantum mechanical structure. For a generic operator $`B_n`$ acting on the site $`n`$ the consistency with (2.21) requires that the twist acts as:
$$B_{N+n}=e^{i\omega A}B_ne^{i\omega A}.$$
(2.22)
This expression can be used for the determination of the $`t_{ab}`$ operator acting on the last column:
$$\begin{array}{c}{}_{N}{}^{}a\left|t^{(N)}(\lambda )\right|b_{N+1}^{}=_Na\left|t^{(N)}(\lambda )e^{i\omega A}\right|b_1=\\ =e_N^{i\omega b}a\left|t^{(N)}(\lambda )\right|b_1\end{array}$$
(2.23)
where the “positions” of the states and operators are emphasized; in particular observe that only $`b`$ crosses the seam line. Using the identification (2.20) the expression (2.17) can be simply obtained.
The row-to-row transfer matrix is an operator $`t^{(N)}(\lambda ,\{\lambda _i\},\omega )`$ that acts on the vertical space:
$$𝒱^{(N)}=\underset{i=1}{\overset{N}{}}𝒱$$
in the following way:
$$t^{(N)}(\lambda ,\{\lambda _i\},\omega )=\underset{a_1}{}T_{a_1a_1}^{(N)}=\underset{a_1,\mathrm{}a_N}{}t_{a_1a_2}(\lambda \lambda _1)t_{a_2a_3}(\lambda \lambda _2)\mathrm{}t_{a_Na_1}(\lambda \lambda _N)e^{i\omega b}$$
(2.24)
The $`T_{ab}`$ is called monodromy matrix. The expression for the partition function is:
$$Z=Tr_{V^{(N)}}\left[t^{(N)}(\lambda ,\{\lambda _i\},\omega )\right]^M.$$
Define now a matrix acting on $`𝒜𝒜`$ given by:
$$R_{a\alpha }^{\beta b}(\lambda )=\left[t_{ab}(\lambda )\right]_{\alpha \beta }=w(ab|\alpha \beta ,\lambda ).$$
(2.25)
(observe that in the operator in (2.20) the Latin and Greek index belong to different spaces, while in R matrix they belong to the same space and this differentiation of notation is redundant). An obvious generalization is needed for the last column. This definition (2.25) shows that the R matrix is defined as a Boltzmann weight, as required in the section 2.2. In the following will be explained the importance of (2.25) for the integrability.
It is possible to show, at this point, that the transfer matrix (2.24) can be used to express the operators defined in (2.15, 2.16). Consider the following form for the inhomogeneity:
$$\lambda _i=(1)^{i+1}\mathrm{\Theta }$$
(2.26)
where $`\mathrm{\Theta }`$ is a positive real number, and calculate a matrix element of the transfer matrix on vectors of $`V^{(2N)}`$ (clearly, in this case, the periodicity requires an even number of horizontal sites):
$$\begin{array}{c}\alpha _1,\alpha _2,\mathrm{},\alpha _{2N}\left|t^{(2N)}(\mathrm{\Theta },\{\lambda _i\})\right|\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2N}^{}=\\ \\ =\underset{a_1\mathrm{}a_{2N}}{}t_{a_1a_2}(0)_{\alpha _1\alpha _1^{}}t_{a_2a_3}(2\mathrm{\Theta })_{\alpha _2\alpha _2^{}}\mathrm{}t_{a_{2N}a_1}(2\mathrm{\Theta })_{\alpha _{2N}\alpha _{2N}^{}}=\\ \\ =R_{\alpha _1\alpha _2}^{\alpha _2^{}\alpha _3^{}}R_{\alpha _3\alpha _4}^{\alpha _4^{}\alpha _5^{}}\mathrm{}R_{\alpha _{2N1}\alpha _{2N}}^{\alpha _{2N}^{}\alpha _1^{}}\end{array}$$
that is exactly the expression (2.16). In an operatorial form this means that:
$$t^{(2N)}(\mathrm{\Theta },\{\lambda _i\})=U_L.$$
(2.27)
A similar picture holds for the operator $`U_R`$, but to take into account this case an assumption on R is necessary. Assume that R matrix is hermitian analytic, that means:
$$R_{cd}^{ab}(\lambda )^{}=R_{ab}^{cd}(\lambda )$$
that in S matrix language is $`S^{}(s)=S(s^{}).`$ This is a well known property of an S matrix and, as explained in the section 2.2, the same requirement is necessary for R. Then, as for $`U_L`$, the following expression can be obtained:
$$t^{(2N)}(\mathrm{\Theta },\{\lambda _i\})^{}=U_R.$$
(2.28)
Here the adjoint conjugation is only on the vertical space.<sup>3</sup><sup>3</sup>3 The definition of adjoint is $`\beta \left|t_{ab}\right|\alpha ^{}=\alpha \left|t_{ab}^{}\right|\beta `$ It is given by:
$$\left[t_{ab}(\lambda )^{}\right]_{\alpha \beta }=w(ab|\beta \alpha ,\lambda )^{}.$$
Transfer matrix can express the evolution operators.
As introduced in the section 2.2, if the R matrix satisfies the Yang-Baxter equation (2.13), then it is possible to show that the transfer matrix commute with itself at different values of the spectral parameter. In this case the system is integrable. The same conclusion holds also for the inhomogeneous and twisted models. Assume that the R matrix satisfies the Yang-Baxter equation, that in the explicit form is:
$$R_{a_2a_1}^{c_1c_2}(\vartheta _1\vartheta _2)R_{a_3c_1}^{b_1c_3}(\vartheta _1\vartheta _3)R_{c_3c_2}^{b_2b_3}(\vartheta _2\vartheta _3)=R_{a_3a_2}^{c_2c_3}(\vartheta _2\vartheta _3)R_{c_3a_1}^{c_1b_3}(\vartheta _1\vartheta _3)R_{c_2c_1}^{b_1b_2}(\vartheta _1\vartheta _2).$$
(2.29)
Using (2.25) it is possible to obtain another form, specific for the $`t`$ operator:
$$R_{ab}^{ef}(\vartheta \vartheta ^{})\left[t_{ec}(\vartheta )\right]_{\alpha \gamma }\left[t_{fd}(\vartheta ^{})\right]_{\gamma \beta }=\left[t_{ae}(\vartheta ^{})\right]_{\alpha \gamma }\left[t_{bf}(\vartheta )\right]_{\gamma \beta }R_{ef}^{cd}(\vartheta \vartheta ^{})$$
(2.30)
For consistency, the R matrix must be of the regular type, that means that at the origin it is the unit operator (or proportional to) on $`𝒜𝒜`$:
$$R_{a\alpha }^{\beta b}(0)=\delta _{a\beta }\delta _{\alpha b}.$$
This equation fixes a normalization for the R matrix. Then, to fit with the equations (2.27, 2.28) a normalization factor is required. It is fixed by the unitarity requirement. The (2.30) holds also if an identical shift is performed on both the spectral parameters:
$$\vartheta \vartheta \alpha $$
(the same for $`\vartheta ^{}`$) because only the difference appears in the R term. This shift is an inhomogeneity for the $`t`$ operator.
Consider now a group of matrices $`g𝒢`$ of dimension $`dim𝒜\times dim𝒜`$ and determinant one, such that they commute with R:
$$[gg,R]=0.$$
(2.31)
Then define a transformed (“gauged” or twisted) $`t`$ operator:
$$t_{(g)ab}=g_{ac}t_{cb}.$$
The twist introduced in (2.23) is the special case of $`𝒢=U(1)^{dim𝒜}.`$ It is simple to verify that also the gauged $`t_{(g)}`$ satisfies the same (2.30). This is because of the commutation relation (2.31). The group $`𝒢`$ is a symmetry group for the Yang-Baxter equation and for the vertex model defined in this way. With some simple algebra, using many times eq. (2.30), it is possible to obtain the form for the monodromy matrix, defined as in (2.24) but with a possibly different gauge $`g_i`$ in every site of a row (all the column has the same gauge) and clearly an inhomogeneity. If the following synthetic notation is introduced, $`G=(g_1,\mathrm{},g_N)`$, then the final result is:
$$R(\vartheta \vartheta ^{})\left[T_G^{(N)}(\vartheta )_𝒜T_G^{(N)}(\vartheta ^{})\right]=\left[T_G^{(N)}(\vartheta ^{})_𝒜T_G^{(N)}(\vartheta )\right]R(\vartheta \vartheta ^{}).$$
Taking the trace on the horizontal space $`𝒜`$ yields:
$$[t_G^{(N)}(\vartheta ,\{\lambda _i)),t_G^{(N)}(\vartheta ^{},\{\lambda _i))]=0$$
(2.32)
(observe that the gauge must be exactly the same). At the various values of $`\vartheta `$, $`t_G^{(N)}(\vartheta ,\{\lambda _i))`$ describes an infinite family of conserved charges. Then the system is integrable. This holds in particular in the inhomogeneous and the twisted cases.
The interesting fact, at this point, is that in some cases the transfer matrix can be exactly diagonalized by Bethe Ansatz methods. This gives an exact expression for eigenstates and eigenvalues of the operator $`U`$.
### 2.4 6 vertex model: main results
The theory developed until now is general and not referred to a specific model. The simplest non trivial case to take into account in the previous framework is the $`4\times 4`$ R matrix, corresponding to the choice $`𝒜=𝒱`$. As shown in , the most general solution is the so called 8 vertex model. This name means that only 8, between the 16 entries of the R matrix, are nonzero. A special case is the 6 vertex model, for which many results have been obtained in the light-cone description: this will be the principal object of this dissertation. The R matrix has the form (lower index are rows and upper index are columns)
$$R(\vartheta ,\gamma )=\left(\begin{array}{cccc}a& & & \\ & c& b& \\ & b& c& \\ & & & a\end{array}\right)$$
(2.33)
As Boltzmann weights, all this nonvanishing entries must be real and nonnegative. But Yang-Baxter equation is solved for all the values of this variables.<sup>4</sup><sup>4</sup>4 Observe that the upper and lower entries are the same; this is the assumption of parity invariance. As explained in (2.2), the scattering interpretation is possible only assuming the properties indicated in that section.
There is a well known mapping between vertex models and spin chains (see ), i.e. the transfer matrix is the exponential of the quantum hamiltonian of the chain:
$$t^{(N)}=e^H.$$
In the case of 8 vertex model, the hamiltonian is the XYZ(1/2) chain, while in the special case of 6 vertex, is the XXZ(1/2) chain. In what follows, this identification can be useful to interpret some facts connected with Bethe Ansatz. The XXZ(1/2) chain hamiltonian is given by:
$$H=\underset{i=1}{\overset{N}{}}\left[\sigma _x^i\sigma _x^{i+1}+\sigma _y^i\sigma _y^{i+1}+(1\mathrm{cos}\gamma )\sigma _z^i\sigma _z^{i+1}\right]$$
and $`\gamma `$ is the anisotropy. The $`\sigma `$ are Pauli matrices. The total z-component of the spin will play an important role in Bethe Ansatz.
This is a statistical approach, but it is possible to give a particle interpretation to the same R matrix (2.33) on the light-cone lattice.
The simplest case of particles obeying Pauli exclusion principle (fermions) and without internal degrees of freedom (color number) is assumed. This means that in an event only one particle of type R and one L at most can take place. In other words, one link has two states: empty or occupied. At every point there are four links, that means 16 possible configurations associated to it. In terms of events, these are the possible configurations that connect a point with the nearest neighbours in the future.
Assume now that only amplitudes that conserves the total number of particles (R+L) are nonvanishing. This reduces to 6 the permitted configurations, as it is shown in the figure 2.10. This is simply the 6 vertex model whose R matrix is written in (2.33).
The assumption of integrability for this amplitudes gives the general six-vertex model. The requirement of symmetry under parity transformation implies that $`\omega _3=\omega _4`$ and $`\omega _5=\omega _6`$. The convention adopted in the figure 2.9 and in (2.20, 2.25, 2.33) shows that
$$b=\omega _3=\omega _4,c=\omega _5=\omega _6.$$
Now, this R matrix can be written in an operatorial form, by defining a lattice chiral fermion $`\psi _{A,n}`$, with $`A=R,L`$, and $`n`$ labels the sites. The anticommutation rules are the canonical ones:
$$\{\psi _{A,n},\psi _{B,m}\}=0,\{\psi _{A,n},\psi _{B,m}^{}\}=\delta _{AB}\delta _{nm}.$$
(2.34)
This fermion has some interesting properties, that are exposed in the paper , and are sketched in the following list:
1. the R matrix and all the other operators $`U_{any}`$ can be written in an operatorial form in terms of the fermion;
2. the lattice hamiltonian, in the free case $`\omega _1=\omega _2=b=1,c=0`$, can be explicitly written; by this, the following dispersion relation can be obtained: $`E=\pm k`$; this is the dispersion relation for a free massless particle; the unusual fact is that it is monotonous, then there is no doubling of fermions; this is a consequence of the nonlocality of the hamiltonian
3. the lattice hamiltonian admits a continuum limit $`N\mathrm{}`$ and $`a\mathrm{\hspace{0.17em}0}`$ but with $`L=Na`$ fixed; the locality is recovered in this limit; the continuum equations of motion are those of the massive Thirring model. The space is compactified on a cylinder.
In the continuum limit, the massive Thirring model emerges as the field theory characterising the scaling behavior of the dynamics on the lattice (remember that $`L`$ is finite; the scaling behaviour is understood in terms of this $`L`$).
### 2.5 6 vertex model: Bethe Ansatz
The assumptions of unitarity and hermitian analyticity will be taken into account, for the R matrix, and this requires that the variables in (2.33) must have the specific form:
$$a=a(\vartheta ,\gamma )=\mathrm{sinh}(i\gamma \vartheta ),b=b(\vartheta ,\gamma )=\mathrm{sinh}\vartheta ,c=c(\vartheta ,\gamma )=i\mathrm{sin}\gamma $$
(2.35)
The transfer matrix defined by this R matrix (2.24, 2.25) can be diagonalized with Bethe Ansatz method. In terms of the spin chain, this means that there are two operators, usually indicated by $`B(\vartheta )`$ and $`C(\vartheta )`$, whose expression is known, and there is a “reference state”<sup>5</sup><sup>5</sup>5 The “reference state” at this point is only a mathematical object. Physically speaking, it corresponds to the ferromagnetic state with all the spins up. $`|\mathrm{\Omega }`$ such that:
$$C(\vartheta )|\mathrm{\Omega }=0$$
and
$$B(\vartheta _1)\mathrm{}B(\vartheta _M)|\mathrm{\Omega },$$
(2.36)
for appropriate values of $`\vartheta _j`$ is an eigenstate of the transfer matrix. The “appropriate values” of $`\vartheta _j`$ can be obtained as the solution of a set of $`M`$ coupled nonlinear equations, that are called Bethe Ansatz equations. In general, because of (2.32), the transfer matrix contains all the conserved charges, in particular the hamiltonian. Then the Bethe Ansatz eigenstates are also eigenstates of the hamiltonian. The Hilbert space of the theory and the action of conserved charges on it are then perfectly know.
All this computations for the 6 vertex model were obtained in ; the final results are written here for the eigenvalues of the inhomogeneous and twisted transfer matrix:
$$\begin{array}{c}\tau (\vartheta ,\mathrm{\Theta },\omega )=e^{i\omega }\left[a(\vartheta \mathrm{\Theta })a(\vartheta +\mathrm{\Theta })\right]^N\underset{j=1}{\overset{M}{}}\frac{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{\pi }{2}}+\vartheta _j+\vartheta \right]}{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{\pi }{2}}\vartheta _j\vartheta \right]}+\\ +e^{i\omega }\left[b(\vartheta \mathrm{\Theta })b(\vartheta +\mathrm{\Theta })\right]^N_{j=1}^M\frac{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{3\pi }{2}}\vartheta _j\vartheta \right]}{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{\pi }{2}}+\vartheta _j+\vartheta \right]}\end{array}$$
and the values of $`\vartheta _j`$ are defined by the set of coupled nonlinear equations called Bethe Ansatz equations:
$$\left(\frac{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j+\mathrm{\Theta }+{\displaystyle \frac{i\pi }{2}}\right]\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j\mathrm{\Theta }+{\displaystyle \frac{i\pi }{2}}\right]}{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j+\mathrm{\Theta }{\displaystyle \frac{i\pi }{2}}\right]\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j\mathrm{\Theta }{\displaystyle \frac{i\pi }{2}}\right]}\right)^N=e^{2i\omega }\underset{k=1}{\overset{M}{}}\frac{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j\vartheta _k+i\pi \right]}{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[\vartheta _j\vartheta _ki\pi \right]}$$
(2.37)
where $`2N`$ is the length of the chain and $`N`$ the number of sites in a row of the light-cone lattice. The $`\vartheta _j`$ are *called Bethe roots*, and in principle, can take any complex value. But there is a periodicity in their imaginary part:
$$\vartheta _j\vartheta _j+\frac{\pi ^2}{\gamma }i$$
(2.38)
then only a strip around the real axis must be taken into account for the Bethe roots:
$$\vartheta _j\times i]\frac{\pi ^2}{2\gamma },\frac{\pi ^2}{2\gamma }].$$
(2.39)
Moreover, only the range
$$0<\gamma <\pi $$
will be examined.
From Bethe Ansatz it is known that the whole spectrum of the theory can be obtained using all the Bethe configurations having $`MN`$ and $`\vartheta _j\vartheta _k`$ for every $`jk.`$ In general for a state with $`M`$ roots the third component of the spin of the chain is
$$S=NM,$$
(2.40)
because every operator $`B(\vartheta _i)`$ counts as $`1`$ spin.
This XXZ(1/2) chain has 2 states in every site, then $`2^{2N}`$ states. Then the energy spectrum is upper and lower bounded. Changing the sign of the hamiltonian gives another permitted physical system. This means that there are two possible vacua. The first one is the so called ferromagnetic ground state, corresponding to $`M=0`$ that is the reference state $`|\mathrm{\Omega }.`$ It has spin $`S=N`$. The second one is the antiferromagnetic ground state, that can be obtained with $`M=N`$ and all the roots $`\vartheta _i`$ real. It has spin $`S=0`$.
In what follows, only the antiferromagnetic ground state will be considered, because it has one important property: in the thermodynamic limit ($`N\mathrm{}`$) it can be interpreted as a Dirac vacuum (a sea of particles created by $`B`$) and the excitations on this vacuum behave as particles.
The energy $`E`$ and momentum $`P`$ of a state can be read out by the transfer matrix eigenvalues by use of the equations (2.27, 2.28). The final form is:
$$e^{i{\displaystyle \frac{a}{2}}\left(E\pm P\right)}=e^{\pm i\omega }\underset{j=1}{\overset{M}{}}\frac{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{\pi }{2}}\mathrm{\Theta }\pm \vartheta _j\right]}{\mathrm{sinh}{\displaystyle \frac{\gamma }{\pi }}\left[i{\displaystyle \frac{\pi }{2}}+\mathrm{\Theta }\vartheta _j\right]}.$$
(2.41)
Other integrals of motion can be obtained in a similar way by transfer matrix (see (2.32)). Observe that the second term in the transfer matrix expression vanishes because $`b(0)=0`$.
For future analysis, it is more convenient to express the coupling constant $`\gamma `$ in terms of a different variable $`p`$:
$$p=\frac{\pi }{\gamma }1,0<p<\mathrm{}.$$
Then all the expressions will be in terms of $`p.`$
## Chapter 3 A NONLINEAR EQUATION FOR BETHE ANSATZ
In this chapter the fundamental nonlinear integral equation (on the lattice) for the Bethe Ansatz is derived. In literature it is known as Destri-de Vega equation. It was obtained first in and, in the final form, in . On this equation the continuum limit is performed.
### 3.1 Counting function
It is possible to write the Bethe equations (2.37) in terms of a counting function, that will be called $`Z_N(\vartheta ).`$ To do that, one important preliminar definition is needed:
$$\varphi (\vartheta ,\nu )=i\mathrm{log}\frac{\mathrm{sinh}\frac{1}{p+1}(i\pi \nu +\vartheta )}{\mathrm{sinh}\frac{1}{p+1}(i\pi \nu \vartheta )},\varphi (\vartheta ,\nu )=\varphi (\vartheta ,\nu )$$
where the primary interest is in values $`\nu =1/2,\mathrm{\hspace{0.17em}1}`$ and the oddness on the analyticity strip around the real axis defines a precise logarithmic branch: the so called fundamental branch. In the appendix B the complete structure of the cuts and all the relevant properties of this function are exposed. The counting function<sup>1</sup><sup>1</sup>1 The name counting function will became clear in the next paragraphs, where the counting equation will be obtained. can be defined by
$$Z_N(\vartheta )=N[\varphi (\vartheta +\mathrm{\Theta },\frac{1}{2})+\varphi (\vartheta \mathrm{\Theta },\frac{1}{2})]\underset{k=1}{\overset{M}{}}\varphi (\vartheta \vartheta _k,1)+2\omega $$
(3.1)
used to express the logarithm of the Bethe equations, one obtains simply:
$$Z_N(\vartheta _j)=2\pi I_j,I_j+\frac{1+\delta }{2},\delta =(M)_{mod\mathrm{\hspace{0.17em}2}}=(NS)_{mod\mathrm{\hspace{0.17em}2}}\{0,1\}.$$
(3.2)
The number $`I_j`$ plays the role of a quantum number for the Bethe basic vectors (2.36) and it must be chosen depending on the value of $`\delta .`$ Notice that $`\delta `$ and $`\omega `$ play a similar role, because both produce a shift in the quantum numbers $`I_j`$ (if $`\omega `$ is absorbed in the definition of $`I_j`$): in the first case the shift is exactly $`\pi ,`$ in the second case it is a real (possibly complex) number. This means that the variable $`\delta `$ can be absorbed in $`\omega `$ but the most convenient choice is to use them both.
Observe that Bethe roots can be obtained as zeros of the equation:
$$1+(1)^\delta e^{iZ_N(\vartheta _j)}=0$$
(3.3)
### 3.2 Classification of Bethe roots
From Bethe Ansatz it is known that a Bethe state (2.36) is uniquely characterized by the set of quantum numbers $`\left\{I_j\right\}_{j=1,\mathrm{},M},MN`$ that appear in (3.2). Notice that $`MN`$ means $`S0.`$ The values of $`\vartheta _j`$ to put in (2.36) can be obtained solving Bethe equations. It is also known that only states with
$$\vartheta _j\vartheta _i,ji$$
are necessary to form a basis for the space of states. It is a sort of fermionic character for Bethe states .
As usual in many cases, Bethe roots can be either real or in complex conjugate pairs (for large $`N`$). In the specific case (2.37), there is another possibility, due to periodicity (2.38): if a complex solution has imaginary part $`\mathrm{}m\vartheta =\frac{\pi }{2}(p+1)`$ it appears as a single (in (2.37) it produces the left hand side real, then its complex conjugate is not required). It is called *self-conjugate root*. Remember now that the maximal number of real roots ($`M=N`$) describes the antiferromagnetic ground state. From the point of view of the counting function, a more precise classification of roots is required:
* *real roots*; they are real solutions of (3.2) that appear in the vector (2.36); their number is $`M_R`$;
* *holes*; real solutions of (3.2) that do NOT appear in the vector (2.36); their number is $`N_H`$;
* *special roots/holes* (special objects); they are real roots or holes where the derivative $`Z_N^{}(\vartheta _j)`$ is negative<sup>2</sup><sup>2</sup>2 The characteristics of this type of solutions will be completely clarified in the next sections. ; their number is $`N_S`$; they must be counted both as “normal” and as “special” objects;
* *close pairs*; complex conjugate solutions with imaginary part in the range $`0<|\mathrm{}m\vartheta |<\pi \mathrm{min}(1,p)`$; their number is $`M_C`$;
* *wide roots in pairs*; complex conjugate solutions with imaginary part in the range $`\pi \mathrm{min}(1,p)<|\mathrm{}m\vartheta |<\pi \frac{p+1}{2}`$;
* *self-conjugate roots* (wide roots appearing as single); $`\mathrm{}m\vartheta =\pi \frac{p+1}{2}`$; their number is $`M_{SC}`$.
The total number of wide roots appearing in pairs or single is $`M_W`$. The following notation will be used (sometimes) for later convenience, to indicate the position of the solutions: $`h_j`$ for holes, $`y_j`$ for special objects, $`c_j`$ for close roots, $`w_j`$ for wide roots.
Complex roots with imaginary part larger than the self-conjugates are not required because of the periodicity of Bethe equations. This classification is not at all academic; it will play an important role in the physical interpretation of the final equation. A graphical representation of the various types of solutions is given in figure 3.1.
An important remark must be made: from the definition itself of $`Z_N`$ (3.1) it is obvious that only for states without complex roots the fundamental strip for $`\varphi `$ (see (B.3)), that is the largest strip around the real axis without singularities, is the fundamental strip for $`Z_N.`$ In all the other cases the analyticity strip for $`Z_N`$ is narrower, and depends on the imaginary parts of the complex roots.
An important property follows from this classification: the $`Z_N`$ function is *real analytic* if $`\omega `$ is a real number
$$Z_N\left(\vartheta ^{}\right)=\left(Z_N(\vartheta )\right)^{}$$
(3.4)
### 3.3 Counting equation
It is possible to obtain an equation relating the numbers of all the various types of solutions. The path is simple, and makes use of the asymptotic values obtained in (B.5) to calculate the limits $`\vartheta \pm \mathrm{}`$ in $`Z_N`$ (3.1). Observe that the term $`\varphi (\vartheta \vartheta _j,1)`$, in the case of the wide roots, takes contributions from the horizontal strips next to the fundamental, as in figure B.1. This contributions depend from $`M_W`$ (number of wide roots below the real axis) and $`M_W`$ (number of wide roots over the real axis). The asymptotic limits are then:
$$\begin{array}{c}Z_N(+\mathrm{})=N\pi +\pi \frac{p1}{p+1}S+2\pi \mathrm{sign}(p1)M_W+2\omega \\ \\ Z_N(\mathrm{})=N\pi \pi \frac{p1}{p+1}S2\pi \mathrm{sign}(p1)M_W+2\omega \end{array}$$
(3.5)
Observe that the number of self-conjugate roots is: $`M_WM_W=M_{SC}.`$ On the lattice, all the values for $`N,S`$ are permitted. But in view of the continuum limit, $`N`$ is expected to be larger than $`S`$: $`NS`$ (also the number of holes and complex roots is expected to be much smaller than $`N`$). Then
$$Z_N(+\mathrm{})>Z_N(\mathrm{}).$$
The function is globally increasing (but not necessarily monotonous: in some points it can be a decreasing function). Define now two variables $`\zeta _\pm `$, choosing $`k_\pm `$ in such a way that they take value in the indicated interval:
$$\zeta _\pm =\pm Z_N(\pm \mathrm{})+\pi \delta 2\pi k_\pm ,\pi <\zeta _\pm \pi $$
Using the previous expressions for $`Z_N`$ the result is:
$$\zeta _\pm =2\pi \left(\frac{S}{p+1}\pm \frac{\omega }{\pi }+\frac{1}{2}+\frac{S}{p+1}\frac{\omega }{\pi }\right)$$
where the new symbol $`x`$ is the integer part of $`x`$, that is the largest integer smaller or equal to $`x`$. In this way, the following expression holds:
$$\begin{array}{c}Z_N(+\mathrm{})=2\pi I_{max}+\pi +\zeta _+\\ Z_N(\mathrm{})=2\pi I_{min}\pi \zeta _{}\end{array}$$
(3.6)
where $`I_{max}`$ is the largest quantum number satisfying $`I_{max}<{\displaystyle \frac{Z_N(+\mathrm{})}{2\pi }}`$ and $`I_{min}`$ is the smallest quantum number satisfying $`I_{min}>{\displaystyle \frac{Z_N(\mathrm{})}{2\pi }}`$. This very precise definition is required to take into account for the special roots/holes that can appear in the tails (this can happen when $`Z_N^{}<0`$ for very large values of $`|\vartheta |`$).
The total number of real solutions is the sum of real roots and holes, and in terms of $`I_{max},I_{min}`$ it can be written as:
$$M_R+N_H=I_{max}I_{min}+1+2N_S.$$
Using the equations (2.40, 3.6) and the expression (3.5) the result is:
$$\begin{array}{c}N_H2N_S=2S+M_C+2\theta (p1)M_W+\\ \\ \frac{1}{2}+\frac{S}{p+1}+\frac{\omega }{\pi }\frac{1}{2}+\frac{S}{p+1}\frac{\omega }{\pi }.\end{array}$$
(3.7)
It is the so called *lattice counting equation*. Remember that $`S`$ is a nonnegative integer. In the case of $`\omega =0`$, it turns out that $`N_H`$ is even (remember that $`M_C`$ is the number of close roots, and is even).
The most important fact is that in this equation doesn’t appear the number of real roots. This fact, together to what will be explained in the next paragraph, allows to consider the real roots as a sea of particles (Dirac vacuum) and all other types of solutions (holes, complex) as excitations on this sea.
### 3.4 Non linear integral equation (I)
In this section an equation generating $`Z_N`$ will be obtained. The counting function (3.1) can be written in the following way:
$$\begin{array}{c}Z_N(\vartheta )=N\left[\varphi (\vartheta +\mathrm{\Theta },\frac{1}{2})+\varphi (\vartheta \mathrm{\Theta },\frac{1}{2})\right]+\underset{k=1}{\overset{N_H}{}}\varphi (\vartheta h_k,1)+\\ \underset{k=1}{\overset{M_C+M_W}{}}\varphi (\vartheta \xi _k,1)\underset{k=1}{\overset{M_R+N_H}{}}\varphi (\vartheta x_k,1)\end{array}$$
(3.8)
In this case $`\xi _k`$ collects close and wide roots, and $`x_k`$ the real roots and holes. Now, it is convenient to deal with the derivative of this expression:
$$\begin{array}{c}Z_N^{}(\vartheta )=N\left[\varphi ^{}(\vartheta +\mathrm{\Theta },\frac{1}{2})+\varphi ^{}(\vartheta \mathrm{\Theta },\frac{1}{2})\right]+\underset{k=1}{\overset{N_H}{}}\varphi ^{}(\vartheta h_k,1)+\\ \underset{k=1}{\overset{M_C+M_W}{}}\varphi ^{}(\vartheta \xi _k,1)\underset{k=1}{\overset{M_R+N_H}{}}\varphi ^{}(\vartheta x_k,1)\end{array}$$
(3.9)
Both the previous equations hold for $`\vartheta `$. Let $`\widehat{x}`$ be a real solution of the Bethe equation. Assume $`Z__N^{}(\widehat{x})0`$ and define a complex neighbour of $`\widehat{x}`$ by a small $`|\nu |1`$: $`\mu =\widehat{x}+\nu `$. Clearly $`(1)^\delta e^{iZ_N(\widehat{x})}=1`$. Consider the expression:
$$1+(1)^\delta e^{iZ_N(\widehat{x}+\nu )}1+(1)^\delta e^{iZ_N(\widehat{x})+i\nu Z_N^{}(\widehat{x})}i\nu Z_N^{}(\widehat{x})$$
(3.10)
then the following identity holds:
$$\frac{1}{\mu \widehat{x}}=\frac{1}{\nu }=\frac{(1)^\delta e^{iZ_N(\widehat{x}+\nu )}iZ_N^{}(\widehat{x}+\nu )}{1+(1)^\delta e^{iZ_N(\widehat{x}+\nu )}}+\mathrm{}$$
(3.11)
(the dots are regular terms in $`\mu \widehat{x}`$). From the Cauchy theorem and from (3.11), given an analytic function $`f(x)`$ on an appropriate strip containing the real axis, yields:
$$f(\widehat{x})=_{\mathrm{\Gamma }_{\widehat{x}}}\frac{d\mu }{2\pi i}\frac{f(\mu )}{\mu \widehat{x}}=_{\mathrm{\Gamma }_{\widehat{x}}}\frac{d\mu }{2\pi i}f(\mu )\frac{(1)^\delta e^{iZ_N(\mu )}iZ_N^{}(\mu )}{1+(1)^\delta e^{iZ_N(\mu )}}$$
(3.12)
where $`\mathrm{\Gamma }_{\widehat{x}}`$ is a anti-clockwise curve encircling $`\widehat{x}`$ and avoiding other singularities of the denominator, i.e. other Bethe solutions (real or complex). It is always possible to find such a closed curve, because Bethe solutions are finite in number. An equation like (3.12) can be written for all the real roots of (3.2), $`x_k,k`$. The derivative $`\varphi ^{}(\vartheta ,1)`$ is an analytic function, if poles are avoided, and applying to that the expression (3.12), the last sum in (3.9) becomes:
$$\begin{array}{c}\underset{k=1}{\overset{M_R+N_H}{}}\varphi ^{}(\vartheta x_k,1)=\underset{k=1}{\overset{M_R+N_H}{}}_{\mathrm{\Gamma }_{x_k}}\frac{d\mu }{2\pi i}\varphi ^{}(\vartheta \mu ,1)\frac{(1)^\delta e^{iZ_N(\mu )}iZ_N^{}(\mu )}{1+(1)^\delta e^{iZ_N(\mu )}}=\\ =_\mathrm{\Gamma }\frac{d\mu }{2\pi i}\varphi ^{}(\vartheta \mu ,1)\frac{(1)^\delta e^{iZ_N(\mu )}iZ_N^{}(\mu )}{1+(1)^\delta e^{iZ_N(\mu )}}\end{array}$$
(3.13)
The sum on the contours was modified to a unique curve $`\mathrm{\Gamma }`$ encircling all the real solutions $`x_k`$, and avoiding the complex Bethe solutions (this is possible because they are finite in number), as in the figure 3.2.
Clearly the $`\mathrm{\Gamma }`$ curve must be contained in the strip
$$0<\eta _+,\eta _{}<\mathrm{min}\{\pi ,\pi p,|\mathrm{}mc_k|k\}$$
Without loss of generality, assume that $`\eta _+=\eta _{}=\eta `$, and deform $`\mathrm{\Gamma }`$ to the contour of the strip characterized by $`\eta `$. The regions at $`\pm \mathrm{}`$ do no contribute because of the vanishing of $`\varphi ^{}`$, then the integral can be performed only on the lines $`\mu =x\pm i\eta `$, where $`x`$ is real<sup>3</sup><sup>3</sup>3 $`=_{\mathrm{}}^+\mathrm{}`$ In all the cases where the limits of integration are not explicitly specified, the integral is taken on the whole real axis. :
$$\begin{array}{c}_\mathrm{\Gamma }\frac{d\mu }{2\pi i}\varphi ^{}(\vartheta \mu ,1)\frac{(1)^\delta e^{iZ_N(\mu )}iZ_N^{}(\mu )}{1+(1)^\delta e^{iZ_N(\mu )}}=\\ \\ =\frac{dx}{2\pi i}\varphi ^{}(\vartheta x+i\eta ,1)\frac{(1)^\delta e^{iZ_N(xi\eta )}iZ_N^{}(xi\eta )}{1+(1)^\delta e^{iZ_N(xi\eta )}}+\\ \\ \frac{dx}{2\pi i}\varphi ^{}(\vartheta xi\eta ,1)\frac{(1)^\delta e^{iZ_N(x+i\eta )}iZ_N^{}(x+i\eta )}{1+(1)^\delta e^{iZ_N(x+i\eta )}}\end{array}$$
(3.14)
The first term at the right can be written as:
$$\begin{array}{c}\frac{(1)^\delta e^{iZ_N(xi\eta )}iZ_N^{}(xi\eta )}{1+(1)^\delta e^{iZ_N(xi\eta )}}=\\ \frac{(1)^\delta e^{iZ_N(xi\eta )}iZ_N^{}(xi\eta )}{1+(1)^\delta e^{iZ_N(xi\eta )}}+iZ_N^{}(xi\eta )\end{array}$$
(3.15)
Putting this in (3.15) the integral contribution from $`Z_N^{}(xi\eta )`$ is independent from $`\eta `$ because no singularities are crossed if $`\eta \mathrm{\hspace{0.17em}0}`$ and at the infinity the integrand is zero. The following convenient redefinition can be introduced (see (B.4)):
$$K(x)\frac{\varphi ^{}(x,1)}{2\pi }=\frac{1}{\pi (p+1)}\frac{\mathrm{sin}{\displaystyle \frac{2\pi }{p+1}}}{\mathrm{cosh}{\displaystyle \frac{2\vartheta }{p+1}}\mathrm{cos}{\displaystyle \frac{2\pi }{p+1}}}$$
(3.16)
In general, for every complex value of $`\vartheta `$:
$$\begin{array}{c}𝑑x\left[\delta (x)Z_N^{}(\vartheta x)+K(\vartheta x)Z_N^{}(x)\right]=N\left[\varphi ^{}(\vartheta +\mathrm{\Theta },\frac{1}{2})+\varphi ^{}(\vartheta \mathrm{\Theta },\frac{1}{2})\right]+\\ +\underset{k=1}{\overset{N_H}{}}\varphi ^{}(\vartheta h_k,1)\underset{k=1}{\overset{M_C+M_W}{}}\varphi ^{}(\vartheta \xi _k,1)+\\ \frac{dx}{2\pi i}\varphi ^{}(\vartheta x+i\eta ,1)\frac{(1)^\delta e^{iZ_N(xi\eta )}iZ_N^{}(xi\eta )}{1+(1)^\delta e^{iZ_N(xi\eta )}}+\\ +\frac{dx}{2\pi i}\varphi ^{}(\vartheta xi\eta ,1)\frac{(1)^\delta e^{iZ_N(x+i\eta )}iZ_N^{}(x+i\eta )}{1+(1)^\delta e^{iZ_N(x+i\eta )}}\end{array}$$
(3.17)
The second term on the left takes the form
$$𝑑xK(\vartheta x,1)Z_N^{}(x)=𝑑xK(x,1)Z_N^{}(\vartheta x);$$
that can be obtained by shifting the integration line: $`xx+\mathrm{}m\vartheta `$. But if poles are crossed (this can happen if the imaginary part is sufficiently large), their contribution must appear in the right hand side of the equation.
Using Fourier transformations, the convolution in the left hand side of (3.17) can be simply handled. Observe first that the Fourier transform of $`\left[\delta (\vartheta x)+K(\vartheta x)\right]`$ is simply $`1+\stackrel{~}{K}`$. In the appendix B the form of $`\stackrel{~}{K}`$ is given (as $`\stackrel{~}{\varphi ^{}}`$ in (B.6)). It is obvious that $`1+\stackrel{~}{K}`$ is nonvanishing, so it can be reversed. Call $`\stackrel{~}{\mathrm{\Delta }}`$ the inverse:
$$\stackrel{~}{\mathrm{\Delta }}(k)\frac{1}{1+\stackrel{~}{K}(k)}.$$
The inverse Fourier transform of $`\stackrel{~}{\mathrm{\Delta }}`$, indicated as $`\mathrm{\Delta }(x)`$, is a distribution, because $`\stackrel{~}{\mathrm{\Delta }}(k)\mathrm{\hspace{0.17em}1}`$ for $`k\pm \mathrm{}.`$ The following obvious equation holds:
$$𝑑x\mathrm{\Delta }(x)\left(\delta (\theta x)+K(\vartheta x)\right)=\delta (\theta )$$
and this allows to invert the convolution in (3.17). Calling globally $`(\vartheta )`$ the right hand side of (3.17) the inversion is:
$$Z_N^{}(\vartheta )=𝑑x\mathrm{\Delta }(\vartheta x)(x).$$
(3.18)
The various terms in $``$ give different contributions, that shall be computed in the following.
For the two terms $`\varphi ^{}(\vartheta \pm \mathrm{\Theta },{\displaystyle \frac{1}{2}})`$ the convolution can be completely performed using (A.3):
$$𝑑x\mathrm{\Delta }(\vartheta x)\varphi ^{}(x\pm \mathrm{\Theta },\frac{1}{2})=\frac{1}{\mathrm{cosh}(\vartheta \pm \mathrm{\Theta })}.$$
(3.19)
The effect of the inversion on $`\varphi ^{}(\vartheta ,1)`$ gives an important object, that will be called $`G:`$
$$𝑑x\mathrm{\Delta }(\vartheta x)\frac{\varphi ^{}(x,1)}{2\pi }=𝑑x\mathrm{\Delta }(\vartheta x)K(x)G(\vartheta )$$
(3.20)
Using Fourier transforms (see A.3) for $`\mathrm{\Delta }`$ and $`\varphi ^{}`$ (B.6) the following expression can be obtained:
$$\begin{array}{c}G(\theta )=\frac{1}{p+1}\frac{dk}{2\pi }e^{ik\theta {\displaystyle \frac{1}{p+1}}}\stackrel{~}{K}(k)\frac{1}{1+\stackrel{~}{K}(k)}=\\ =\frac{1}{2\pi }𝑑ke^{ik\theta }\frac{\mathrm{sinh}{\displaystyle \frac{\pi (p1)k}{2}}}{2\mathrm{sinh}{\displaystyle \frac{\pi pk}{2}}\mathrm{cosh}{\displaystyle \frac{\pi k}{2}}}.\end{array}$$
(3.21)
The expression (B.6) used in the previous computation for the Fourier transform of $`\varphi ^{}`$ holds only for $`\theta `$ in the fundamental strip (B.3), i.e. for real solutions or close pairs, as in section 3.2. This means that contributions to $`Z_N^{}`$ coming from wide roots require a different approach (see later). Instead the terms containing holes and close roots are completely arranged in this way.
The Fourier transform of (3.21) is obviously
$$\stackrel{~}{G}(k)=\frac{\stackrel{~}{K}(k)}{1+\stackrel{~}{K}(k)}.$$
This is an even function, real on the real axis (*real analyticity*), positive if $`p>1`$ and negative in the opposite case. The same properties, clearly, hold for $`G(\theta )`$. The function $`\stackrel{~}{G}`$ vanishes exponentially for large values of $`|k|`$:
$$\stackrel{~}{G}e^{\frac{\pi \mathrm{min}(1,p)}{p+1}|k|}.$$
(consequently also the function $`G(\theta )`$ has a similar asymptotic behaviour). Because of this, the Fourier transform in (3.21) at the points
$$|\mathrm{}m\theta |=\pm \pi \mathrm{min}(1,p)$$
has a singularity. As before, this points correspond exactly to the limit between close roots and wide roots. This confirms exactly that wide roots require a different analysis. Call $`source`$ the contribution to (3.18) by one wide root put in $`w`$:
$$\begin{array}{c}source=𝑑x\mathrm{\Delta }(x)\varphi ^{}(\vartheta xw,1)=\\ =\frac{1}{p+1}𝑑x\frac{dk}{2\pi }e^{ikx{\displaystyle \frac{1}{p+1}}}\stackrel{~}{\mathrm{\Delta }}(k)\varphi ^{}(\vartheta xw,1).\end{array}$$
Using the known expression (B.4) for $`\varphi ^{}`$ the computation can be performed into a closed form. The integral in $`x`$ must be performed first. It is a Fourier transform, but because of the $`w`$ contribution, it is different from the Fourier transform of $`\varphi ^{}`$ calculated in (B.6), that holds for $`\vartheta `$ near the real axis. This is because $`\mathrm{}mw`$ is larger than the position of the first singularity of $`\varphi ^{}`$, and the poles to take into account in the two cases are different. The result can be written in the following form:
$$source=2\pi G_{II}(\vartheta w)$$
(3.22)
where the function $`G_{II}`$ is defined by (for $`|\mathrm{}m\theta |>\pi \mathrm{min}(1,p)`$):
$$2\pi G_{II}(\theta )=\{\begin{array}{c}\frac{i}{p}\left[\mathrm{coth}\frac{\theta }{p}\mathrm{sign}\mathrm{}m(\theta )+\mathrm{coth}\frac{(i\pi \theta \mathrm{sign}\mathrm{}m(\theta ))}{p}\right]\mathrm{if}p>1\\ \\ \mathrm{sign}\mathrm{}m(\theta )i\left[\frac{1}{\mathrm{sinh}\theta }+\frac{1}{\mathrm{sinh}(\theta i\pi p\mathrm{sign}\mathrm{}m(\theta ))}\right]\mathrm{if}p<1\end{array}$$
(3.23)
A quite singular property, pointed out in , is that this function can be interpreted as the so called *second determination* of the previous function $`G`$. The general definition of the second determination is:
$$f_{II}(\theta )=\{\begin{array}{c}f(\theta )+f(\theta i\pi \mathrm{sign}\mathrm{}m(\theta ))\mathrm{if}p>1\\ \\ f(\theta )f(\theta i\pi p\mathrm{sign}\mathrm{}m(\theta ))\mathrm{if}p<1\end{array}\text{ for }|\mathrm{}m\theta |>\pi \mathrm{min}(1,p)$$
(3.24)
Applied to (3.21), it yields exactly (3.22).
Taking into account (3.19, 3.20, 3.22) in (3.18), the equation (3.17) takes the form of an integro-differential equation for $`Z_N^{}`$:
$$\begin{array}{c}Z_N^{}(\vartheta )=N\left[\frac{1}{\mathrm{cosh}(\vartheta +\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(\vartheta \mathrm{\Theta })}\right]+\underset{k=1}{\overset{N_H}{}}2\pi G(\vartheta h_k)+\\ \underset{k=1}{\overset{M_C}{}}2\pi G(\vartheta c_k)\underset{k=1}{\overset{M_W}{}}2\pi G_{II}(\vartheta w_k)+\\ +\frac{1}{i}𝑑xG(\vartheta xi\eta )\frac{(1)^\delta e^{iZ_N(x+i\eta )}iZ_N^{}(x+i\eta )}{1+(1)^\delta e^{iZ_N(x+i\eta )}}+\\ \frac{1}{i}𝑑xG(\vartheta x+i\eta )\frac{(1)^\delta e^{iZ_N(xi\eta )}iZ_N^{}(xi\eta )}{1+(1)^\delta e^{iZ_N(xi\eta )}}\end{array}$$
(3.25)
This equation holds for both the cases of $`\omega `$ real or complex. The most common case is $`\omega `$. This allows to write the last two lines in a more compact form, valid only if $`\vartheta `$ is on the real axis. The real analyticity of $`Z_N`$ and of $`G`$ is used:
$$2\mathrm{}m𝑑xG(\vartheta xi\eta )\frac{(1)^\delta e^{iZ_N(x+i\eta )}iZ_N^{}(x+i\eta )}{1+(1)^\delta e^{iZ_N(x+i\eta )}}\text{if }\vartheta .$$
Observe that the second factor in the integral terms seems to be the derivative of $`\mathrm{log}\left[1+(1)^\delta e^{\pm iZ_N(x\pm i\eta )}\right]`$, but this substitution requires some care, because of the polydromy of the logarithmic function. Call, for the sake of simplicity, the argument of the $`\mathrm{log}`$:
$$f(x)=1+(1)^\delta e^{iZ_N(x+i\eta )}$$
(3.26)
and use the fundamental branch as in the appendix B. If $`f(x)`$ cross the cut of the $`\mathrm{log}`$ on the negative real axis, the function $`\mathrm{log}_{FD}f(x)`$ has a jump by $`\pm 2\pi i`$. The exact real point $`y_{}`$ where this happens (suppose for simplicity only one such point; the extension is trivial) is given (see the figure 3.3)
by the condition $`f(y_{})]\mathrm{},0[`$, that is:
$$1+(1)^\delta e^{i\mathrm{}eZ_N(y_{}+i\eta )}=0\mathrm{and}\mathrm{}mZ_N(y_{}+i\eta )<0.$$
(3.27)
Only the expression
$$\mathrm{log}_{FD}f(x)2\pi i\theta (xy_{})ϵ$$
gives rise to a continuous function. $`ϵ=\pm 1`$ is a sign that is positive if the cross is in the clockwise direction (from below to up) as in figure 3.4 and negative in the opposite case. An important observation can be made by the following Taylor expansion:
$$\mathrm{}eZ_N(y_{}+i\eta )=Z_N(y_{})\frac{\eta ^2}{2}Z_N^{\prime \prime }(y_{})+\mathrm{}$$
(3.28)
$$\mathrm{}mZ_N(y_{}+i\eta )=\eta \left(Z_N^{}(y_{})\frac{\eta ^2}{6}Z_N^{\prime \prime \prime }(y_{})+\mathrm{}\right).$$
(3.29)
It is clear that if and only if $`\eta \mathrm{\hspace{0.17em}0}`$ the condition (3.27) becomes the definition of a special solution, i.e. $`y_{}`$ is a special object only in the limit of $`\eta `$ going to zero because in this case the second terms in (3.28, 3.29) are negligible. Then the condition (3.27) becomes $`Z_N^{}(y_{})<0`$ and $`Z_N=2\pi I`$. The sign must be chosen positive: $`ϵ=1.`$<sup>4</sup><sup>4</sup>4 In what follows, always this case will be assumed. For the general case $`\eta 0`$, $`y_{}`$ is a special object $`y`$ shifted by a little bit.
A completely analogous analysis can be performed for the lower integral, obtaining a $`y_{}`$. From (3.28), because of the real analyticity of $`Z_N`$, the following equation holds: $`y_{}=y_{}\widehat{y}.`$ The result is:
$$\frac{(1)^\delta e^{\pm iZ_N(x\pm i\eta )}(\pm i)Z_N^{}(x\pm i\eta )}{1+(1)^\delta e^{\pm iZ_N(x\pm i\eta )}}=\frac{d}{dx}\mathrm{log}_{FD}\left(1+(1)^\delta e^{\pm iZ_N(x\pm i\eta )}\right)2\pi i\delta (x\widehat{y}).$$
(3.30)
Using this expression in (3.25), a new source term appears every times there is a such point $`\widehat{y}`$:
$$\begin{array}{c}Z_N^{}(\vartheta )=N\left[\frac{1}{\mathrm{cosh}(\vartheta +\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(\vartheta \mathrm{\Theta })}\right]+\underset{k=1}{\overset{N_H}{}}2\pi G(\vartheta h_k)+\\ \underset{k=1}{\overset{M_C}{}}2\pi G(\vartheta c_k)\underset{k=1}{\overset{M_W}{}}2\pi G_{II}(\vartheta w_k)2\pi \underset{k=1}{\overset{N_S}{}}\left(G(\vartheta \widehat{y}_k+i\eta )+G(\vartheta \widehat{y}_ki\eta )\right)\\ +\frac{1}{i}𝑑xG(\vartheta xi\eta )\frac{d}{dx}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(x+i\eta )}\right]+\\ \frac{1}{i}𝑑xG(\vartheta x+i\eta )\frac{d}{dx}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(xi\eta )}\right].\end{array}$$
(3.31)
The sum on $`\widehat{y}_k`$ has been indicated as sum on the specials, even if they are not exactly in the position of specials, because in the usual computations made on the NLIE $`\eta `$ takes a small value. The condition to be used for the correct position is (3.27).
In the following, we omit the label FD, intending always the $`\mathrm{log}`$ in the fundamental branch. It is very important to note that special objects are not completely independent degrees of freedom, unlike the holes and complex solutions, that are fixed “a priori”. Special objects appear when the derivative of $`Z_N`$ on a root or on a hole becomes negative. A special object then apparently appears two times in (3.31), both as a real root/hole and as a special root/hole.
### 3.5 Non linear integral equation (II)
From a Bethe Ansatz point of view, only the function $`Z_N`$ (and not the derivative) has a physical meaning, because it is required to obtain the position of various types of solution, given the quantum numbers (3.2). To obtain $`Z_N`$, an integration in $`\vartheta `$ must be performed in (3.31). All the source terms must be integrated. For the case of holes, specials and closes this is given by the function:
$$\chi (\theta )=_0^\theta 𝑑x2\pi G(x)$$
(3.32)
It is and odd and real analytic function. It asymptotic value is known by direct integration:
$$\chi _{\mathrm{}}=\chi (+\mathrm{})=\frac{\pi }{2}\frac{p1}{p}$$
(3.33)
Because $`G`$ has always the same sign for every $`x`$, the $`\chi (\theta )`$ is a monotonous function.
The integration of the source for wide roots is simply the integral of (3.23). But this integration produces logarithms, then there is a problem of definition of the appropriate branch. The $`Z_N`$ function is continuous in the analyticity strip, then the most correct choice for the wide sources satisfies this continuity. The forms that are proposed in the following do so (they have no jumps) for $`\theta `$ with imaginary part $`\pi \mathrm{max}(1,p)>|\mathrm{}m\theta |>\pi \mathrm{min}(1,p)`$ (the $`\mathrm{log}`$ is in the fundamental determination):
$$\chi _{II}(\theta )=\{\begin{array}{c}i\text{sign}(\mathrm{}m\theta )\left(\mathrm{log}\mathrm{sinh}\left(\frac{\theta }{p}\right)\mathrm{log}\mathrm{sinh}\left(\frac{\theta i\pi \text{sign}\mathrm{}m\theta }{p}\right)\right)\mathrm{if}p>1\\ \\ i\text{sign}(\mathrm{}m\theta )\left(\mathrm{log}\left(\mathrm{tanh}\frac{\theta }{2}\right)+\mathrm{log}\left(\mathrm{tanh}\frac{\theta i\pi p\text{sign}\mathrm{}m\theta }{2}\right)\right)\mathrm{if}p<1\end{array}$$
(3.34)
The notation $`\chi _{II}`$ is used to indicate the integral of the $`G_{II}`$. As in the previous case of $`G`$, this is the same as taking the second determination of $`\chi `$ in (3.32). The two forms corresponding to the opposite signs of the variable are related by the oddity:
$$\chi _{II}(\theta )=\chi _{II}(\theta )\text{if}|\mathrm{}m\theta |>\pi \mathrm{min}(1,p).$$
This function can be continued to smallest values of the imaginary part, but obviously it has a line of discontinuity. The asymptotic limits are (now $`\vartheta `$ is real):
$$\underset{\vartheta \pm \mathrm{}}{lim}\chi _{II}(\vartheta w)=\{\begin{array}{c}\pm \left(2\chi _{\mathrm{}}\pi \right)\mathrm{if}p>1\\ \\ \pm \pi \mathrm{if}p<1\end{array}$$
(3.35)
The final form obtained in this way is the so called fundamental non linear integral equation, also known as *Destri-de Vega equation* (see ) for the counting function:
$$\begin{array}{c}Z_N(\vartheta )=2N\mathrm{arctan}\frac{\mathrm{sinh}\vartheta }{\mathrm{cosh}\mathrm{\Theta }}+g(\vartheta |\vartheta _j)+\\ \\ +\frac{dx}{i}G(\vartheta xi\eta )\mathrm{log}\left(1+(1)^\delta e^{iZ_N(x+i\eta )}\right)+\\ \\ \frac{dx}{i}G(\vartheta x+i\eta )\mathrm{log}\left(1+(1)^\delta e^{iZ_N(xi\eta )}\right)+\alpha \end{array}$$
(3.36)
where
$$\begin{array}{c}g(\vartheta |\vartheta _k)=\underset{k=1}{\overset{N_H}{}}\chi (\vartheta h_k)\underset{k=1}{\overset{N_S}{}}\left(\chi (\vartheta \widehat{y}_k+i\eta )+\chi (\vartheta \widehat{y}_ki\eta )\right)+\\ \\ \underset{k=1}{\overset{M_C}{}}\chi (\vartheta c_k)\underset{k=1}{\overset{M_W}{}}\chi _{II}(\vartheta w_k)\end{array}$$
(3.37)
Both the previous equations hold for $`\vartheta `$ in the fundamental strip. Out of that, the analytic continuation of $`Z_N`$ is required, because the first singularity of $`G`$ is crossed. To obtain an equation expressing the analytic continuation of $`Z_N`$ beyond this singularity, the second determination of the various terms appearing in the right hand side of (3.36) must be taken, as defined in (3.24). Notice that it is not the second determination of $`Z_N`$ but its analytic continuation:
$$\begin{array}{c}\text{for }|\mathrm{}m\vartheta |>\mathrm{min}(1,p):Z_N(\vartheta )=2N\left[\mathrm{arctan}\frac{\mathrm{sinh}\vartheta }{\mathrm{cosh}\mathrm{\Theta }}\right]_{II}+\\ \\ +g(\vartheta |\vartheta _k)_{II}+\frac{dx}{i}G(\vartheta xi\eta )_{II}\mathrm{log}\left(1+(1)^\delta e^{iZ_N(x+i\eta )}\right)+\\ \\ \frac{dx}{i}G(\vartheta x+i\eta )_{II}\mathrm{log}\left(1+(1)^\delta e^{iZ_N(xi\eta )}\right)+\alpha _{II}\end{array}$$
(3.38)
This explicit form of the analytic continuation requires to compute $`g(\vartheta |\vartheta _k)_{II}`$. Remembering (3.37), it is clear that for holes, specials and close roots the second determination of $`\chi (\vartheta \vartheta _j)`$ (see (3.34)) appears. For the wide roots the second determination of $`\chi _{II}(\vartheta w_j)`$ appears. It must be determined considering that
$$\chi _W(\vartheta ,w_j)\chi _{II}(\vartheta w_j)$$
is a function of $`\vartheta `$:
$$\chi _{II}(\vartheta w_j)_{II}\left(\chi _W(\vartheta ,w_j)\right)_{II}.$$
(3.39)
The explicit forms are not given there but they can be simply computed from the definition of second determination.
The integral of the convolution term is possible because $`G`$ vanishes exponentially at infinity. $`\alpha `$ is the integration constant and must be determined.
There is a useful way to write the log term in (3.36), valid for real $`x`$ and if $`\eta \mathrm{\hspace{0.17em}0}`$:
$$\begin{array}{c}𝒬_N(x)lim_{\eta \mathrm{\hspace{0.17em}0}}2\mathrm{}m\mathrm{log}\left(1+(1)^\delta e^{iZ_N(x+i\eta )}\right)=\\ \\ =lim_{\eta \mathrm{\hspace{0.17em}0}}\frac{1}{i}\mathrm{log}\frac{1+(1)^\delta e^{iZ_N(x+i\eta )}}{1+(1)^\delta e^{iZ_N(xi\eta )}}=\left(Z_N(x)+\pi \delta \right)\text{ mod }2\pi \end{array}$$
(3.40)
The expression $`\left(A\right)\text{ mod }2\pi `$ means exactly $`\pi <A\pi .`$ The second line in (3.40) holds only if $`Z_N^{}(x)>0`$, because the branch cut of the $`\mathrm{log}`$ is not crossed and $`\left|\mathrm{}m\mathrm{log}f(x)\right|<\pi `$ (the notation (3.26) is used). In the opposite case the branch cut is crossed and the previous expression can take values larger than $`\pi `$: $`\left|\mathrm{}m\mathrm{log}f(x)\right|<2\pi .`$ In the limit case $`Z_N^{}(x)=0`$ (for example this happens when $`x\mathrm{}`$) the values permitted are $`\pi \mathrm{}m\mathrm{log}f(x)\pi .`$
At this point it is possible to explicitly calculate, from (3.36), the limit $`Z_N(+\mathrm{})`$, using (3.36, 3.7, 3.33, 3.35). The limit gives:
$$\begin{array}{c}Z_N(+\mathrm{})=N\pi +\chi _{\mathrm{}}\left(N_H2N_SM_C2\theta (p1)M_W\right)+\\ \\ +\pi \text{sign}(p1)M_W+𝑑xG(x)𝒬_N(+\mathrm{})+\alpha .\end{array}$$
(3.41)
The limit $`𝒬_N(+\mathrm{})`$ can be calculated by using in (3.40) the limit of $`Z_N`$ (3.5). By a comparison of the (3.41) with the asymptotic values of $`Z_N`$ calculated in (3.5) the value of $`\alpha `$ can be obtained:
$$\alpha =\omega \frac{p+1}{p}+\chi _{\mathrm{}}\left(\frac{1}{2}+\frac{S}{p+1}+\frac{\omega }{\pi }\frac{1}{2}+\frac{S}{p+1}\frac{\omega }{\pi }\right)$$
(3.42)
In this expression the following term has been omitted: $`\pi \text{sign}(p1)M_{SC}.`$ It is a multiple of $`\pi `$ and can be taken into account simply changing the value of $`\delta `$ from $`0`$ to $`1`$ (or viceversa) if $`M_{SC}`$ is odd, and shifting of an integer quantity the quantum numbers
$$\delta =(M+M_{SC})_{mod\mathrm{\hspace{0.17em}2}}=(NS+M_{SC})_{mod\mathrm{\hspace{0.17em}2}}\{0,1\}.$$
(3.43)
This sort of manipulation on the quantum numbers has been explained at the end of the section 3.1.
Observe that the result is exactly $`\alpha =0`$ if there is no twist. There is another important observation. In (2.37) the twist term $`\omega `$ is invariant for the shift
$$\omega \omega +\pi .$$
The same sort of invariance is required in $`\alpha `$ for the NLIE (which is equivalent to Bethe equations). It is simple to verify that the expression for $`\alpha `$ (3.42) displays the following symmetry
$$\alpha \alpha +2\pi \text{when}\omega \omega +\pi .$$
Note that shifting $`\alpha `$ by $`2\pi `$ is an invariance of the NLIE (3.36), if an appropriate redefinition of the Bethe quantum numbers is made: $`I_jI_j+1`$. This shift does not affects physical quantities, that depend only on the variables $`\vartheta _j`$.
With the given value of the integration constant the equation (3.36) is complete. The quantization condition (3.2) can now be written as:
$$Z_N(\vartheta _j)=2\pi I_j,I_j+\frac{1+\delta }{2}$$
(3.44)
for the various solutions. For the wide roots remember the previous warning about the analytic continuation of the $`Z_N`$ function.
The previous equation is completely equivalent to Bethe equations. No new physics has been introduced, until this point. Simply an equation that generates the counting function $`Z_N`$ has been obtained (the NLIE).
### 3.6 Energy and momentum
The expression (2.41) can be written using the function $`\varphi `$ as:
$$\begin{array}{c}E=\frac{1}{a}\underset{j=1}{\overset{M}{}}\left(\varphi (\mathrm{\Theta }\vartheta _j,1/2)+\varphi (\mathrm{\Theta }+\vartheta _j,1/2)2\pi \right)\\ \\ P=\frac{1}{a}\underset{j=1}{\overset{M}{}}\left(\varphi (\mathrm{\Theta }\vartheta _j,1/2)\varphi (\mathrm{\Theta }+\vartheta _j,1/2)\right)+2\omega \end{array}$$
(3.45)
The choice of the logarithmic branch in the energy ensures that the contribution of each real root is negative definite. This is consistent with the known ground state structure (as in section 2.5), that is given by the maximal number of real roots. It will be clear that excitations give only positive contributions.
It is possible to relate the previous expressions to the counting function. To do that, consider first the following quantity:
$$W(\theta )=\underset{j=1}{\overset{M}{}}\varphi ^{}(\theta \vartheta _j,1/2)\theta .$$
(3.46)
It can be integrated to obtain the pieces appearing in (3.45). Clearly the following expressions hold:
$$\begin{array}{c}\underset{j=1}{\overset{M}{}}\varphi (\mathrm{\Theta }\vartheta _j,1/2)=_{asymp}^\mathrm{\Theta }𝑑xW(x)\\ \underset{j=1}{\overset{M}{}}\varphi (\mathrm{\Theta }+\vartheta _j,1/2)=\underset{j=1}{\overset{M}{}}\varphi (\mathrm{\Theta }\vartheta _j,1/2)=_{asymp}^\mathrm{\Theta }𝑑xW(x).\end{array}$$
(3.47)
The integration requires to fix a constant. This can be simply done by computing for $`\mathrm{\Theta }+\mathrm{}`$ the left hand side in (3.47) and imposing the equality with the primitive of $`W`$ at the same limit. This has been indicated by the symbol
$$_{asymp}^{\pm \mathrm{\Theta }}$$
The function $`W`$ admits an expression in terms of the counting function. The sum over roots in (3.46) can be expressed as in (3.8) and using the same notations:
$$W(\theta )=\underset{j=1}{\overset{M_R+N_H}{}}\varphi ^{}(\theta x_j,1/2)\underset{j=1}{\overset{N_H}{}}\varphi ^{}(\theta h_j,1/2)+\underset{j=1}{\overset{M_C+M_W}{}}\varphi ^{}(\theta \xi _j,1/2).$$
Using now the same trick used for $`Z_N^{}`$, as in (3.13, 3.14, 3.15), and the expression (3.30) generating the “special” contribution, the following expression yields:
$$\begin{array}{c}W(\vartheta )=\frac{dx}{2\pi i}\varphi ^{}(\vartheta x+i\eta ,1/2)\frac{d}{dx}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(xi\eta )}\right]+\\ \frac{dx}{2\pi i}\varphi ^{}(\vartheta xi\eta ,1/2)\frac{d}{dx}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(x+i\eta )}\right]+\\ \underset{k=1}{\overset{N_H}{}}\varphi ^{}(\vartheta h_k,1/2)+\underset{k=1}{\overset{M_C}{}}\varphi ^{}(\vartheta c_k,1/2)+\underset{k=1}{\overset{M_W}{}}\varphi ^{}(\vartheta w_k,1/2)+\\ +\underset{k=1}{\overset{N_S}{}}\left(\varphi ^{}(\vartheta \widehat{y}_k+i\eta ,1/2)+\varphi ^{}(\vartheta \widehat{y}_ki\eta ,1/2)\right)+\frac{dx}{2\pi }\varphi ^{}(\vartheta x,1/2)Z_N^{}(x).\end{array}$$
(3.48)
Now the equation (3.31) can be put in the last integral in the previous expression. The terms “similar” can be collected (i.e. holes with holes, close roots with close roots, and so on) to obtain the following form:
$$\begin{array}{c}W(\vartheta )=\underset{k=1}{\overset{N_H}{}}\left[\varphi ^{}(\vartheta h_k,1/2)𝑑x\varphi ^{}(\vartheta x,1/2)G(xh_k)\right]+\\ +\underset{k=1}{\overset{M_C}{}}\left[\varphi ^{}(\vartheta c_k,1/2)𝑑x\varphi ^{}(\vartheta x,1/2)G(xc_k)\right]+\\ +\underset{k=1}{\overset{M_W}{}}\left[\varphi ^{}(\vartheta w_k,1/2)𝑑x\varphi ^{}(\vartheta x,1/2)G_{II}(xw_k)\right]+\\ +\underset{k=1}{\overset{N_S}{}}[\varphi ^{}(\vartheta \widehat{y}_k+i\eta ,1/2)+\varphi ^{}(\vartheta \widehat{y}_ki\eta ,1/2)+\\ dx\varphi ^{}(\vartheta x,1/2)(G(x\widehat{y}_k+i\eta )+G(x\widehat{y}_ki\eta ))]+\\ +\frac{dx}{2}\varphi ^{}(\vartheta x,1/2)N\left[\frac{1}{\mathrm{cosh}(\vartheta +\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(\vartheta \mathrm{\Theta })}\right]+\\ 2\mathrm{}m\frac{dx}{2\pi }\varphi ^{}(\vartheta xi\eta ,1/2)𝑑y\left[\delta (xy)G(xy)\right]\frac{d}{dy}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(y+i\eta )}\right]=\\ =W_H+W_C+W_W+W_S+W_{bulk}+W_I\end{array}$$
where the notation introduced on the last line is obvious (the meaning of bulk will be clear in the following; label I is for interaction term or, that is the same, for integral).
As in section 3.4, there are different terms to analyze.
First, rearranging a little bit the terms and introducing a delta-function, it is possible to see that in the case of holes, closes, specials and in the integral term there is a contribution by:
$$\delta (xy)G(xy)=\mathrm{\Delta }(xy)$$
(this equality is a simple consequence of the definition of $`\mathrm{\Delta }`$). The corresponding terms contain the integral:
$$𝑑x\varphi ^{}(x\vartheta _j,1/2)\mathrm{\Delta }(\vartheta x)=\frac{1}{\mathrm{cosh}(\vartheta \vartheta _j)}$$
that is well known because it has been used in (3.19). This holds in the case of holes and closes and $`\vartheta _j`$ is the corresponding rapidity. For the integral term and the specials the contribution is exactly:
$$\frac{1}{\mathrm{cosh}(\vartheta yi\eta )}.$$
where $`y`$ is the integration variable or the special rapidity.
The wide computation is more involved. The expression to be computed is
$$W_W(\vartheta )=\underset{k=1}{\overset{M_W}{}}\left[\varphi ^{}(\vartheta w_k,1/2)𝑑x\varphi ^{}(\vartheta x,1/2)G_{II}(xw_k)\right]$$
and the most convenient way to do that is to use the Fourier transformation of all the terms. For $`\varphi ^{}(\vartheta w_k,1/2)`$ and $`G_{II}`$ it has been used in the computation for (3.23). For $`\varphi ^{}`$ on the real axis it is in (B.6). The result is
$$W_W(\vartheta )=\{\begin{array}{c}0\\ \left[\frac{1}{\mathrm{cosh}(\vartheta )}\right]_{II}+O(\vartheta )\end{array}\begin{array}{c}\text{ for }p>1\\ \text{ for }p<1\end{array}$$
In the case of $`p>1`$ it simply cancels out. In the other case, there is no such cancellation, for generic wide positions. The explicit form of the contribution $`O(\vartheta )`$ is quite long and shall not be written out at this point. The important fact is that it is a fast vanishing contribution for large $`\vartheta `$. Observe that the second determination (3.24) of a whatever trigonometric hyperbolic function, for $`p>1`$, is exactly zero, then the notation $`[1/\mathrm{cosh}(\vartheta )]_{II}`$ will be used also in this case.
The final form is:
$$\begin{array}{c}W(\vartheta )=\underset{k=1}{\overset{N_H}{}}\frac{1}{\mathrm{cosh}(\vartheta h_k)}+\underset{k=1}{\overset{N_S}{}}\left(\frac{1}{\mathrm{cosh}(\vartheta \widehat{y}_ki\eta )}+\frac{1}{\mathrm{cosh}(\vartheta \widehat{y}_k+i\eta )}\right)+\\ +\underset{k=1}{\overset{M_C}{}}\frac{1}{\mathrm{cosh}(\vartheta c_k)}+\underset{k=1}{\overset{M_W}{}}\left[\frac{1}{\mathrm{cosh}(\vartheta w_k)}\right]_{II}+O(\vartheta )+\\ +\frac{dx}{2\pi }\varphi ^{}(\vartheta x,1/2)N\left[\frac{1}{\mathrm{cosh}(x+\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(x\mathrm{\Theta })}\right]+\\ 2\mathrm{}m\frac{dx}{2\pi }\frac{1}{\mathrm{cosh}(\vartheta xi\eta )}\frac{d}{dx}\mathrm{log}_{FD}\left[1+(1)^\delta e^{iZ_N(x+i\eta )}\right]\end{array}$$
Integrating this function, as in (3.47), the energy and momentum of this lattice system can be obtained. In all the terms appears $`1/\mathrm{cosh}x`$; its primitive is
$$\frac{dx}{\mathrm{cosh}x}=\mathrm{arctan}\mathrm{sinh}x=2\mathrm{arctan}\mathrm{tanh}(x/2).$$
(3.49)
The last form is the most convenient, to calculate the continuum limit. Then:
$$\begin{array}{c}a\frac{E\pm P}{2}=\underset{k=1}{\overset{N_H}{}}2\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }h_k)/2)+\underset{k=1}{\overset{M_C}{}}2\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }c_k)/2)+\\ +\underset{k=1}{\overset{N_S}{}}2\left(\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }\widehat{y}_ki\eta )/2)+\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }\widehat{y}_k+i\eta )/2)\right)+\\ +\underset{k=1}{\overset{M_W}{}}2\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }w_k)/2)_{II}+O(\mathrm{\Theta })\frac{dx}{2\pi }\frac{1}{\mathrm{cosh}(\pm \mathrm{\Theta }x)}𝒬_N(x)+S\pi +\\ +\frac{dx}{2\pi }\varphi (\mathrm{\Theta }x,1/2)N\left[\frac{1}{\mathrm{cosh}(x+\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(x\mathrm{\Theta })}\right]\pm \omega N\pi \end{array}$$
(3.50)
The term $`O(\mathrm{\Theta })`$ vanishes for large $`\mathrm{\Theta }`$.
The function (3.40) has been used, by putting $`\eta \mathrm{\hspace{0.17em}0}`$, and observing that $`x`$ is real.
### 3.7 Continuum limit
At the conclusion of the previous long computations there is an important observation: in the final expressions for energy and momentum there is no explicit contribution from the real roots. This also happens in the NLIE (3.36). On the lattice it is not a relevant observation, but in this section it shall be shown that it can justify a continuum limit procedure, $`N\mathrm{}`$ and $`a\mathrm{\hspace{0.17em}0}`$, and a particle interpretation. In the lattice energy expression there is a bulk term (see later) i.e. a term increasing with $`N`$ and completely independent from the solutions $`\vartheta _j`$ of Bethe equations. This bulk term can be subtracted and the remaining terms (due only to holes and complex roots) can be chosen in such a way that they describe particle excitations.
To do that, in the limit procedure, they must be chosen in finite number (order $`O(1)`$) and also its “rapidity” $`\vartheta _j`$ must be finite. From (3.2) this means that only real roots can appear in the asymptotic tails of $`Z_N`$. All the other lattice states must be discarded (as usual, lattice theory contains more states than the continuum theory).
A state where holes and complex roots are not considered, is the hamiltonian vacuum. If they are considered, they give positive contribution to the energy, behaving as particle excitations on a vacuum state. Real roots are completely disappeared, in this limit. They can be interpreted (in the limit procedure) as a sort of Dirac sea that is the hamiltonian vacuum on with holes and complex roots built particle excitations.
This is the structure required to have a consistent quantum field theory.
All this interpretation is possible only because of the antiferromagnetic vacuum choice, made in section (2.5).
What will be proved now is that there is a consistent way to do continuum limit in the NLIE and in the expression of the energy.
In it was shown that the correct way to do the continuum limit for this model on the light-cone is to send $`N,\mathrm{\Theta }\mathrm{}`$ connected in the following way:
$$\mathrm{\Theta }\mathrm{log}\frac{4N}{L}.$$
(3.51)
$`L=Na`$ is the spatial dimension of the lattice (as in section 2.1) and stay fixed in the limit. $``$ is the renormalized physical mass. As in section 2.5, the light-cone lattice is periodic in space direction. Then its continuum limit yields a compactified space of length $`L`$ (space-time is a cylinder of circumference $`L`$). The lattice 6 vertex model becomes a field theory defined on this cylinder. The interpretation of what field theory is defined by this continuum limit, at the various values of the twist, is the job of the next chapter.
To do the limit, consider first the counting equation. It takes the following form, on the continuum
$$N_H2N_S=2S+M_C+2\theta (p1)M_W$$
(3.52)
because the structures in the tails are extremely simplified. The other terms appearing in the lattice counting equation (3.7) in fact came from the configurations where in the tails there are special root/holes or self-conjugate of the so called first class (see the original paper ). Such configurations completely disappear in the continuum limit.
The principal interest is to obtain the limit of the energy and momentum eigenvalues (3.50). Clearly, to fix the positions of the roots $`\vartheta _j`$ the limit of (3.44) is required. This can be made using a *continuum counting function* that is defined by:
$$Z(\vartheta )=\underset{N\mathrm{}}{lim}Z_N(\vartheta ).$$
(3.53)
This limit computation is made on the sequence of counting functions that are implicitly defined in (3.36) at the various values of $`N`$; the (3.51) must be also taken onto account. The interesting fact is that this limit can be done explicitly in all the terms appearing in (3.36), i.e. a *continuum NLIE* can be obtained. To show that, the various terms are analyzed.
The first term and the integral in (3.36) are trivial matter to compute. The source terms for holes, complex solutions and specials don’t have any explicit dependence from $`N`$. Only the positions that appear in this terms must be determined, obviously, by the continuum NLIE. Then the term $`g(\vartheta |\vartheta _j)`$ is unchanged. The result is:
$$\begin{array}{c}Z(\vartheta )=L\mathrm{sinh}\vartheta +g(\vartheta |\vartheta _j)+\\ +\frac{dx}{i}G(\vartheta xi\eta )\mathrm{log}\left(1+(1)^\delta e^{iZ(x+i\eta )}\right)+\\ \frac{dx}{i}G(\vartheta x+i\eta )\mathrm{log}\left(1+(1)^\delta e^{iZ(xi\eta )}\right)+\alpha \end{array}$$
(3.54)
The quantization conditions (3.2, 3.44) can now be written as:
$$Z(\vartheta _j)=2\pi I_j,I_j+\frac{1+\delta }{2}$$
(3.55)
for the various solutions. This terminates the limit procedure on NLIE. Observe that on the continuum only the implicit “definition” of $`Z`$ by (3.54) is available, instead of and explicit expressions as in the case (3.1).
Consider now in the energy expression (3.50) the terms with the form $`2/a\mathrm{arctan}\mathrm{tanh}((\mathrm{\Theta }\theta )/2)`$. Using the following hyperbolic trigonometric identity
$$\mathrm{arctan}\mathrm{tanh}(x)=\frac{\pi }{4}\mathrm{arctan}e^x$$
the following asymptotic behaviour can be obtained:
$$\frac{N}{L}\mathrm{\hspace{0.17em}2}\mathrm{arctan}\mathrm{tanh}(\mathrm{\Theta }\theta )\frac{N\pi }{L2}\frac{1}{2}e^\theta .$$
Then, collecting in the expression for $`(E\pm P)/2`$ the contribution of this type, yields (remember that there is a $`S\pi `$ in (3.50)):
$$\frac{\pi }{2}(2SN_H+2N_S+M_C+2\theta (p1)M_W)=0$$
(3.56)
(the (3.52) has been used).
One of the two integral contributions is
$$\begin{array}{c}\underset{N\mathrm{}}{lim}\frac{N}{L}\frac{dx}{2\pi }\frac{1}{\mathrm{cosh}(\pm \mathrm{\Theta }x)}𝒬_N(x)=\\ =\frac{}{2}\frac{dx}{2\pi }e^{\pm x}𝒬(x).\end{array}$$
The final result has been obtained by exchanging the limit and the integral.
The other one is the last line in (3.50). It can be handled by shifting the integration variable and observing that one term in the square brackets is odd and gives a vanishing integral:
$$\begin{array}{c}E_{N,bulk}^\pm =\frac{N^2}{L}\left[\frac{dx}{2\pi }\varphi (\mathrm{\Theta }x,1/2)\left[\frac{1}{\mathrm{cosh}(x+\mathrm{\Theta })}+\frac{1}{\mathrm{cosh}(x\mathrm{\Theta })}\right]\pi \right]\pm \omega =\\ =\frac{N^2}{L}\left[\frac{dx}{2\pi }\left[\frac{\varphi (x,1/2)}{\mathrm{cosh}(2\mathrm{\Theta }x)}\right]\pi \right]\pm \omega .\end{array}$$
This energy contribution is “source independent”. It has been labeled by “bulk” because it will be shown that it diverges as $`N`$. It is interesting to compute the contribution to $`E`$ and $`P`$, instead of the diverging contribution to $`(E\pm P)/2`$:
$$\begin{array}{c}E_{Nbulk}=E_{N,bulk}^++E_{N,bulk}^{}=\frac{N^2}{L}\left[\frac{dx}{\pi }\left[\frac{\varphi (x,1/2)}{\mathrm{cosh}(2\mathrm{\Theta }x)}\right]2\pi \right]=\frac{N^2}{L}(\pi +\gamma )\\ \\ P_{Nbulk}=E_{Nbulk}^+E_{Nbulk}^{}=2\omega \end{array}$$
The first integral can be computed observing that $`\varphi `$ is a constant for very large values of $`x`$, and can be put out of the integral sign. The second one contains the difference of two equal contributions and only $`\omega `$ is not deleted.
Collecting all the terms follows:
$$\begin{array}{c}\frac{E\pm P}{2}=\frac{}{2}(\underset{k=1}{\overset{N_H}{}}e^{\pm h_k}\underset{k=1}{\overset{N_S}{}}(e^{\pm \widehat{y}_k+i\eta }+e^{\pm \widehat{y}_ki\eta })+\\ \underset{k=1}{\overset{M_C}{}}e^{\pm c_k}\underset{k=1}{\overset{M_W}{}}e_{II}^{\pm w_k}\frac{dx}{2\pi }e^x𝒬_N(x))+E^\pm _{Nbulk}\end{array}$$
(3.57)
or, subtracting all the diverging contributions,
$$\begin{array}{c}E=(\underset{k=1}{\overset{N_H}{}}\mathrm{cosh}h_k\underset{k=1}{\overset{N_S}{}}(\mathrm{cosh}(\widehat{y}_k+i\eta )+\mathrm{cosh}(\widehat{y}_ki\eta ))+\\ \underset{k=1}{\overset{M_C}{}}\mathrm{cosh}c_k\underset{k=1}{\overset{M_W}{}}\mathrm{cosh}_{II}w_k\frac{dx}{2\pi }\mathrm{sinh}x𝒬(x))\end{array}$$
(3.58)
and
$$\begin{array}{c}P=(\underset{k=1}{\overset{N_H}{}}\mathrm{sinh}h_k\underset{k=1}{\overset{N_S}{}}(\mathrm{sinh}(\widehat{y}_k+i\eta )+\mathrm{sinh}(\widehat{y}_ki\eta ))+\\ \underset{k=1}{\overset{M_C}{}}\mathrm{sinh}c_k\underset{k=1}{\overset{M_W}{}}\mathrm{sinh}_{II}w_k+\frac{dx}{2\pi }\mathrm{cosh}x𝒬(x))\end{array}$$
(3.59)
Observe that for $`p>1`$ the second determination of hyperbolic functions vanishes. This means that wide roots contribute to the energy only in an implicit way, because they contribute to the position of other objects, by (3.54).
### 3.8 Physical interpretation
The limit procedure described in the previous section is mathematically consistent, but the question is if from the physical point of view it describes a consistent quantum theory and allows for a meaningful physical interpretation.
Before to propose it, an important remark must be made about the allowed values for $`S`$. It is clear from (2.40) that on the lattice only integer and nonnegative values can be taken into account for $`S`$. But on the continuum the definition of $`S`$ is no more related to the Bethe state (that is undefined), instead it is given implicitly by (3.52). Then, *“a priori”*, there are no arguments that fixes its values to be integers. As our group showed in , the half-integer choice for $`S`$ is necessary (and gives completely consistent results) to describe odd numbers of particles.
At this point the following physical interpretation can be proposed. It will be refined to describe the correspondence with particles. Also it will be supported by many arguments that will be clarified in the next chapters.
Physical interpretation:
* the physical vacuum (hamiltonian ground state) corresponds to absence of sources (i.e. holes, complex); all the sources are excitations on this vacuum
* for $`\omega =0`$ and at the various values of $`S`$ this theory describes the sine-Gordon/massive Thirring model on a finite space of size $`L`$; $`2S`$ is the topological charge and can take nonnegative integer values.
* for the values
$$\omega =\frac{k\pi }{s},k=1,\mathrm{},q^{}1$$
(3.60)
it describes the quantum reduction of sine-Gordon model, i.e. the massive integrable theory obtained perturbing the minimal model $`Vir(r,s)`$ with the operator $`\mathrm{\Phi }_{(1,3)}`$.
Observe that it has been assumed that only nonnegative values of $`S`$ are required to describe the whole Hilbert space of the theory. Indeed the theory is assumed charge-conjugation invariant then negative topological charge states have the same energy and momentum as their charge conjugate states. The assumption that all the states can be described by the NLIE is absolutely not trivial, or better still until this moment it is not available a general proof of this fact, but only a number of specific cases supports this conjecture.
All this things are the argument of the next chapters.
## Chapter 4 ANALYSIS OF THE CONTINUUM THEORY
### 4.1 Principal questions
In the next sections, a carefull analysis of the vacuum and of some excited states from NLIE will be performed at the various scales (various $`L`$), to completely understand the physical interpretation suggested at the end of the section 3.7. The principal question is what theory is described by equations (3.54, 3.58, 3.59).
As scale parameter, as appears in (3.54), can be chosen equivalently the size $`L`$ or the *adimensional size* $`l=L`$, where $``$ plays the role of a mass scale. The limit of very large $`l`$ can be interpreted both as large size (that reproduces infinite Minkowski space-time) or large mass scale, that is an *infrared point* in a renormalization flow (IR). Finite size effects do not appear, in that limit.
At the opposite limit of small $`l`$, important finite size effects are expected; a small value of $`l`$ can be obtained with a small mass $``$, then this case is an *ultraviolet point* in a renormalization flow (UV). The complete control on the scaling functions can be obtained analyzing the whole range of values $`l>0`$.
### 4.2 Connection with sine-Gordon/massive Thirring
The section 3.7 ends with a conjecture about the physical meaning of the model described by the continuum NLIE and the corresponding energy and momentum expressions. There are many arguments to support this interpretation. Some of them came from the well known properties of the 6-vertex Bethe Ansatz, other from the analysis of the NLIE itself.
Consider the $`R`$ matrix in (2.33, 2.35). It has the same entries as the $`R_{sG}`$ matrix that appears by putting sine-Gordon on a lattice, and obtaining the corresponding Bethe Ansatz equations (lattice Thirring B.A. gives exactly the same). The equality of $`R`$ matrix and of Bethe Ansatz equations means that there is a common integrable structure in the two systems.
Moreover, the lattice sine-Gordon B.A. can be obtained as the continuum limit $`N\mathrm{}`$ (and with $`\mathrm{\Theta }=\mathrm{log}{\displaystyle \frac{4N}{L}}`$) of the inhomogeneous 6-vertex Bethe equations (2.37). The Bethe Ansatz techniques, then, strongly suggest the previous interpretation.
The 6-vertex model in its thermodynamics limit is critical, as shown in , then there is a conformal field theory describing this critical point. From the eigenvalues of the transfer matrix (obtained by B.A.) it is possible to extract the conformal properties, as shown in . As shown in , for 6-vertex model they reproduce exactly the equation (1.4) for $`\omega =0`$ (that is the UV structure of sine-Gordon and massive Thirring) and the minimal models conformal weights for $`\omega `$ chosen as in (3.60).
It has been shown, in , that a fermion satisfying massive Thirring equations of motion can be constructed, as described in section (2.4). It describes the scaling behaviour of the inhomogeneous 6-vertex light-cone lattice dynamics. A similar argument can be made on the XXZ chain quantum hamiltonian, whose exponential is the transfer matrix of 6-vertex model. In a (classical) field satisfying sine-Gordon equations is built from this XXZ chain.
There is one more suggestion, that comes from the NLIE itself: the $`\chi `$ function obtained in the derivation of the NLIE (3.32) is well known in literature, because it is exactly the logarithm of the soliton-soliton sine-Gordon $`S`$ matrix:
$$\chi (\vartheta )=i\mathrm{log}S_{++}^{++}(\vartheta )$$
(4.1)
as in .
This arguments are the starting points for the analysis of NLIE itself that will be performed in the next sections. They completely confirm the physical interpretation given in section 3.7.
### 4.3 Some general facts about the IR limit
This IR analysis is important, because it can connect NLIE with the (minkowskian) scattering theory. The limit $`L\mathrm{}`$ is not at all a thermodynamic limit, because the size of the system is scaled, but the number of particles is fixed. Then, the interaction among the physical particles is expected to cease affecting the energy and momentum, because their density is vanishing. Hence, $`E`$ and $`P`$ should approach finite limits equal to a free massive spectrum. In fact, for large $`l`$, the dominant term in (3.54) is the $`l\mathrm{sinh}\vartheta `$. Using it as first iteration in the convolution terms of (3.54, 3.58, 3.59), they become small (of order $`O\left(e^l\right)`$) and can be dropped. The surviving part in NLIE can be seen as a dressed Bethe Ansatz giving constraints on the asymptotic states:
$$Z(\vartheta )=l\mathrm{sinh}\vartheta +g(\vartheta |\vartheta _k)+\alpha \text{for}l\mathrm{}.$$
(4.2)
Out of the fundamental strip where this equation holds, a similar expression can be written, using the second determination. Assume for the moment $`\alpha =0`$. Taking the exponential of the previous equation and imposing the quantization condition yields:
$$e^{il\mathrm{sinh}\vartheta _j}e^{ig(\vartheta _j|\vartheta _k)}=\pm 1$$
that has a structure similar to the usual quantization condition of particles in a box:
$$e^{iLP(\lambda _j)}\underset{kj}{}S(\lambda _j\lambda _k)=\pm 1$$
(4.3)
($`+`$ is referred to periodic and $``$ to antiperiodic boundary conditions). The equation (4.3) is not a Bethe Ansatz equation<sup>1</sup><sup>1</sup>1 It is called Dressed Bethe Ansatz because it contains real particles instead of pseudoparticles appearing in usual B.A. , in general, because the $`\lambda _k`$ that appear in it are real rapidities of particles, instead of the rapidities of the pseudoparticles ($`\vartheta _j`$) required to built Bethe states. But in (4.2) can appear complex “rapidities”. The correct correspondence between the box quantization (4.3) and the infrared NLIE (4.2) will be carefully analyzed later. The important fact, for the moment, is that the comparison between this two equations allows to read out scattering data from NLIE and this can be compared with the known S-matrix of the models to which the NLIE is referred.
Another important IR observation is that the derivative of the counting function is a very large positive number and no one special root/hole can takes place
$$Z^{}(\vartheta )=l\mathrm{cosh}\vartheta +g^{}(\vartheta |\vartheta _k)\text{for}l\mathrm{}.$$
This suggest that special roots/holes are not physical particles. Instead, as it is clearly shown in section 3.4, they are a mathematical artifact of the logarithmic contribution that appears in the convolution term.
The final important observation in IR analysis comes from quantization conditions (3.44)
$$2\pi I_j=l\mathrm{sinh}\vartheta _j+g(\vartheta _j|\vartheta _k)+\alpha $$
The number $`I_j`$ is real. Distinguishing the real and imaginary part of this equation gives:
$$\begin{array}{c}2\pi I_j=l\mathrm{sinh}\mathrm{}e\vartheta _j\mathrm{cos}\mathrm{}m\vartheta _j+\mathrm{}e(g(\vartheta _j|\vartheta _k)+\alpha )\\ 0=l\mathrm{cosh}\mathrm{}e\vartheta _j\mathrm{sin}\mathrm{}m\vartheta _j+\mathrm{}m(g(\vartheta _j|\vartheta _k)+\alpha )\end{array}$$
To deal with the first equation, the real part of the $`\vartheta _j`$ for increasing values of the large scale $`l`$ moves toward the origin as
$$\mathrm{}e\vartheta _j\frac{1}{l\mathrm{cos}\mathrm{}m\vartheta _j}.$$
From the second equation, the imaginary part must develops a singularity from $`\mathrm{}m(g(\vartheta _j|\vartheta _k))`$ to annihilate the rapidly diverging term $`l\mathrm{sin}\mathrm{}m\vartheta _j`$. Because the positions of singularities are known, this simple argument can be used to fix (at the IR) the imaginary part of the sources. Being the real part zero, the source position at IR is completely fixed.
### 4.4 The intermediate regions
The NLIE (3.54) contains a source term $`g(\vartheta |\vartheta _j)`$ which is specified by giving the root/hole structure of the given state and depends on the positions of the holes, special roots/holes and complex roots. These positions in turn are fixed by the Bethe quantization conditions (3.55). The NLIE supplemented by the quantization conditions gives a set of coupled nonlinear equations in the function $`Z(\vartheta )`$ and the variables $`\vartheta _j`$ that can be solved numerically by an iterative procedure. First one chooses a starting position for the sources $`\vartheta _j`$. Then one iterates the integral equation (using fast Fourier transform to evaluate the convolution) to update the counting function $`Z`$. Using this new $`Z`$, an improved determination of position of the $`\vartheta _j`$ can be obtained, and is fed back into the integral equation for a new iteration cycle. The process is repeated until the solution is found to a prescribed precision (usually $`10^6`$). As explained in section 4.3, for large $`l`$ the source term dominates, while the correction coming from the integral term is exponentially small. Therefore it is reasonable to expect that the further one goes to the IR regime the faster the iteration converges which is in fact what happens in the computations. Hence it is preferable to start iterating at the largest desired value of $`l`$ and decrease the volume gradually, always taking as a starting point at the next value of the volume the solution found at the previous value.
As a preliminar analysis, the general behaviour of the $`Z(x)`$ function is that of $`\mathrm{sinh}x`$ for large $`l`$. Decreasing the scale, the tangent in the central portion of $`Z(x)`$ for $`x`$ becomes more and more horizontal; instead for large $`|x|`$ the leading term is always the hyperbolic sine. This central position broadens to a single or double plateau system that extends to all $`|x|`$’s smaller than $`\mathrm{log}(2/l),`$ and rapidly disappears for greater values, in favor of the dominant exponential growth.
It can happens, when decreasing $`l`$, that special holes appear. But there is the so called *number of effective holes*, that is given by
$$N_{H,eff}=N_H2N_S$$
(4.4)
that is a constant independent of $`l`$. This is a consequence of one simple fact: when a globally increasing function (as $`Z`$ is) makes a small oscillation as in figure 4.1 there are an odd number of points that intersect the horizontal line corresponding to a certain quantum number.
Then to every real solution (when $`Z`$is increasing) correspond three (of five or seven or …) real solutions in the case of specials. Of this odd number of solutions, only one is a real root (no Bethe roots can have the same quantum numbers). All the other must be holes (specials or normal). Then (as in figure 4.1)
$$1\text{ root for large }l=1\text{ root}+2\text{ holes}2\times 1\text{ special r}.\text{ or h}.\text{ for small }l.$$
The general case is (4.4). The physical meaning of the holes, then, must be related to the effective number $`N_{H,eff}`$, not to $`N_H.`$various
For the finite values of the scale, it is interesting to made the comparison of the numerical data from the NLIE predictions with those obtained with the Truncated Conformal Space Approach (TCSA, see section (1.4)) at several values of the parameter $`p`$. For illustration, some interesting cases will be presented. For the various values of the coupling constant, our group obtained a spectacular agreement between the results obtained by the two methods, up to deviations of order $`10^410^3`$(see ). The deviation grows with the volume $`L`$, exactly as expected for truncation errors. As also told in (1.4), in the attractive regime ($`p<1`$) the TCSA convergence is fast. The opposite happens for the NLIE, i.e. it converges faster in the repulsive regime ($`p>1`$) and for large volume $`L`$, but in general it is possible to have precisions of order $`10^7`$ (and higher) in both the regimes.
By studying other parameters such as the mass gap, the breather-soliton mass ratios and the rate of convergence of the energy levels with increasing the value of $`E_{cut}`$, the small differences between the two methods can be clearly attributed to the inaccuracy in the TCS data.
### 4.5 The UV limit computation
The ultraviolet limit of the states described by the NLIE will be examined, in order to compare it with the known facts about the UV limit of sG/mTh theory and conformal minimal models, outlined in sections (1.2, 1.3).
More explicitly, it is known, from conformal perturbation theory, that the behaviour of energy and momentum for very small $`L`$ are given by:
$$\begin{array}{c}E(L)=\frac{\pi \stackrel{~}{c}(l)}{6L}=\frac{\pi }{6L}\left(c12(\mathrm{\Delta }+\overline{\mathrm{\Delta }})\right)+\mathrm{}\\ P(L)=\frac{2\pi }{L}(\mathrm{\Delta }\overline{\mathrm{\Delta }})+\mathrm{}\end{array}$$
(4.5)
where the ellipsis are most regular terms in the space size $`L`$. Therefore, to determine the UV behaviour, in $`E`$ and $`P`$ only terms containing $`{\displaystyle \frac{1}{L}}`$ must be retained. Consider, e.g., the hole contribution to the energy: $`\mathrm{cosh}h`$. In order to have this hole, this term must behaves as $`1/L`$, that is possible only if
$$h=\text{finite}\pm \mathrm{log}\frac{2}{l}.$$
(4.6)
This means that only holes with this rapidity can contribute to UV. The same argument applies to specials and complex solutions.
The behaviour of the sources for $`l\mathrm{\hspace{0.17em}0}`$ can be classified by three possibilities: they can remain finite (they are called *central*), or they can move towards the two infinities as $`\pm \mathrm{log}{\displaystyle \frac{2}{l}}`$ (*right/left-movers*). The finite parts of their positions can be obtained by extracting the divergent part:
$$\left\{\vartheta _j(l)\right\}=\{\vartheta _j^\pm \pm \mathrm{log}\frac{2}{l},\vartheta ^0\}.$$
(4.7)
(the $`l`$ dependence of the roots has been made manifest). The number of right/left moving resp. central holes is indicated by $`N_H^{\pm ,0}`$ and similarly the numbers $`N_S^{\pm ,0}`$, $`M_C^{\pm ,0}`$ and $`M_W^{\pm ,0}`$ are introduced. The finite parts in (4.7) satisfy a modified version of the NLIE, known as *kink equation*. To obtain that, observe that $`Z`$ has an implicit dependence from $`l`$ that can be made manifest writing $`Z(\lambda ,l).`$ Define the *kink functions* with the following expression:
$$Z_\pm (\vartheta )=\underset{l0}{lim}Z(\vartheta \pm \mathrm{log}\frac{2}{l},l)$$
(4.8)
The term “kink” has its origin in the fact that the asymptotic form of the function $`Z(\vartheta )`$ has two plateaus stretching between the central region and the regions of the left/right movers. The functions $`Z_\pm (\vartheta )`$ describe the interpolation between the plateaus and the asymptotic behaviour of $`Z(\vartheta )`$ at the corresponding infinity. Note that if there are no central objects then the two plateaus merge into a single one stretching from the left movers’ region to the right movers’ one. Using (4.7, 4.8) in the NLIE (3.54), the source terms may behave in different ways. Consider e.g. one hole $`h(l)`$ of the type $`h^+`$. Its source term for the kink “+” is:
$$\underset{l0}{lim}\chi \left(\vartheta +\mathrm{log}\frac{2}{l}h(l)\right)=\chi (\vartheta h^+).$$
There are $`N_H^+`$ terms of this of type. For a $`h^{}`$ hole the source contribution to the same “+” kink is
$$\underset{l0}{lim}\chi \left(\vartheta +\mathrm{log}\frac{2}{l}h(l)\right)=\underset{l0}{lim}\chi \left(\lambda h^{}+2\mathrm{log}\frac{2}{l}\right)=\chi _{\mathrm{}}$$
($`\chi _{\mathrm{}}`$ is defined in (3.33)). One $`h^0`$ hole behaves in the same way. Then there are $`N_HN_H^+=N_H^{}+N_H^0`$ such terms. Analogous arguments apply to “–” kink and to the other roots. For wide roots the limits of second determination are given in (3.35).
The following definitions are introduced, according to (3.52):
$$S^{\pm ,0}=\frac{1}{2}[N_H^{\pm ,0}2N_S^{\pm ,0}M_C^{\pm ,0}2M_W^{\pm ,0}\theta (p1)]$$
Observe that $`S`$ is always integer, whereas $`S^{\pm ,0}`$ can be half integer; moreover, in some cases, $`S^0`$ can be negative. All this numbers are interpreted as the spin of the right, left movers and fixed solutions. Clearly
$$S=S^++S^{}+S^0.$$
Using all this computations in the definition (4.8), the NLIE for the “+” and “–” kinks reads:
$$Z_\pm (\vartheta )=\pm e^{\pm \vartheta }+\alpha +g_\pm (\vartheta )+𝑑xG(\vartheta x)𝒬_\pm (x)$$
(4.9)
where the following definitions are used:
$$\begin{array}{c}g_\pm (\vartheta )=\underset{l0}{lim}g\left(\vartheta \pm \mathrm{log}\frac{2}{l}|\vartheta _k\right)=\pm 2\chi _{\mathrm{}}(SS^\pm )+2\pi l_W^\pm +\underset{k=1}{\overset{N_H^\pm }{}}\chi (\vartheta h_k^\pm )+\\ \\ \underset{k=1}{\overset{N_S^\pm }{}}\left(\chi (\vartheta \widehat{y}_k^\pm +i\eta )+\chi (\vartheta \widehat{y}_k^\pm i\eta )\right)\underset{k=1}{\overset{M_C^\pm }{}}\chi (\vartheta c_k^\pm )\underset{k=1}{\overset{M_W^\pm }{}}\chi _{II}(\vartheta w_k^\pm );\\ \\ l_W^\pm =\pm \text{sign}(p1)\frac{1}{2}(M_WM_W^\pm )\end{array}$$
(4.10)
and the function $`𝒬_\pm (x)`$ is related to $`Z_\pm `$ as $`Z`$ to $`𝒬`$ in (3.40).
This equations allow to write the quantization conditions as:
$$Z_\pm \left(\vartheta _j^\pm \right)=2\pi I_j^\pm .$$
(4.11)
It is a matter of convenience to put the apex $`\pm `$ on the quantum numbers. Indeed, they do not change in the limit procedure: they are exactly the same as for finite $`l`$.
By simply taking the limit $`l\mathrm{\hspace{0.17em}0}`$ in NLIE reads an equation for the fixed objects $`\vartheta ^0`$:
$$Z_0(\vartheta )=\underset{l0}{lim}Z(\vartheta ,l)=\alpha +g_0(\vartheta )+𝑑xG(\vartheta x)𝒬_0(x)$$
(4.12)
where
$$\begin{array}{c}g_0(\vartheta )=\underset{l0}{lim}g\left(\vartheta |\vartheta _k\right)=2\chi _{\mathrm{}}(S^{}S^+)+2\pi l_W^0+\underset{k=1}{\overset{N_H^0}{}}\chi (\vartheta h_k^0)+\\ \\ \underset{k=1}{\overset{N_S^0}{}}\left(\chi (\vartheta \widehat{y}_k^0+i\eta )+\chi (\vartheta \widehat{y}_k^0i\eta )\right)\underset{k=1}{\overset{M_C^0}{}}\chi (\vartheta c_k^0)\underset{k=1}{\overset{M_W^0}{}}\chi _{II}(\vartheta w_k^0);\\ \\ l_W^0=\text{sign}(p1)\frac{1}{2}(M_W^{}M_W^+).\end{array}$$
As in the previous case the function $`𝒬_0(x)`$ is related to the corresponding $`Z_0(x)`$ by the usual expression (3.40).
In the following the asymptotic values of $`Z_{\pm ,0}`$ and the corresponding values for $`𝒬`$ will play an important role. Equations (4.9) simply gives
$$Z_+(+\mathrm{})=+\mathrm{}=Z_{}(\mathrm{}),$$
(4.13)
consequently the value for $`𝒬(x)`$ can be obtained using (3.40):
$$𝒬_\pm (\pm \mathrm{})=0$$
(4.14)
(to compute it, before do the limit in $`x`$ and later the limit in $`\eta `$, because the NLIE holds with a finite $`\eta `$). Moreover it can be shown that
$$\begin{array}{c}g_+(\mathrm{})=g_0(+\mathrm{})=2\chi _{\mathrm{}}(S2S^+)+2\pi k_W^+\\ g_{}(+\mathrm{})=g_0(\mathrm{})=2\chi _{\mathrm{}}(S2S^{})+2\pi k_W^{}\end{array}$$
and
$$k_W^\pm =\pm \frac{1}{2}\text{ sign }(p1)(M_W2M_W^\pm )=\frac{M_{SC}}{2}+\text{integer}.$$
A direct consequence of this is that
$$Z_+(\mathrm{})=Z_0(+\mathrm{})\text{and}Z_{}(+\mathrm{})=Z_0(\mathrm{})$$
(4.15)
because they satisfy the same equation. Remember that for small $`l`$ the counting function gives, in the central region, one or two plateaus that stretch to the infinities. Then the previous equations simply mean that $`Z_0(+\mathrm{})`$ is the right plateau height and and $`Z_0(\mathrm{})`$ the left plateau height. If there are no central objects, they take the same value, yielding so a one plateau system.
The explicit values can be computed, using the kink equations. The integral can be handled because the $`𝒬`$ in the asymptotic limit takes a constant value and the kernel $`G`$ rapidly converges to zero:
$$\begin{array}{c}Z_+(\mathrm{})=\alpha +g_+(\mathrm{})+\frac{\chi _{\mathrm{}}}{\pi }𝒬_+(\mathrm{})\\ Z_{}(+\mathrm{})=\alpha +g_{}(+\mathrm{})+\frac{\chi _{\mathrm{}}}{\pi }𝒬_{}(+\mathrm{}).\end{array}$$
(4.16)
They are called *plateau equations*. From the definition of $`𝒬`$ in (3.40), and making attention to the warning concerning the range of values allowed for it, it follows that
$$Z_\pm (\mathrm{})=𝒬_\pm (\mathrm{})+\pi \delta +2\pi k_\pm $$
with an appropriate choice of the integers $`k_\pm `$ such that the condition $`\pi 𝒬_+(\mathrm{})\pi `$ holds (remember that $`k_W^\pm `$ can be half-integer).The solution of the plateau equation is then
$$\omega _\pm 𝒬_\pm (\mathrm{})=\pm 2\pi \frac{p1}{p+1}(S2S^\pm )+2\pi \frac{p}{p+1}\left(\frac{\alpha }{\pi }\delta +2k_W^\pm 2k_\pm \right).$$
(4.17)
Notice that for $`p>1`$ there can be cases without solution. If instead it is there, there is no ambiguity in the choice of $`k_\pm `$, because its contribution is a multiple of $`4\pi {\displaystyle \frac{p}{p+1}}>2\pi `$, that is larger than the range of values allowed for $`𝒬_\pm `$. Instead, for $`p<1`$ the solution always exists, but it can happens that two or more choices of $`k_\pm `$ satisfy the condition on the range of $`𝒬_\pm `$. It will be clear later, that this phenomenon is related to the fact that wide roots become excitations independent of the other one’s, which is manifestly evident from the fact that they no longer contribute to the spin of the state $`S`$ and that the asymptotic value of their source contribution $`\chi _{II}`$ in the attractive regime is simply a $`\pi `$ (see 3.35).
Before to continue the general computation, the kink equations for the wide roots will be explicitly written. They can be obtained by using the same procedure used for (4.9), applied to the correct expression for $`Z(\vartheta )`$ given in (3.38):
$$\begin{array}{c}Z_\pm (\vartheta )=\pm e_{II}^\vartheta +\underset{k=1}{\overset{N_H^\pm }{}}\chi _{II}(\vartheta h_k^\pm )\underset{k=1}{\overset{N_S^\pm }{}}\left(\chi _{II}(\vartheta \widehat{y}_k^\pm +i\eta )+\chi _{II}(\vartheta \widehat{y}_k^\pm i\eta )\right)\\ \\ \underset{k=1}{\overset{M_C^\pm }{}}\chi _{II}(\vartheta c_k^\pm )\underset{k=1}{\overset{M_W^\pm }{}}\chi _{II}(\vartheta w_k^\pm )_{II}\pm \theta (p1)4\chi _{\mathrm{}}(SS^\pm )+2\pi L_W^\pm (\text{sign}\mathrm{}m\vartheta )+\\ \\ +𝑑xG_{II}(\vartheta x)𝒬(x)+\alpha _{II}\text{for }|\mathrm{}m\vartheta |>\mathrm{min}(1,p)\end{array}$$
(4.18)
where the notation introduced for $`L_W`$ is so complicated because it depends on the sign of the imaginary part of $`\vartheta `$:
$$L_W^\pm (\text{sign}\mathrm{}m\vartheta )=\underset{k=1}{\overset{M_WM_W^\pm }{}}\underset{x\pm \mathrm{}}{lim}\left[\chi _{II}(\vartheta w_k+x)_{II}4(\chi _{\mathrm{}}\pi )\theta (p1)\right]=\pi \text{ }\text{integer}$$
(4.19)
The explicit form is quite complicated. Then this contribution must be computed case by case. Observe that the term $`e_{II}^\vartheta `$ vanishes for $`p>1`$ and $`\alpha _{II}`$ vanishes for $`p<1`$.
Now, the ultraviolet limit on the energy and momentum expressions will be done. To this end, substitute in the energy and momentum expressions (3.58, 3.59) the ultraviolet behaviour of the roots (4.7) and retain only terms in $`{\displaystyle \frac{1}{L}}`$. The kinetic terms give contributions as
$$\mathrm{cosh}\vartheta _k\frac{1}{L}e^{\pm \vartheta _k^\pm }.$$
(4.20)
The integral terms must be calculated separately. Their form in the energy and momentum expressions is, respectively:
$$\begin{array}{c}\frac{dx}{2\pi }\mathrm{sinh}x𝒬(x)=\frac{dx}{2\pi }\frac{1}{2}[e^x𝒬(x)e^x𝒬(x)],\\ \frac{dx}{2\pi }\mathrm{cosh}x𝒬(x)=\frac{dx}{2\pi }\frac{1}{2}[e^x𝒬(x)+e^x𝒬(x)],\end{array}$$
where the contribution depending on $`e^x`$ is called *“+” kink contribution*, and the other term is the *“–” kink contribution*. From the definition itself of $`𝒬_\pm `$ observe that, in the limit of very small $`l`$, it is possible to write for $`𝒬`$:
$$𝒬(x\pm \mathrm{log}\frac{2}{l},l)𝒬_\pm \left(x\right)+q_\pm (x,l),$$
(4.21)
where the two functions $`q_\pm (x,l)`$ vanishes in the $`l\mathrm{\hspace{0.17em}0}`$ limit. The expression (4.21) can be substituted in the integral form for the two kinks. The integration variable can be shifted by $`xx+\mathrm{log}(2/l)`$ then
$$\frac{dx}{2\pi }\frac{1}{2}e^x𝒬(x)\frac{1}{L}\frac{dx}{2\pi }e^x𝒬_+(x),$$
(4.22)
where the symbol $``$ means that only the terms of order $`1/L`$ are retained. Similarly for the “–” kink term
$$\frac{dx}{2\pi }\frac{1}{2}e^x𝒬(x)\frac{1}{L}\frac{dx}{2\pi }e^x𝒬_{}(x).$$
(4.23)
At this point, it is possible to express energy and momentum in a way dependent only on quantities which are finite in the UV limit. Using then (4.20, 4.22, 4.23) in (4.5), and rearranging the expression in terms of $`\mathrm{\Delta },\overline{\mathrm{\Delta }}`$ (i.e. restoring the light-cone coordinates) follows:
$$\begin{array}{c}\mathrm{\Delta }^\pm =\frac{c}{24}+\frac{1}{2\pi }(\underset{j=1}{\overset{N_H^\pm }{}}e^{\pm h_j^\pm }\underset{j=1}{\overset{N_S^\pm }{}}(e^{\pm \widehat{y}_j^\pm +i\eta }+e^{\pm \widehat{y}_j^\pm i\eta })+\\ \\ \underset{j=1}{\overset{M_C^\pm }{}}e^{\pm c_j^\pm }\underset{j=1}{\overset{M_W^\pm }{}}e_{II}^{\pm w_j^\pm }\frac{dx}{2\pi }e^{\pm x}𝒬_\pm (x)).\end{array}$$
As a notation, $`\mathrm{\Delta }=\mathrm{\Delta }^+`$ and $`\overline{\mathrm{\Delta }}=\mathrm{\Delta }^{}`$. The NLIE can be used now to express the various terms appearing in the sums and the integrals. Consider, for example, the quantization condition (4.11) for the holes $`h_j^\pm `$. Using the oddity of the $`\chi (x)`$ function, it can be arranged in a more convenient way:
$$\begin{array}{c}\pm \underset{j=1}{\overset{N_H^\pm }{}}e^{\pm h_j^\pm }=2\pi \underset{j=1}{\overset{N_H^\pm }{}}I_{h_j}^\pm +\underset{j=1}{\overset{N_H^\pm }{}}(\underset{k=1}{\overset{N_S^\pm }{}}(\chi (h_j^\pm \widehat{y}_k^\pm +i\eta )+\chi (h_j^\pm \widehat{y}_k^\pm i\eta ))+\\ \\ +\underset{k=1}{\overset{M_C^\pm }{}}\chi (h_j^\pm c_k^\pm )+\underset{k=1}{\overset{M_W^\pm }{}}\chi _{II}(h_j^\pm w_k^\pm ))N_H^\pm (\alpha \pm 2\chi _{\mathrm{}}(SS^\pm )+2\pi l_W^\pm )+\\ \\ \underset{j=1}{\overset{N_H^\pm }{}}2\mathrm{}m\frac{dx}{i}G(h_j^\pm xi\eta )\mathrm{log}\left(1+()^\delta e^{iZ_\pm (x+i\eta )}\right)\end{array}$$
(the nonzero $`\eta `$ has been restored for later convenience). For the other objects, similar forms hold (remember that the $`\widehat{y}_k`$ are defined by (3.27)). All the equations so obtained must be substituted in the expression for $`\mathrm{\Delta }^\pm `$. Now, observe that all the sums of $`\chi (\vartheta )`$, $`\chi _{II}(\vartheta )`$ and $`\chi _{II}(\vartheta )_{II}`$ cancel completely for the oddity of the functions. Moreover in the integral term the following substitution (for all the type of sources) can be made
$$2\pi G(h_j^\pm x)=\chi ^{}(xh_j^\pm ).$$
(4.24)
The result is:
$$\begin{array}{c}\mathrm{\Delta }^\pm =\frac{c}{24}\pm \left(I_H^\pm 2I_S^\pm I_C^\pm I_W^\pm \right)\frac{2\chi _{\mathrm{}}}{\pi }\left(SS^\pm \right)S^\pm S^\pm \left(\frac{\alpha }{\pi }+l_W^\pm \right)+\\ \pm _W^\pm 2\mathrm{}m\frac{dx}{(2\pi )^2i}\phi _\pm ^,(x+i\eta )\mathrm{log}\left(1+()^\delta e^{iZ_\pm (x+i\eta )}\right).\end{array}$$
The following notation has been introduced:
$$I_H^\pm =\underset{j=1}{\overset{N_H^\pm }{}}I_{h_j}^\pm ,I_C^\pm =\underset{j=1}{\overset{M_C^\pm }{}}I_{c_j}^\pm ,I_W^\pm =\underset{j=1}{\overset{M_W^\pm }{}}I_{w_j}^\pm \mathrm{and}I_S^\pm =\underset{j=1}{\overset{N_S^\pm }{}}I_{y_j}^\pm .$$
The new constant appearing in that equation takes into account for wide roots, and is given by (see also (4.18)):
$$\begin{array}{c}_W^\pm =\underset{j=1}{\overset{M_W^\pm }{}}L_W^\pm (\text{sign}\mathrm{}mw_j^\pm )=\pi \text{sign}(p1)M_W^\pm 2(SS^\pm )+\\ \underset{j=1}{\overset{M_W^\pm }{}}\underset{k=1}{\overset{M_WM_W^\pm }{}}\underset{x\pm \mathrm{}}{lim}\left[\chi _{II}(w_j^\pm w_k+x)_{II}4(\chi _{\mathrm{}}\pi )\theta (p1)\right]=\pi \text{ }\text{integer}\end{array}$$
It must be computed case by case.
The following function has been introduced:
$$\begin{array}{c}\phi _\pm (\vartheta )=\pm e^{\pm \vartheta }+\underset{k=1}{\overset{N_H^\pm }{}}\chi (\vartheta h_k^\pm )\underset{k=1}{\overset{N_S^\pm }{}}\left(\chi (\vartheta y_k^\pm +i\eta )+\chi (\vartheta y_k^\pm i\eta )\right)\\ \\ \underset{k=1}{\overset{M_C^\pm }{}}\chi (\vartheta c_k^\pm )\underset{k=1}{\overset{M_W^\pm }{}}\chi _{II}(\vartheta w_k^\pm )=Z_\pm (\vartheta )𝑑xG(\vartheta x)𝒬_\pm (x)\alpha \end{array}$$
(4.25)
Now, with the help of the computations shown in appendix C the integral becomes:
$$2\mathrm{}m\frac{dx}{(2\pi )^2i}\phi _\pm ^,(x+i\eta )\mathrm{log}\left(1+()^\delta e^{iZ_\pm (x+i\eta )}\right)=\pm \left(\frac{1}{24}\frac{𝒬_\pm ^2(\mathrm{})}{16\pi ^2}\frac{p+1}{p}\right)$$
and the UV limit computation ends with a close form for the conformal dimensions:
$$\begin{array}{c}\mathrm{\Delta }^\pm =\frac{c1}{24}\frac{\alpha }{\pi }S^\pm \pm \left(I_H^\pm 2I_S^\pm I_C^\pm I_W^\pm \right)S^\pm 2l_W^\pm +\\ \pm _W^\pm \frac{2\chi _{\mathrm{}}}{\pi }\left(SS^\pm \right)S^\pm +\frac{𝒬_\pm ^2(\mathrm{})}{16\pi ^2}\frac{p+1}{p}\end{array}$$
(4.26)
This means that the UV limit admits an exact computation of the spectrum. The identification of the UV states is one of the fundamental steps to understand the physical interpretation of the continuum model so far defined.
An interesting phenomenon is that the conformal weights obtained depend only on very generic features of the source configuration such as the asymptotics of the left and right moving sources, the total spin and the number of self-conjugate roots. This means that if a certain source configuration is given, one can add new sources separately to the right and the left moving part in such a way that they are separately neutral (i.e. do not change $`S^+`$ and $`S^{}`$, the total spin and the number of self-conjugate roots). In this way the primary weights do not change, but generally the term $`I_H^\pm 2I_S^\pm I_C^\pm I_W^\pm `$ is increased and so descendents of the initial state are created. An example: states which have $`S=0=M_{SC}`$ and $`S^\pm =0`$ are all descendants of the vacuum, however complicated their actual source configurations are.
It is important also to drive attention to the well known fact (see the section 3.8) that (at the moment) there is no proof that the scaling functions obtained by the method of NLIE span the complete space of states. This is an extremely difficult problem due to the following two circumstances: (1) the dependence of the UV conformal weights from the parameters of the source configuration is rather complicated and (2) to obtain the allowed values of the complex roots one has to carefully examine the IR limit as suggested in section 4.3. It will be also clear that the same state can be realized by different root configurations depending on the regime (e.g. the $`(s\overline{s})_\pm `$ states: scattering of soliton and antisoliton in even and odd wave functions; see later).
### 4.6 Sine-Gordon and massive Thirring
As suggested in section (3.7), the choice $`\omega =\alpha =0`$ is supposed to describe the sine-Gordon and massive Thirring models on a cylinder. In this section the corresponding NLIE will be analyzed, starting with the vacuum state. Obviously, the coupling $`p`$ is expected to be the same that appears in the s-G and mTh lagrangians via the equation (1.12).
Starting from a lattice with $`2N`$ sites one finds that the antiferromagnetic ground state which has spin $`S=0`$, when written in terms of Bethe vectors depends on $`M=N`$ roots of the Bethe equations, all of which are real, as indicated in (2.5). This ground state is expected to correspond to the vacuum of the field theory. However, there are two possibilities: one for $`N`$ even and the other for $`N`$ odd, corresponding in the continuum to the choices $`\delta =0`$ or $`\delta =1`$. As it will be shown in the sequel, only one of these states can be identified with the vacuum for a local field theory having a $`c=1`$ UV limiting CFT. The UV dimensions of the vacuum are $`\mathrm{\Delta }^\pm =0`$ (because the theory is assumed unitary). Choosing $`\delta =0`$, the expression (4.26) gives a value consistent with the interpretation proposed, i.e.
$$c=1,$$
as obtained in the paper . At the IR, because no holes neither complex roots are considered, the energy and momentum (3.58, 3.59) vanishes, that is what is expected for the vacuum. The numerical iteration of the NLIE gives the scaling energy shown in figure 4.2.
In that figure, the value of the predicted bulk energy (1.21) has been subtracted from TCSA data in order to normalize them at the same way as the NLIE data are. The agreement of the two data is to order $`10^410^3`$.
As a consequence of this vacuum analysis, the equations for the asymptotic behaviour (4.17) and for the conformal dimensions (4.26) take a simpler form ($`\alpha =0`$ and $`c=1`$). If the physical interpretation is consistent, the expected conformal dimensions must have the form (1.4), that is a sum of powers of $`R^2=\frac{p+1}{2p}`$ (as is (1.12)). Then it is convenient to express $`\mathrm{\Delta }^\pm `$ in powers of $`\frac{p+1}{p}`$, using also the expression (4.17) for $`𝒬_\pm ^2(\mathrm{})`$. The final form is
$$\begin{array}{c}\mathrm{\Delta }^\pm =\frac{p}{p+1}n_\pm ^2+\frac{m^2}{16}\frac{p+1}{p}\pm \frac{n_\pm m}{2}+N_\pm =\\ =\frac{1}{2}\left(\frac{n_\pm }{R}\pm \frac{1}{2}mR^2\right)^2+N_\pm \end{array}$$
(4.27)
where the following identifications have been made:
$$\begin{array}{c}m=2S\\ n_\pm =\left(\frac{\delta }{2}k_W^\pm +k_\pm \right)\left(S2S^\pm \right)\end{array}$$
(4.28)
and
$$N_\pm =\pm \left(I_H^\pm 2I_S^\pm I_C^\pm I_W^\pm \right)S^\pm 2l_W^\pm \pm _W^\pm 2\left(S^\pm \right)^22S^\pm \left(\frac{\delta }{2}k_W^\pm +k_\pm \right).$$
(4.29)
The similarity with (1.4) is manifest. To match an exact correspondence two conditions are required. The first one requires that
$$|n_+|=|n_{}|.$$
(4.30)
The second condition is that one must be able to chose the sign of $`n`$ in such a way that the left and right descendent numbers are integer. Consider the case where $`n_+=n_{}`$: then (using (4.27)) follows that both $`N_\pm `$ must be integers. Instead, in the case $`n_+=n_{}`$ a combination of $`N_\pm `$ and $`\pm {\displaystyle \frac{n_\pm m}{2}}`$ must be integer.
Since the $`R^2`$ contribution comes exclusively from the last term in (4.26), the one-plateau systems (i.e. when $`𝒬_+(\mathrm{})=𝒬_{}(+\mathrm{})`$) satisfy trivially the condition (4.30), but even in that case it is not clear whether the descendent numbers are integers. At the time of this writing the analysis of our group cannot exclude the possibility that there exist some states for which the conformal weights cannot be interpreted within the framework of $`c=1`$ CFT. However, in the numerous explicit examples we have calculated so far we have not found any configuration of sources which lead to such behaviour. If it happens for some configuration of the sources then the corresponding scaling function must be excluded from the spectrum.
From the second line of (4.28) one obtains a useful relation ($`m=2S`$)
$$n_\pm =\frac{m+\delta +M_{SC}}{2}+\text{integers}.$$
Using now the classification of UV conformal operators for s-G and mTh in the two algebras $`𝒜_b`$ and $`𝒜_f`$ given in section 1.3, one can obtain the following rule:
$$\begin{array}{c}\frac{m+\delta +M_{SC}}{2}\text{for}\text{sine Gordon}\\ \frac{\delta +M_{SC}}{2}\text{for}\text{massive Thirring}.\end{array}$$
(4.31)
With reference to figure 1.2, from this rule follows that all the four sectors can be accessed by NLIE, because the sectors I and II describe sine-Gordon, I and III describe massive Thirring and the sector IV can be obtained by $`m`$ even and $`{\displaystyle \frac{\delta +M_{SC}}{2}}`$ half-integer. If there are no self-conjugate roots, the rule simplifies: $`\delta =0`$ for mTh (I and III) and $`\delta =1`$ for the sectors II and IV. The sector IV, that also is accessible by NLIE, contains non local states. Then NLIE describes also this non local sector.
#### 4.6.1 “–“ vacuum
Consider the case of $`g(\vartheta |\vartheta _j)=0`$ (i.e. no sources). There are two possibilities $`\delta =0,\mathrm{\hspace{0.17em}1}`$. One is the vacuum, considered previously, the other is $`\delta =1`$. The (4.17) yields
$$𝒬_\pm (\mathrm{})=2\pi \frac{p}{p+1}(1+2k_\pm ),$$
(4.32)
which admits solution only in the attractive regime when
$$|𝒬_+(\mathrm{})|=|𝒬_{}(+\mathrm{})|=2\pi \frac{p}{p+1}.$$
(4.33)
The value is not determined uniquely due to the contribution of wide roots, but it can be fixed using information from the repulsive regime.
In that case, by performing the usual iteration procedure described in section 4.4, one finds that only for $`l>l_0`$ the iteration converges. For smaller values of $`l`$ there is no convergence. The “breakdown” volume $`l_0`$ depends by $`p`$: for $`p=1.5`$ it is $`l_0=0.3`$ (no breakdown in attractive regime). What happens is that the real root at the origin is a special one with $`Z^{}(0)<0`$. By (3.52), because $`S=0`$, also two holes must appear. All of which sources have Bethe quantum number zero. One of the holes moves to the left and the other one to the right and so $`S^\pm ={\displaystyle \frac{1}{2}}`$, while the special root is central (the configuration is completely symmetric, then $`Z(x)=Z(x)`$). This allows for a unique solution of the plateau equations:
$$𝒬_+(\mathrm{})=\pm \frac{2\pi }{p+1},Z_+(\mathrm{})=Z_{}(+\mathrm{})=\pi \frac{1p}{p+1}<0.$$
Observe that $`Z_{}(+\mathrm{})>Z_+(\mathrm{})`$. This compared with (4.15) means that there are two plateaus and the left is higher than the right, i.e. $`Z^{}(0)<0`$. The picture is completely consistent. This phenomenon can be thought of as the splitting of the root at the origin into a special root and two holes and is one of the general mechanisms in which the special sources are generated (see figure 4.3). Another mechanism will be explored when the UV limit of the soliton-antisoliton states will be examined.
Let us suppose that the counting function $`Z`$ is analytic as a function of $`p`$ which is plausible because all terms in the NLIE for this state are analytic in the coupling. This determines which branch to choose for the plateau values in the attractive regime. One obtains the final result (for both the attractive and the repulsive regime)
$$\mathrm{\Delta }^\pm =\frac{1}{4}\frac{p}{p+1}=\frac{1}{8R^2}=\mathrm{\Delta }_{\pm 1/2,0}^\pm ,$$
which are the conformal weights of the vertex operators $`V_{(\pm 1/2,0)}`$. The actual UV limit can be a linear combination of these operators; in any case, it is not contained in the UV spectrum of the sG/mTh theory, but is a state in the nonlocal sector (sector IV in figure 1.2).
#### 4.6.2 Pure hole states
The analysis of pure holes states is the simplest one. In the IR limit one such state gives (using (4.2):
$$Z(\vartheta )=l\mathrm{sinh}\vartheta +\underset{j=1}{\overset{N_H}{}}\chi (\vartheta h_j),Z(h_j)=2\pi I_j.$$
(4.34)
The holes “rapidities” are real. Using the observation (4.1) the equations (4.34) show the same structure of the box quantization equations for physical particles (4.3). Then the holes can be interpreted as particles with rapidity $`h_j`$, i.e. they are the solitons of sine-Gordon. Pure holes states are states with only solitons.
The IR limit puts no restriction on the quantum numbers $`I_k`$ of the holes; instead the rapidities converge to zero. This seems to be a very general feature that remains true even when complex roots are allowed.
The UV calculations, for one hole, give for $`\delta =1`$ (the hole can be put fixed in the origin):
$$\begin{array}{cc}S=\frac{1}{2};\hfill & \hfill S^\pm =0.\end{array}$$
Solving the equation for $`𝒬_\pm (\mathrm{})`$ gives:
$$\begin{array}{c}m=1\\ n_+=n_{}=0\end{array}$$
that correspond to the operators $`V_{0,\pm 1}`$. They create the sine-Gordon soliton (see also ). For one hole with $`\delta =0`$, using
$$\begin{array}{cc}S^+=\frac{1}{2};\hfill & \hfill S^{}=S=0,\end{array}$$
the UV computation gives the dimensions of $`V_{1/2,\mathrm{\hspace{0.17em}1}}`$ or $`V_{1/2,1}`$ that describe the Thirring fermion. The UV picture is consistent with the IR. The behaviour at intermediate values of $`l`$ is shown in figure 4.4
where it appears that the interpretation of the hole with the soliton in consistent. This is the simplest case where the NLIE was used, in , for an odd number of particles.
Table 4.1 gives an idea about the numerical magnitude of the difference. Note that the deviations are extremely small for small values of $`l`$ and they grow with the volume, exactly as expected for truncation errors.
For two solitons case the two possibilities, corresponding to $`\delta =0,1`$ must be checked. The simplest cases are:
1. Two holes quantised with $`I_{1,2}=\pm \frac{1}{2}`$, $`\delta =0`$. It yields:
$$\mathrm{\Delta }^\pm =\frac{p+1}{4p}=\frac{R^2}{2}=\mathrm{\Delta }_{0,2}^\pm ,$$
corresponding to the operator $`V_{(0,2)}`$, which is the UV limit of the lowest-lying two-soliton state. as it can also be seen from the TCS data. If instead of the minimal choice $`I_{1,2}`$ we take a nonminimal one $`I_+=3/2,\mathrm{\hspace{0.17em}5}/2,\mathrm{}`$, $`I_{}=3/2,5/2,\mathrm{}`$, we obtain
$$\mathrm{\Delta }^\pm =\mathrm{\Delta }_{0,2}^\pm \pm I_\pm \frac{1}{2},$$
which corresponds to descendents of $`V_{(0,2)}`$. This is a general phenomenon: the “minimal” choice of quantum numbers yields the primary state, while the nonminimal choices give rise to descendents.
2. Two holes quantized with $`I_{1,2}=\pm 1`$, $`\delta =1`$, as proposed in . In the repulsive regime a special root is required, as in the case of the $`\delta =1`$ vacuum state. The result (in both regimes) is:
$$\mathrm{\Delta }^\pm =1+\frac{1}{4}\frac{1}{p(p+1)}=\mathrm{\Delta }_{1/2,2}^+=\mathrm{\Delta }_{1/2,2}^{}+1.$$
This state is a linear combination of $`\overline{a}_1V_{(1/2,2)}`$ and $`a_1V_{(1/2,2)}`$ which means that it is *not* contained in the local operator algebras of sG/mTh theories.
3. Two holes quantized with $`I_1=0,I_2=\pm 1`$, $`\delta =1`$. Consider in detail the case of $`I_2=+1`$ since the other one is similar. Suppose that the hole with $`I_1=0`$ is a left mover and the other one is a right mover (the other possibilities lead to a contradiction). We find a solution to the plateau equation only in the attractive regime:
$$𝒬_+(\mathrm{})=\pm 2\pi \frac{p}{p+1},𝒬_{}(+\mathrm{})=\pm 2\pi \frac{p}{p+1}.$$
In the repulsive regime the hole with quantum number $`0`$ becomes a special hole $`y`$ and emits other two ordinary holes each quantised with zero. We obtain for the plateau values
$$𝒬_+(\mathrm{})=𝒬_{}(+\mathrm{})=2\pi \frac{1}{p+1}.$$
The conformal weights turn out to be:
$$\begin{array}{cc}\hfill \mathrm{\Delta }^+=& 1+\frac{1}{4}\frac{1}{p(p+1)}=\mathrm{\Delta }_{1/2,2}^+,\hfill \\ \hfill \mathrm{\Delta }^{}=& \frac{1}{4}\frac{1}{p(p+1)}=\mathrm{\Delta }_{1/2,2}^+.\hfill \end{array}$$
These are the conformal weights of the vertex operator $`V_{(1/2,2)}`$. Performing a similar calculation for $`I_2=1`$ we obtain the weights of $`V_{(1/2,2)}`$.
The calculation for a generic number of holes proceeds in complete analogy with the cases treated until now, we only give a summary of the results. The conformal families which we obtain depend on how many holes move to the left and to the right. The primaries are obtained for the minimal choice of the hole quantum numbers; increasing the quantum numbers we obtain secondary states, as it was pointed out in the case with two holes. The states we obtain are in the conformal family of a vertex operator $`V_{n,m}`$ with
$$m=N_{H,eff},n+\frac{\delta }{2}.$$
As a consequence, all states with $`\delta =0`$ are contained in the UV spectrum of the sG/mTh theory while the ones with $`\delta =1`$ are not, in agreement with the relation (4.31). The complete expression for $`n`$ is somewhat complicated, but we give it in the case of symmetric ($`N_H^+=N_H^{}=N_{H,eff}/2`$) configurations:
$$\begin{array}{cc}n=0,\hfill & \delta =0,\hfill \\ n=\pm \frac{1}{2},\hfill & \delta =1.\hfill \end{array}$$
The comparison with TCSA completely confirms the analysis made until now. The example of four solitons states in repulsive regime is in figure 4.5.
As usual in TCSA, one can see that the truncation errors become larger; at values of the volume $`l`$ close to $`10`$ the deviation can be observed even from the figures. For $`l<5`$ the agreement is still within an error of order $`10^3`$. More comments and examples can be found in .
#### 4.6.3 Holes and close roots
Let us now extend our investigation to situations with complex roots and consider the two particle states in more detail. Forgetting for the moment the breathers, we have to consider the two-soliton states. The soliton-antisoliton come in doublets so there are four different polarizations for two particle states of which the $`ss`$ and $`\overline{s}\overline{s}`$ have topological charge $`Q=\pm 2`$ instead $`s\overline{s}`$ and $`\overline{s}s`$ have zero topological charge. The first two states are expected to have exactly the same scaling function for energy and momentum because the sG/mTh theory is charge conjugation invariant. Instead there are two different situations for the neutral $`s\overline{s}`$ state, which have spatially symmetric and antisymmetric wavefunctions (denoted by $`(s\overline{s})_+`$ and $`(s\overline{s})_{}`$, respectively). To separate the symmetric and antisymmetric part one simply has to diagonalize the $`4\times 4`$ SG two particle S-matrix and see that it has 2 coinciding eigenvalues (equal to $`e^{i\chi (\vartheta )}`$, corresponding to $`ss`$ and $`\overline{s}\overline{s}`$), and two different eigenvalues in the $`Q=0`$ channel.
Now we proceed to demonstrate that the IR limit restricts the possible quantum numbers of the complex roots. To simplify matters we consider only the repulsive regime $`p>1`$.
In the repulsive regime $`p>1`$, the scaling function $`(s\overline{s})_{}`$ is realized as the solution to the NLIE with two holes (at positions $`h_{1,2}`$) and a complex pair at the position $`\rho \pm i\sigma `$ . In the IR limit we have
$$\begin{array}{c}\begin{array}{ccc}Z(\vartheta )\hfill & =& l\mathrm{sinh}\vartheta +\chi (\vartheta h_1)+\chi (\vartheta h_2)\chi (\vartheta \rho i\sigma )\chi (\vartheta \rho +i\sigma ),\hfill \end{array}\hfill \\ Z(h_{1,2})=2\pi I_{1,2},\hfill \\ Z(\rho \pm i\sigma )=2\pi I_C^\pm .\hfill \end{array}$$
(4.35)
The quantization condition for the complex roots explicitly reads (we write down the equation only for the upper member of the complex pair, since the other one is similar)
$$l\mathrm{sinh}(\rho +i\sigma )+\chi (\rho +i\sigma h_1)+\chi (\rho +i\sigma h_2)\chi (2i\sigma )=2\pi I_{C.}^+$$
(4.36)
Now observe that as $`l\mathrm{}`$, the first term on the left hand side acquires a large imaginary part, but the right hand side is strictly real. The imaginary contribution should be cancelled by some other term. The function $`\chi `$ is bounded everywhere except for isolated logarithmic singularities on the imaginary axis. For the cancellation of the imaginary part the argument $`2i\sigma `$ of the last term on the left hand side should approach one of these singularities (similarly to the analysis in TBA ).
In the repulsive regime, taking into account that for a close pair $`\sigma <\pi `$, the only possible choice for $`\sigma `$ is to approach $`{\displaystyle \frac{\pi }{2}}`$ as $`l\mathrm{}`$. The soliton-soliton scattering amplitude has a simple zero at $`\vartheta =i\pi `$ with a derivative which we denote by $`C`$ (the exact value does not matter). To leading order in $`l`$, the cancellation of the imaginary part reads
$$l\mathrm{cosh}\rho +\mathrm{}e\mathrm{log}C\left(\sigma \frac{\pi }{2}\right)=0,$$
(4.37)
from which we deduce
$$\left|\sigma \frac{\pi }{2}\right|\mathrm{exp}\left(l\mathrm{cosh}\left(\rho \right)\right),$$
(4.38)
so the imaginary part of the complex pair approaches its infrared limit exponentially fast. This approach is modified by taking into account the finite imaginary contributions coming from the source terms of the holes and from the convolution term. These contributions lead to corrections of the order $`e^l`$ and so they modify only the value of the constant $`C`$.
The behaviour of the real part near the singularity is:
$$\mathrm{}e\chi (zi\pi )\mathrm{\hspace{0.17em}0}\text{ if }\mathrm{}ez<0.$$
Assume, for the moment, that $`\sigma `$ moves toward $`\frac{\pi }{2}`$ from below (i.e. $`\sigma \frac{\pi }{2}<0`$).
For the real part we get, again to the leading order
$$\mathrm{}e\chi \left(\rho +i\frac{\pi }{2}h_1\right)+\mathrm{}e\chi \left(\rho +i\frac{\pi }{2}h_2\right)=2\pi I_{C.}^+$$
(4.39)
It can be shown that (in the repulsive regime $`p>1`$)
$$\xi \left(\vartheta \right)\mathrm{}e\chi \left(\vartheta +i\frac{\pi }{2}\right)=\frac{i}{2}\mathrm{log}\frac{\mathrm{sinh}{\displaystyle \frac{1}{p}}\left(i{\displaystyle \frac{\pi }{2}}\vartheta \right)}{\mathrm{sinh}{\displaystyle \frac{1}{p}}\left(i{\displaystyle \frac{\pi }{2}}+\vartheta \right)},$$
(4.40)
where the the fundamental branch of the logarithm is taken into account. $`\xi `$ is an odd monotonous function bounded by
$$\left|\xi (\vartheta )\right|\left|\xi (\mathrm{})\right|=\frac{\pi (p1)}{2p}$$
(4.41)
from below and above. This means that for any allowed value of $`I_C^\pm `$ the real position of the complex pair is determined uniquely and that
$$\left|I_C^\pm \right|<\left|\frac{p1}{2p}\right|,$$
(4.42)
and since in the repulsive regime $`p>1`$, the only possible choice is $`I_C^\pm =0`$. Then the solution for $`\rho `$ is
$$\rho =\frac{h_1+h_2}{2},$$
(4.43)
so it approaches the central position between the two holes (the corrections to this asymptotic are also exponentially small for large $`l`$). In fact, for the symmetric hole configuration $`I_1=I_2`$ we expect $`h_1=h_2`$ and $`\rho =0`$ to be valid even for finite values $`l`$.
However, the above derivation is valid only if $`\sigma <{\displaystyle \frac{\pi }{2}}`$ and so we do not cross the boundary of the analyticity strip of the $`\chi (\vartheta )`$ function which is at $`\mathrm{}m\vartheta =\pi `$. If $`\sigma >{\displaystyle \frac{\pi }{2}}`$, we must use a consider that:
$$\mathrm{}e\chi (zi\pi )\pi \text{ if }\mathrm{}ez>0.$$
The conclusion that $`\sigma `$ approaches $`\frac{\pi }{2}`$ exponentially fast as $`l\mathrm{}`$ is unaffected by this change, but now the real part of $`\chi (2i\sigma )`$ is not zero but instead $`\pm \pi `$ (modulo $`2\pi `$). We choose the branches of the logarithm in such a way that $`\mathrm{}e\chi (\pm 2i\sigma )=\pi `$ (the structure of the cuts is displayed in figure 4.6). In this way we obtain that
$$I_C^\pm =\frac{1}{2},$$
and the asymptotic value of $`\rho `$ is just as before.
The S-matrix of the $`(s\overline{s})_{}`$ configuration can be obtained by substituting the asymptotic values
$$\rho =\frac{h_1+h_2}{2},\sigma =i\frac{\pi }{2}$$
(4.44)
into the expression for $`Z(\vartheta )`$ and considering now the quantization rules of the holes. We obtain
$$Z(h_1)=l\mathrm{sinh}\left(h_1\right)i\mathrm{log}S_{}(h_1h_2)=2\pi I_1,$$
(4.45)
and a similar equation for the second hole, with
$$S_{}(\vartheta )=\frac{\mathrm{sinh}\left({\displaystyle \frac{\vartheta +i\pi }{2p}}\right)}{\mathrm{sinh}\left({\displaystyle \frac{\vartheta i\pi }{2p}}\right)}S_{++}^{++}(\vartheta ),$$
(4.46)
which is the correct scattering amplitude for the antisymmetric configuration of the soliton-antisoliton system. The amplitude has poles at $`\vartheta =i\pi (12kp),k=1,2,\mathrm{}`$ , corresponding to breathers with mass
$$m_{2k}=2\mathrm{sin}kp\pi .$$
The equation (4.45) describes an approximation to the full NLIE valid for large $`l`$ which is called the dressed or asymptotic Bethe Ansatz. There are three levels of approximations to the scaling functions: the full NLIE (which is in fact exact), the higher level Bethe Ansatz (HLBA) obtained from the NLIE by dropping the convolution term and the asymptotic BA. The difference between the BA and the HLBA is that while the former keeps the complex pair in its asymptotic position, the latter takes into account the corrections coming from the fact that $`\sigma `$ changes with $`l`$ (to first order as given in (4.38) ). However, since the convolution term is of the same order in $`l`$ as the dependence of $`\sigma `$ derived from the HLBA, the HLBA is not a self-consistent approximation scheme. Therefore we are left with the exact NLIE, which is valid for all scales and with the BA as its IR asymptotic form.
The above conclusions about the $`(s\overline{s})_{}`$ state an be extended into the attractive regime as long as $`p>\frac{1}{2}`$. At the point $`p=\frac{1}{2}`$, however, the pair situated at the asymptotic position $`\pm i\frac{\pi }{2}`$ hits the boundary of the fundamental analyticity strip situated at $`i\pi p`$. This phenomenon has a physical origin: this point is the threshold for the second breather bound state which is the first pole entering the physical strip in the $`(s\overline{s})_{}`$ channel. To continue our state further, it requires to go to configurations having an array of complex roots of the first type (see ). Without going into further details, let us mention only that any such array contains a close pair plus some wide pairs depending on the value of $`p`$.
The symmetric state $`(s\overline{s})_+`$ can be obtained from a configuration with two holes and one selfconjugate wide root. Considerations similar to the above lead to the scattering amplitude
$$S_+(\vartheta )=\frac{\mathrm{cosh}\left({\displaystyle \frac{\vartheta +i\pi }{2p}}\right)}{\mathrm{cosh}\left({\displaystyle \frac{\vartheta i\pi }{2p}}\right)}S_{++}^{++}(\theta ),$$
(4.47)
which agrees with the prediction from the exact S-matrix. The poles are now at $`\vartheta =i\pi (1(2k+1)p),k=0,1,\mathrm{}`$ and correspond to breathers with mass
$$m_{2k+1}=2\mathrm{sin}\frac{(2k+1)p\pi }{2}.$$
This configuration extends down to $`p=1`$, where the selfconjugate root collides with the boundary of the fundamental analyticity strip at $`i\pi `$. This is exactly the threshold for the first breather which is the lowest bound state in the $`(s\overline{s})_+`$ channel. For $`p<1`$ the $`(s\overline{s})_+`$ state contains a degenerate array of the first kind (see ), which consists of a close pair, a selfconjugate root and some wide pairs, again depending on the value of $`p`$.
We remark that it is easy to extend the above considerations to a state with four holes and a complex pair. In this case the only essential modification to the above conclusions is that the limit for the complex pair quantum number becomes
$$\left|I_c^\pm \right|<\left|\frac{p1}{p}\right|$$
(4.48)
when $`\sigma <{\displaystyle \frac{\pi }{2}}`$, and
$$\left|I_c^\pm \pm \frac{1}{2}\right|<\left|\frac{p1}{p}\right|,$$
(4.49)
when $`\sigma >{\displaystyle \frac{\pi }{2}}`$, since now there are four $`\chi `$ sources coming from the holes. In the range $`1<p<2`$ this gives the same constraints
$$I_c^\pm =\{\begin{array}{cc}0,\hfill & \sigma <\frac{\pi }{2},\hfill \\ \frac{1}{2},\hfill & \sigma >\frac{\pi }{2}\hfill \end{array},$$
as for the $`(s\overline{s})_{}`$ state (for $`p>2`$ it allows for more solutions, which however we will not need in the sequel).
From the UV point of view, the results shown in this paragraph about the quantum numbers of the close roots are quite interesting. Consider the conformal dimensions (4.26) or the specific form (4.27). There is no proof that the terms $`N_\pm `$ give always the correct descendent numbers. In particular there is no proof that the sum of quantum numbers $`\left(I_H^\pm 2I_S^\pm I_C^\pm I_W^\pm \right)`$ is bounded from below (and nonnegative). Then we don’t know if the so called “vacuum” is really the hamiltonian ground state. The IR analysis done until now shows exactly that states with only holes and with holes and close roots have the correct hamiltonian behaviour.
We now turn ourselves to the UV limits of the previous states. In particular we consider the state of two holes and a close pair (in the repulsive regime), which describes the antisymmetric soliton-antisoliton two-particle state. According to the results from the IR asymptotic, the quantum numbers of the complex roots can be $`0`$ or $`\pm \frac{1}{2}`$. Therefore, we are left two possibilities: either the holes are quantised with integers or with half-integers.
Let us start with the half-integer choice $`\delta =0`$ and suppose that one of the holes is right moving (which is the case if its quantum number $`I_1=I_+`$ is positive) and the other one is left moving (i.e. $`I_2=I_{}<0`$). Using the general formula (4.17) and $`S^\pm =\frac{1}{2}`$, $`S=0`$, by a simple calculation we obtain
$$𝒬_+(\mathrm{})=Z_+(\mathrm{})=2\pi \frac{p1}{p+1},$$
which is valid for $`1<p<4`$ (the complex pair remains central). The other plateau values follow by oddity of the function $`Z`$. For the other half of the repulsive regime we have to include two new normal holes, one moving to the left and the other to the right, and two special holes that remain central (a justification for this will be given shortly) so that $`S^\pm =1`$. We then find
$$𝒬_+(\mathrm{})=4\pi \frac{1}{p+1},Z_+(\mathrm{})=2\pi \frac{p1}{p+1},$$
which is valid for $`p4`$. In both cases the result for the UV weights turns out to be
$$\mathrm{\Delta }^\pm =\frac{p}{p+1}\pm I_\pm \frac{1}{2}=\frac{1}{2R^2}\pm I_\pm \frac{1}{2}=\mathrm{\Delta }_{(\pm 1,0)}\pm I_\pm \frac{1}{2}.$$
The primary state is obtained for the minimal choice $`I_+=I_{}=\frac{1}{2}`$ and coincides with a linear combination of the vertex operators $`V_{(\pm 1,0)}`$ which is correct for the state to be included in the spectrum of the sG/mTh theory and agrees with the behaviour of the $`(s\overline{s})_{_{}}`$ state observed from TCS. If both of the holes move to the right or to the left, we obtain descendents of the identity operator, i.e. states in the vacuum module of the UV CFT.
Figure 4.7 shows how the special holes are generated in this case. Starting from $`p<4`$ and increasing the value of $`p`$ the UV asymptotic form of the counting function varies analytically. Since the real roots/holes are quantised by half-integers, they are at positions where the function $`Z`$ crosses a value of an odd multiple of $`\pi `$. As the plots demonstrate, the behaviour of $`Z`$ is in fact continuous at the boundary $`p=4`$: it is our interpretation in terms of the sources that changes, exactly because we try to keep the logarithm in the integral term of the NLIE in its fundamental branch. The price we pay is that we have to introduce two new normal holes (one moves left and the other one moves right) and two special holes which are central.
Let us comment on the extension to the attractive regime. The complex root configuration changes only by the possible presence of wide pairs, which however have no effect on the plateau values (4.17) as their contribution can be absorbed in redefinition of the values $`k_\pm `$. The other effect the wide roots have is to shift the terms $`\mathrm{\Sigma }_\pm `$ by multiples of $`2\pi `$, but this affects only the integers $`N_\pm `$ in (4.29). Therefore we again obtain states which are descendents of $`V_{(\pm 1,0)}`$.
For the state with integer quantization of the holes, we only give the result. If one of the holes is moving left and the other one to the right, we obtain
$$\mathrm{\Delta }^\pm =\frac{1}{4}\left(\frac{p}{p+1}\right)\pm I_\pm =\frac{1}{8R^2}\pm I_\pm =\mathrm{\Delta }_{(\pm 1/2,0)}\pm I_\pm ,$$
i.e. some linear combinations of descendents of the vertex operators $`V_{(\pm 1/2,0)}`$. If both of the holes move to the left or to the right, we obtain other descendants of the same primary states. Again, this excludes the integer quantised states from the spectrum of sG/mTh theory.
##### 4.6.3.1 Two holes and a selfconjugate complex root
Taking the integer quantised state, the plateau values are given by the same formulas as for the case with the close pair above, since the plateau equation is identical. The only difference is that the numbers $`l_W^\pm `$ (4.10) take a nonzero value
$$l_W^\pm =S^\pm .$$
If the one of the holes is a right mover, while the other one is a left mover, the conformal weights turn out to be
$$\mathrm{\Delta }^\pm =\frac{p}{p+1}\pm I_\pm 1=\mathrm{\Delta }_{(\pm 1,0)}\pm I_\pm 1.$$
If the two holes move in the same direction, we again obtain secondaries of the vacuum state.
Concerning the extension to the attractive regime, one can again check that the plateau solution remains unchanged for the root configuration of two holes and two close roots, therefore the conformal family to which the state belongs remains the same, similarly to the case of the $`(s\overline{s})_{}`$ state.
One can ask what happens if we quantize with half-integers. The result is, similarly to the case with the complex pair, that we obtain descendants of the operators $`V_{(\pm 1/2,0)}`$. Therefore we conclude that in this case the integer quantised configuration must be accepted, while the half-integer one is ruled out, in full accordance with the rule (4.31).
#### 4.6.4 Breather $`S`$-matrices and IR limit
In the infrared limit $`l\mathrm{}`$ the term $`l\mathrm{sinh}(\vartheta )`$ develops a large imaginary part in the first determination away from the real axis, forcing the close complex roots to fall into special configurations called *arrays* (we use the terminology of )*.* An array is a set of complex roots in which the roots are placed at specific intervals in the imaginary direction and have the same real part. In the attractive regime $`l\mathrm{sinh}(\vartheta )_{II}`$ is nonzero and so this is true for nonselfconjugate wide pairs as well (in the repulsive case wide roots do not have such driving force), while self-conjugate roots have a fixed imaginary part anyway. The deviation of the complex roots from their positions in the array decays exponentially with $`l`$ (section 4.6.3).
For the rest of this subsection, whenever it is not explicitly stated, we restrict ourselves to the attractive regime $`p<1`$. The possible arrays fall into two classes:
1. *Arrays of the first kind* are the ones containing close roots, which describe the polarization states of solitons.
There are two degenerate cases: *odd degenerate* arrays, which have a self-conjugate root at
$$\vartheta _0=\vartheta +i\frac{\pi (p+1)}{2}$$
and accompanying complex pairs at
$$\vartheta _k=\vartheta \pm i\frac{\pi (1(2k+1)p)}{2},k=0,\mathrm{},\left[\frac{1}{2p}\right]$$
and *even degenerate* ones, which only contain complex pairs, at the positions
$$\vartheta _k=\vartheta \pm i\frac{\pi (12kp)}{2},k=0,\mathrm{},\left[\frac{1}{2p}\right]$$
These arrays always contain exactly one close pair. The odd degenerate arrays in the repulsive regime reduce to single self-conjugate roots and the even degenerate ones to a single close complex pair.
2. *Arrays of the second kind* describe breather degrees of freedom. The odd ones contain a self-conjugate root
$$\vartheta _0=\vartheta +i\pi (p+1)/2$$
and wide pairs as follows:
$$\mathrm{}m\vartheta _k=\vartheta \pm i\frac{\pi (1(2k+1)p)}{2},k=0,\mathrm{},s,$$
where
$$0s\left[\frac{1}{2p}\right]1,$$
while the even ones only contain wide pairs
$$\mathrm{}m\vartheta _k=\vartheta \pm i\frac{\pi (12kp)}{2},k=0,\mathrm{},s,$$
and $`s`$ runs in the same range. They correspond to the $`(2s+1)`$-th breather $`B_{2s+1}`$ and the $`(2s+2)`$-th breather $`B_{2s+2}`$, respectively.
As one can see, arrays of the second kind become degenerate ones of the first kind, if we analytically continue increasing $`p`$. The reason is that breathers are of course soliton-antisoliton bound states, while degenerate arrays of the first kind describe scattering states of a soliton and antisoliton, as we will see shortly.
One can compute the energy and momentum contribution of a array of the second kind corresponding to the breather $`B_s`$. The energy-momentum contribution turns out to be
$$2\mathrm{sin}\frac{\pi sp}{2}(\mathrm{cosh}\vartheta ,\mathrm{sinh}\vartheta ),$$
(4.50)
where $`\vartheta `$ is the common real part of the roots composing the array. This is just the contribution of a breather $`B_s`$ moving with rapidity $`\vartheta `$. Arrays of the first kind do not contribute to the energy-momentum in the infrared limit, which lends support to their interpretation as polarization states of solitons.
Now we proceed to show that with the above interpretation the NLIE correctly reproduces the two-body scattering matrices of sine-Gordon theory including breathers.
Let us start with breather-soliton matrices. The Bethe quantization conditions for a state containing a soliton (i.e. a hole) with rapidity $`\vartheta _1`$ and a breather $`B_s`$ with rapidity $`\vartheta _2`$ take the following form in the infrared limit. For the hole we get
$$Z(\vartheta _1)=\mathrm{sinh}\vartheta _1\underset{k=0}{\overset{s}{}}\chi _{II}(\vartheta _1\vartheta _2i\rho _k)=2\pi I_1,$$
where we denoted the prescribed imaginary parts of the roots in the array $`B_s`$ by $`\rho _k`$. Here we used $`\chi (0)=0`$ to eliminate the source term for the hole. Now we can compute the second determination of $`\chi `$ as in (3.34). Now it is a matter of elementary algebra to arrive at
$$Z(\vartheta _1)=\mathrm{sinh}\vartheta _1i\mathrm{log}S_{SB_s}(\vartheta _1\vartheta _2)=2\pi I_1,$$
where $`S_{SB_s}(\vartheta _1\vartheta _2)`$ is the soliton-breather $`S`$-matrix conjectured in .
One can start with the breather quantization conditions, too. Writing
$$Z(\vartheta _2+i\rho _k)=\mathrm{sinh}(\vartheta _2+i\rho _k)_{II}+\chi _{II}(\vartheta _2\vartheta _1+i\rho _k)+\mathrm{}=2\pi I_2^{(k)},k=0,\mathrm{},s$$
(4.51)
(the dots are terms due to wide root sources themselves, which cancel out up to multiples of $`2\pi `$ in the next step). Summing up these equations one arrives at
$$2\mathrm{sin}\frac{sp\pi }{2}\mathrm{sinh}(\vartheta _2)i\mathrm{log}S_{SB_s}(\vartheta _2\vartheta _1)=2\pi I_2,$$
where $`I_2`$ is essentially minus the sum of the quantum numbers of the wide roots composing the array (shifted by some integer coming from summing up the terms omitted in equation (4.51)).
Using a similar line of argument we also reproduced the breather-breather $`S`$-matrices by writing down the Bethe quantization conditions for a state with two degenerate strings $`B_s`$ and $`B_r`$ of the second kind. One has to be careful that when $`Z(\vartheta )`$ contains wide root sources which are expressed in terms of $`\chi _{II}(\vartheta )`$, the second determinations of these terms will appear in $`Z(\vartheta )`$for large $`\mathrm{}m\vartheta `$, i.e. terms that can be written roughly like $`\left(\chi _{II}(\vartheta )\right)_{II}`$, as in (3.39).
Scattering state of a soliton and an antisoliton can be described by taking two holes and a degenerate array of the first kind. There are two possibilities now, corresponding to scattering in the parity-odd and parity-even channels. Following the procedure outlined in section 4.6.3, we were once again able to reproduce the corresponding scattering amplitudes. The results presented here together with those of section 4.6.3 exhaust all two-particle scattering amplitudes of sine-Gordon theory, both in the repulsive and attractive regime.
#### 4.6.5 Some examples of breather states
So we proceed to take a look at the simplest neutral excited state, which is the one containing a first breather $`B_1`$ at rest. The source term to be written into the NLIE turns out to be
$$g(\vartheta w_R)=\chi _{II}(\vartheta w)=i\mathrm{log}\frac{\mathrm{cos}{\displaystyle \frac{\pi p}{2}}i\mathrm{sinh}(\vartheta w_R)}{\mathrm{cos}{\displaystyle \frac{\pi p}{2}}+i\mathrm{sinh}(\vartheta w_R)},$$
(4.52)
where $`w_R`$ is the real part of the position of the self-conjugate root (the imaginary part is $`w_I=\pi (p+1)/2`$. In this case $`w_R=0`$ since the breather has zero momentum. There is no need to look at the Bethe quantization condition as the root does not move due to the left-right symmetry of the problem. This state is quantized with integer Bethe quantum numbers, i.e. $`\delta =1`$, as in (4.31).
Table 4.2 presents the energy values obtained by iterating the NLIE in comparison to results coming from TCS, at the values $`p=2/7`$ and $`p=2/9`$.
The table shows that iteration of the NLIE fails for values of $`l`$ less than $`3`$ (the actual limiting value is around $`2.5`$). What is the reason?
We plot the counting function $`Z`$ on the real line for $`l=3`$ and $`l=2`$ in figure 4.8. In a first approximation we can safely neglect the integral term for these values of the volume to see the qualitative features that we are interested in.
What we see is that the behaviour of the function changes: its derivative changes sign at the origin. As a result, two new holes appear where the new real zeros of the function are. But the topological charge remains zero, due to the fact that now we have a special root at the origin and so $`N_S=1`$ and $`N_H=2`$. The two new holes do not give us any new dynamical degrees of freedom: their quantum numbers are fixed to be $`0`$ and so their positions are uniquely determined.
Of course when we calculated the UV dimension, the appearance of the new sources had to be taken into account. It turns out that for $`{\displaystyle \frac{1}{3}}<p<1`$ the two holes are left/right movers, while for $`p<{\displaystyle \frac{1}{3}}`$ they remain central. Calculating the UV conformal dimensions we obtain
$$\mathrm{\Delta }^\pm =\frac{p}{p+1},$$
(4.53)
so the ultraviolet limit of this state is a linear combination of the vertex operators $`V_{\pm 1,0}`$ of the $`c=1`$ UV CFT. This is in perfect agreement with the TCS calculations performed by us as in table 4.2.
How does this change of sign in the derivative of $`Z`$ affect the iteration scheme for the NLIE (see section 4.4)? The two new zeros of $`Z(\vartheta )`$, which is a complex analytic function apart from logarithmic branch cuts, actually correspond to singularities of the logarithmic term in the NLIE (3.54). They come along the imaginary axis in the $`\vartheta `$ plane as we decrease $`l`$, and at a certain point they cross our integration contour which runs parallel to the real axis at distance $`\eta `$. As they make the logarithmic term in our NLIE (3.54) singular, they blow up our iteration scheme. After reaching the origin of the $`\vartheta `$ plane (at exactly the radius where the derivative of $`Z`$ becomes $`0`$), they continue to move along the real axis, which corresponds to crossing a square root branch cut.
This problem of numerically solve NLIE appears in all the cases where special roots/holes are taken into account (e.g. the “-“ vacuum, the two holes-close quantized with “-“ and so on).
We do not go into details here as this problem is currently under investigation<sup>2</sup><sup>2</sup>2 Work in progress in collaboration with P. E. Dorey and C. Dunning, Durham. . We just remark that these issues prove to be highly nontrivial and for the time being, unfortunately, they prevent us from having a reliable numeric scheme for the NLIE below the critical volume.
One can estimate the volume where the slope of the counting function changes sign by neglecting the integral of the logarithmic term. The result is
$$l_{critical}=\frac{2}{\mathrm{cos}{\displaystyle \frac{p\pi }{2}}}$$
(4.54)
which gives a value of around $`2.22`$ for $`p=2/7`$ and $`2.13`$ for $`p=2/9`$. The actual limiting value is a bit higher, partly due to the finite value of $`\eta `$ used in the iteration program and partly because the iteration already destabilizes when the singularities come close enough to the contours. It must also be noted that the integral term cannot eventually be neglected when the singularities are close to the contour, which is an additional reason why (4.54) is just a crude estimate.
We make a short digression to examine the UV limit of the second breather $`B_2`$. The second breather at rest is described by a wide pair at positions
$$\vartheta =\pm i\frac{\pi }{2}.$$
Calculating the UV conformal dimension we get
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=\frac{p}{p+1},$$
which turns out to be the same as that of the first one (4.53), i.e. this state must originate from the other linearly independent combination of the vertex operators $`V_{\pm 1,0}`$ in the ultraviolet. This is again in perfect agreement with TCS and confirms a result by Pallua and Prester who used $`XXZ`$ chain in transverse magnetic field to regularize sine-Gordon theory. They calculated scaling functions numerically on a finite lattice for several concrete values of $`p`$, and arrived at this conclusion by looking at the numerical data. However, our method to compute UV dimensions gives us an *exact analytic formula* and therefore much stronger evidence. This result is interesting because it invalidates a conjecture made previously by Klassen and Melzer who identified the second breather as a linear combination of $`V_{\pm 2,0}`$.
To close this section, we present the lowest lying example of a two-breather state, containing two $`B_1`$ particles with zero total momentum. It turns out that this is a state for which the numerical iteration of the NLIE is not plagued with the problem found above for the first breather. Locality constrains the state to be quantized with half-integers and for lowest energy the quantum numbers of the self-conjugate roots must take the values
$$I_1=\frac{1}{2},I_2=\frac{1}{2}$$
We remark that in contrast to the case of holes, the self-conjugate root with $`I>0`$ moves to the left, while the one with $`I<0`$ moves to the right. This is due to the fact that the second determination of $`Z`$ is in general a monotonically decreasing function on the self-conjugate line. In order to determine the position of the two self-conjugate roots we need the second determination of $`Z`$. The second determination of the self-conjugate root source turns out to be
$$\left(\chi _{II}(\vartheta )\right)_{II}=i\mathrm{log}\frac{i\mathrm{sin}\pi p\mathrm{sinh}\vartheta }{i\mathrm{sin}\pi p+\mathrm{sinh}\vartheta }.$$
Up to some signs, this is just the phase shift which arises when two breathers scatter on each other, which is exactly why the IR analysis gives the correct scattering amplitude. Using this formula, we obtained the numerical data presented in table 4.3.
The UV limit for this state can be calculated from NLIE to be a symmetric first level descendent of the vacuum with weights
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=1,$$
which agrees with TCS. (Note that this descendent exists due to the fact that there is a $`\widehat{U}(1)_L\times \widehat{U}(1)_R`$ Kac-Moody symmetry at $`c=1`$: this state exactly corresponds to the combination of the left and right moving currents $`J\overline{J}`$.)
### 4.7 $`\alpha `$ twist and minimal models (the ground state)
It is known from that by twisting the Bethe Ansatz Equations of the six vertex model as indicated in section 2.5, the twisted model shows the critical behaviour of conformal minimal models (the untwisted critical behaviour corresponds to $`c=1`$). Moreover, in recent papers, Al. Zamolodchikov has put forward the idea of modifying sine-Gordon theory by a twist $`\alpha `$ to deal with conformal minimal models. As a consequence of this two ideas, the Bethe equations used to obtain NLIE contained the twist (2.37) $`\omega `$. In the NLIE, as shown in section 3.5, the twist is parametrized by $`\alpha `$ (3.42).
Looking at the ground state, that in analogy with sine-Gordon is expected to be a sea of real roots, the source in NLIE is put to zero $`g(\vartheta )=0`$ (3.37).
In analogy with the sine-Gordon ground state, we can choose half-integer quantization rule with $`\delta =0`$. The expression for conformal dimensions (4.26) gives:
$$\mathrm{\Delta }^\pm =\frac{c1}{24}+\frac{1}{4}\frac{p}{p+1}\left(\frac{\alpha }{\pi }\right)^2$$
corresponding to an effective central charge
$$\stackrel{~}{c}=1\frac{6p}{p+1}\left(\frac{\alpha }{\pi }\right)^2$$
(4.55)
(the effective central charge is defined in (4.5)). Only in the unitary models the $`\stackrel{~}{c}`$ is the Virasoro central charge. Furthermore, it is well-known that the perturbation of the Virasoro minimal model $`Vir(r,s)`$ by its relevant primary operator $`\mathrm{\Phi }_{(1,3)}`$ is integrable and is described by an RSOS restriction of sine-Gordon theory with
$$p=\frac{r}{sr}.$$
(4.56)
We will use for this model the shorthand notation $`Vir(r,s)+\mathrm{\Phi }_{(1,3)}`$. Using the rule suggested in (3.60) for $`\omega `$ and the expression (3.42) for $`\alpha `$ gives: $`\alpha =\pi /r`$ (the additional integer terms are excluded, for the moment). Then
$$\stackrel{~}{c}=1\frac{6}{rs},$$
which is exactly the *effective central charge* of the minimal model $`Vir(r,s)`$. Therefore one can expect that the twisted equation describes the ground state of the model $`Vir(r,s)+\mathrm{\Phi }_{(1,3)}`$. In fact, Fioravanti et al. calculated these scaling functions for the unitary case $`s=r+1`$ and showed that they match perfectly with the TBA predictions already available. Moreover, choosing the following values for the twist
$$\alpha =\pm \frac{k\pi }{r},k=1\mathrm{}r1$$
(4.57)
they obtained the conformal weights of the operators $`\mathrm{\Phi }_{(k,k)},k=1\mathrm{}r1`$ in the UV limit (the sign choice is just a matter of convention). In our notation, $`\mathrm{\Phi }_{(q,q^{})}`$ denotes the primary field with conformal weights
$$h^+=h^{}=\frac{(qsq^{}r)^2(sr)^2}{4sr}.$$
(4.58)
The models $`Vir(r,s)+\mathrm{\Phi }_{(1,3)}`$ have exactly $`r1`$ ground states. In fact, one can see from the fusion rules that the matrix of the operator $`\mathrm{\Phi }_{(1,3)}`$ is block diagonal with exactly $`r1`$ blocks in the Hilbert space made up of states with the same left and right primary weights. In each of these blocks, there is exactly one ground state and for the unitary series $`s=r+1`$, it was conjectured in that their UV limits are the states corresponding to $`\mathrm{\Phi }_{(k,k)}`$. One can check that in the general nonunitary case the twists (4.57) correspond in the UV limit to the lowest dimension operators among each of the $`r1`$ different blocks of primaries (see explicit examples later).
These ground states are degenerate in infinite volume, but for finite $`l`$ they split; their gaps decay exponentially as $`l\mathrm{}`$. In the unitary case, they were first analyzed in the context of the NLIE in where it was shown that the NLIE predictions perfectly match with the TBA results already available for the unitary series.
However, ground states for nonunitary models have not been treated so far and therefore now we proceed to give examples of that. The models we select are the ones that will be used for comparison in the case of excited states as well. The first is for the scaling Lee-Yang model $`Vir(2,5)+\mathrm{\Phi }_{(1,3)}`$, for which we have also given data from TCS and TBA for comparison (table 4.4).
$$M_B=2\mathrm{sin}\frac{\pi p}{2}=\sqrt{3}$$
is the mass of the fundamental particle of the Lee-Yang model (this is more natural here than using the mass $``$ of the soliton of the unrestricted sine-Gordon model as a scale, since the soliton disappears entirely from the spectrum after RSOS restriction) and call
$$l_B=M_BL=\sqrt{3}l$$
where $`l`$ is the variable appearing in NLIE. There is only one independent value of the twist, which we choose to be
$$\alpha =\frac{\pi }{2}.$$
Here and in all other subsequent calculations the TCS data were normalized using the analogue of the coupling-mass gap relation (1.18) from .
There is only one ground state in this model, which corresponds to the primary field with conformal weights
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=\frac{1}{5},$$
which is in agreement with TBA and TCS predictions. We have also found a perfect agreement for the models $`Vir(2,7)+\mathrm{\Phi }_{(1,3)}`$ and $`Vir(2,9)+\mathrm{\Phi }_{(1,3)}`$, but we do not present those data here. We remark that the TCS for the minimal models converges much better than the one for $`c=1`$ theories: all TCS data in table 4.4 and subsequent ones were produced by taking a few hundred states and in some fortunate cases (e.g. the ground state of the scaling Lee-Yang model for small values of $`l`$) we were able to produce data with up to $`910`$ digits of accuracy! The better convergence meant that all the computation could be done with the computer algebra program *Mathematica*, greatly simplifying the programming work.
All the models of the class $`Vir(2,2n+1)+\mathrm{\Phi }_{(1,3)}`$ have only one ground state. For models with two ground states, we can take a look at $`Vir(3,5)`$ ($`Vir(3,7)`$ was taken into account in ). For $`Vir(3,5)`$ the ultraviolet spectrum is defined by the following Kac table, where the weight (4.58) of the field $`\mathrm{\Phi }_{(k,l)}`$ is found in the $`k`$-th row and $`l`$-th column.
| $`0`$ | $`\frac{1}{20}`$ | $`\frac{1}{5}`$ | $`\frac{3}{4}`$ |
| --- | --- | --- | --- |
| $`\frac{3}{4}`$ | $`\frac{1}{5}`$ | $`\frac{1}{20}`$ | $`0`$ |
The two blocks of the perturbing operator $`\mathrm{\Phi }_{(1,3)}`$ are defined by the fields $`\{\mathrm{\Phi }_{(1,2)},\mathrm{\Phi }_{(1,4)}\}`$ and $`\{\mathrm{\Phi }_{(1,1)},\mathrm{\Phi }_{(1,3)}\}`$, respectively. The ground states correspond in the UV to the operators $`\mathrm{\Phi }_{(1,2)}`$ and $`\mathrm{\Phi }_{(1,1)}`$, as can be checked directly using formulae (4.55), (4.56) and (4.57).
We also have TBA data to compare with, using the TBA equation written by Christe and Martins . The lower-lying ground state is obtained directly from their TBA, while for the other we used the idea of Fendley of twisting the TBA equation . The numerical results are presented in tables 4.5 and 4.6.
In the case $`Vir(3,7)`$ treated in , it was possible only to have a comparison with TCS results, but it still looked pretty convincing.
To summarize, we now have sufficient evidence to believe that the $`\alpha `$-twisted NLIE describes the correct scaling functions for ground states of minimal models perturbed by $`\mathrm{\Phi }_{(1,3)}`$ *in unitary and nonunitary case*. However, the NLIE for sine-Gordon is known to work for excited states as well. But how do we get the excited state spectrum of the minimal models now?
### 4.8 $`\alpha `$ twist and minimal models: excited states
#### 4.8.1 The choice of $`\alpha `$
From now on we restrict ourselves to the case of neutral (i.e. $`S=0`$) states. It is easy to see that even for states with a zero charge the relation between $`\alpha `$ and $`\omega `$ is highly nontrivial.
Now we proceed to show that choosing the value of $`\omega `$ as
$$\omega =\frac{k\pi }{p+1}$$
where $`k`$ is integer, we can reproduce all the required values of $`\alpha `$ listed in equation (4.57). Observe that this expression for $`\omega `$ is different from (3.60), but totally equivalent, as will be clear soon. First of all, we substitute the value of $`p`$ from (4.56) to obtain
$$\omega =\frac{k(sr)\pi }{s}.$$
Since $`r`$ and $`s`$ are relative primes, the independent values of $`\omega mod\pi `$ can be written as
$$\omega =\frac{l\pi }{s},l=0,\mathrm{},s1.$$
For $`S=0`$, we can rewrite the formula (3.42) as follows
$$\alpha =\frac{l\pi }{r}+\frac{2rs}{2r}\pi \left(\left[\frac{1}{2}+\frac{l}{s}\right]\left[\frac{1}{2}\frac{l}{s}\right]\right).$$
We are interested only in the value of $`\alpha mod\pi `$, since using the parameter $`\delta `$ one can effectively shift $`\alpha `$ by $`\pi `$. This leaves us with the formula
$$\alpha =\frac{l\pi }{r}\frac{s\pi }{2r}\left(\left[\frac{1}{2}+\frac{l}{s}\right]\left[\frac{1}{2}\frac{l}{s}\right]\right).$$
The first possibility is that $`l<\frac{s}{2}`$, which simply gives us the values
$$\alpha =\frac{l\pi }{r}.$$
When $`s`$ is even, we can have $`l=\frac{s}{2}`$, which gives us $`\alpha =0`$. Finally, when $`l>\frac{s}{2}`$, we get the values
$$\alpha =\frac{(ls)\pi }{r}.$$
It is easy to check that these formulae reproduce every value
$$\alpha =\frac{n\pi }{r}mod\pi $$
at least once, as required by (4.57), using the fact that $`s>r`$ and that the values above form an uninterrupted sequence of $`s`$ numbers (or when $`s`$ is even, of $`s1`$ numbers, the zero repeated) with equal distances $`\frac{\pi }{r}`$.
As we have already seen in the previous section, all the values
$$\alpha =\frac{k\pi }{r},k=1,\mathrm{},r1$$
(4.59)
are necessary to reproduce correctly the $`r1`$ ground states of the model $`Vir(r,s)+\mathrm{\Phi }_{(1,3)}`$.
The twisted lattice Bethe Ansatz was analyzed by de Vega and Giacomini in . On the lattice, passing from the sine-Gordon model to the perturbed Virasoro model amounts to going from the six-vertex model to a lattice RSOS model. In it was shown that to obtain all the states of the RSOS model it is necessary to take all the twists
$$\omega =\frac{k\pi }{p+1}mod\pi $$
into account. The fact that not all these twists correspond to inequivalent values of $`\alpha `$ and so to different physical states is a consequence of the RSOS truncation.
To close this section we remark that the parameter $`\alpha `$ drops out of the second determination of $`Z`$ in the attractive regime. This is important because as a consequence the IR asymptotics of the breather states does not depend on $`\alpha `$ and so the $`S`$-matrices involving breathers are unchanged. In fact, scattering amplitudes between solitons and breathers remain unchanged too, as can be seen from examining the argument that we used to derive them in section 4.6.4. This matches with the fact that the RSOS restriction from sine-Gordon theory to perturbed minimal models does not modify scattering amplitudes that involve two breathers or a breather and soliton .
#### 4.8.2 The UV limit
There is in fact a very simple intuitive argument to show that the states we get from NLIE with an $`\alpha `$ twist are related to the minimal models in the UV limit. Let us first recall that the UV limit of the $`\alpha =0`$ NLIE yields the vertex operators $`V_{(n,m)}`$ and their descendants (4.27, 4.28).
Let us look at the conformal dimensions. First note that because the value of $`\alpha `$ for a minimal model is never a multiple of $`\pi `$, one does not expect central sources in the UV limit but only right/movers or left movers. Indeed, in all the examples of sine-Gordon with central sources, the left-right symmetry (i.e. $`Z(\vartheta )=Z(\vartheta )`$ on the real axis) of the NLIE was crucial. This symmetry, however, only holds for $`\alpha =0`$ or $`\pi `$. As a consequence we have
$$S^0=0,S=S^++S^{},$$
and in addition $`𝒬_+(\mathrm{})=𝒬_{}(+\mathrm{})`$, which means that only *one-plateau systems* are allowed in the twisted case.
Using the formulas (4.26, 4.27) for the UV limit of the NLIE, one can see that introducing $`\alpha 0`$ is equivalent to shifting the quantum number $`n_\pm `$ to $`n_\pm +\frac{\alpha }{2\pi }`$. One has to be careful that since the value of the central charge is shifted from $`1`$ to the one of the minimal model, we have to take this shift into account when computing the conformal weight from the leading UV behaviour of the energy level (4.5). We put in (4.26) the value
$$\alpha =\frac{j\pi }{r},j=1,\mathrm{},r1,$$
as in (4.59), and the central charge
$$c=1\frac{6(rs)^2}{rs}$$
of the minimal model $`Vir(r,s)`$, as in section 4.7. By a computation similar to the one done in section 4.6 to obtain (4.27, 4.28, 4.30) one can obtain the following expression for conformal dimensions:
$$\begin{array}{c}\mathrm{\Delta }^\pm =\frac{(lsl_\pm ^{}r)^2(sr)^2}{4rs}+N_\pm \\ l=S\\ l_\pm ^{}=\pm \left(\delta \frac{j}{r}+2k_\pm 2k_W^\pm \right)2\left(S2S^\pm \right).\end{array}$$
(4.60)
The expression of $`N_\pm `$ is exactly the same that appears in (4.30). Moreover, from the formulae (4.26, 4.60) it is also clear that in general
$$2\left(\mathrm{\Delta }^+\mathrm{\Delta }^{}\right)=2S\frac{\alpha }{\pi }mod\mathrm{\hspace{0.17em}\hspace{0.17em}1}$$
because only one plateau systems can take place, and so we see that general charged states will have fractional Lorentz spin. Actually, it is known that charged states in the models $`Vir(r,s)+\mathrm{\Phi }_{(1,3)}`$ generally have fractional Lorentz spin .
Comparing the first line of (4.60) with (4.58) one observes that it can represent the conformal dimensions of an operator $`\mathrm{\Phi }_{(q,q^{})}`$ only if some conditions are verified. The most general one is
$$\left(lsl^{}r\right)^2=\left(qsq^{}r\right)^2+2rs\text{integer},1qr1,\mathrm{\hspace{0.17em}1}q^{}s1$$
(the $`l,l^{}`$ are not integer numbers in general). Observe that it is an arithmetic (Diofantine) equation in $`q,q^{}`$ and the unknown integer.
It is convenient to treat some specific cases, because, as for the sine-Gordon case, there is no proof about the general behaviour of $`l_\pm ^{}`$ and $`N_\pm `$.
In the unitary $`p=r,s=r+1`$ and neutral case $`S=0`$ the resulting conformal weights take the form
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=\frac{(2np+l)^21}{4p(p+1)}.$$
(4.61)
We see that this is the weight of the field $`\mathrm{\Phi }_{(l,l2n)}`$ in the minimal model $`Vir(p,p+1)`$, however, in order not to overflow the Kac table, the range of $`n`$ must be restricted as
$$1l2np.$$
For charged states, a similar calculation can be performed.
By an inspection of the formula (4.26), for any state with $`S=0`$, the property to have only one plateau implies that
$$\mathrm{\Delta }^+\mathrm{\Delta }^{}$$
is integer or half-integer. In fact, choosing the quantization rule and the parameter $`\delta `$ in an appropriate way, one can ensure that this difference is integer (see ). This means that the UV limit of any neutral state is either a field occurring in the ADE classification of modular invariants or (in case we choose $`\delta `$ so that $`\mathrm{\Delta }^+\mathrm{\Delta }^{}`$ is half-integer) it is a field from a fermionic version of the minimal model .
Until now it is not clear whether the weights actually stay inside the Kac table, for which in the unitary case one must require
$$1<l2k+4S^+<p.$$
Due to the fact that the configuration of sources in the UV may be very non-trivially related to the one in the IR, this condition is very hard to check in general, but no concrete examples that we calculated have ever violated this bound.
### 4.9 Concrete examples of excited states
#### 4.9.1 The $`Vir(2,2n+1)+\mathrm{\Phi }(1,3)`$ series
Let us start with examining the scaling Lee-Yang model $`Vir(2,5)+\mathrm{\Phi }_{(1,3)}`$. There is only one independent value of the twist which we choose as
$$\alpha =\frac{\pi }{2},$$
since we have a single ground state, and as a result there are no kinks in the spectrum. We fix the value of $`\alpha `$ as above, so we still have a freedom of choosing $`\delta `$. This can be done by matching to the UV dimensions: if for a certain state we choose the wrong value of $`\delta `$, we find a conformal dimension that is not present in the Kac table of the model.
The excited states are multi-particle states of the first breather of the corresponding unrestricted sine-Gordon model, which has
$$p=\frac{2}{3}.$$
Now one can calculate the state containing one particle at rest. We find the numerical data presented in table 4.7.
It turns out that as we decrease $`l`$, the self-conjugate root starts moving to the right. It does not remain in the middle like in the $`\alpha =0`$ case, which is to be expected since for nonzero $`\alpha `$ we have no left/right symmetry. However, the total momentum of the state still remains zero due to a contribution from the integral term in momentum equation (3.59).
One can see that once again we have the phenomenon noticed in the case of the first breather of sine-Gordon theory, namely the appearance of the special root and its two accompanying holes, so the iteration breaks down again around $`l=2.5`$. The (#) in the table 4.7 written after the NLIE result for $`l=2.6`$ means that due to the fact that the singularities corresponding to the new holes and the special root are just about to cross the contour and upset the iteration scheme, the NLIE result becomes less precise. We will use this notation on later occasions too. In any case, the agreement still looks quite convincing.
Let us now look at the UV spectrum of the model. We know that the Lee-Yang model contains only two primary fields, the identity $`𝕀`$ and the field $`\phi `$ with left/right conformal weights
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=\frac{1}{5}.$$
In fact, the ground state of the massive model corresponds to $`\phi `$ in the UV limit. One can compute the UV limit of the first particle from the NLIE too, taking into account the appearance of the special root and the holes. It turns out that the special root and one of the holes moves to the left together with the self-conjugate root, while the other hole moves to the right. The result is
$$\mathrm{\Delta }^+=\mathrm{\Delta }^{}=0,$$
i.e. the identity operator $`𝕀`$, which fits nicely with the TCS data (see also ).
Let us look now at moving breathers. If the self-conjugate root has Bethe quantum number $`I=1`$, the corresponding state will have momentum quantum number $`1`$, i.e.
$$P=\frac{2\pi }{R},$$
and in the UV $`\mathrm{\Delta }^+\mathrm{\Delta }^{}=1`$. One can note from the numerical data presented in table 4.8 that the special root does not appear here. The reason is that the self-conjugate root moves to the left and the real part of its position $`\vartheta `$ is given to leading order by
$$\mathrm{sinh}(\mathrm{}e\vartheta )\frac{2\pi I}{l_B}.$$
As a result, the contribution to the derivative of $`Z`$ from the $`l\mathrm{sinh}\vartheta `$ term remains finite when $`l\mathrm{\hspace{0.17em}0}`$. In the previous example of the particle at rest the left-moving nature of the self-conjugate root when $`I=0`$ does not prevent the occurrence of the breakdown in the iteration scheme: since its Bethe quantum number is zero, it does not move fast enough to the left in order to balance the negative contribution to derivative of $`Z`$ coming from the self-conjugate root source. At the moment we have no way of predicting analytically whether or not there will be specials in the UV limit: we just use the numerical results to establish the configuration for the evaluation of UV weights, supplemented with a study of the self-consistency of the solution of the plateau equation (4.16).
The UV dimensions for the moving breather turn out to correspond to the state $`L_1\phi `$.
One can similarly compute the UV dimensions for some other excited states. For example, the two-particle states with half-integer Bethe quantum numbers $`I_1>0,I_2<0`$ for the two self-conjugate roots are found to have
$$\mathrm{\Delta }^+=\frac{1}{5}+I_1+\frac{1}{2},\mathrm{\Delta }^{}=\frac{1}{5}I_2+\frac{1}{2},$$
in agreement with TCS data which show that they correspond in the UV to descendent states of $`\phi `$. The first such state with quantum numbers
$$I_1=\frac{1}{2},I_2=\frac{1}{2}$$
corresponds in the UV to $`L_1\overline{L}_1\phi `$ and is given numerically in table 4.9.
The lowest lying three-particle state of zero momentum, with Bethe quantum numbers $`(1,0,1)`$ corresponds to the left/right symmetric second descendent of the identity field, i.e. to the field $`T\overline{T}`$, where $`T`$ denotes the energy-momentum tensor. This is very interesting, since from experience with NLIE UV calculations one would naively expect this to be a first descendent (descendent numbers are usually linked to the sum of Bethe quantum numbers of left/right moving particles and this state is the lowest possible descendent of the identity $`𝕀`$). However, the field $`L_1\overline{L}_1𝕀`$ is well-known to be a null field in any conformal field theory.
The above correspondences are again confirmed by comparing to TCS (see the wonderful figures in ). In general, one can establish the rule that states with odd number of particles must be quantized by integers ($`\delta =1`$), while those containing even number of particles must be quantized by half-integers ($`\delta =0`$) in order to reproduce correctly the spectrum of the scaling Lee-Yang model.
We conducted similar studies for the models $`Vir(2,7)+\mathrm{\Phi }_{(1,3)}`$ and $`Vir(2,9)+\mathrm{\Phi }_{(1,3)}`$ and found similarly good agreement with TCS data. For the first one-particle state of the model $`Vir(2,7)+\mathrm{\Phi }_{(1,3)}`$ we also checked our results against the TBA data in the numerical tables of and found agreement with the TBA results.
Given the choice of $`\alpha `$ above, the correct rule of quantization in all of the models $`Vir(2,2n+1)+\mathrm{\Phi }_{(1,3)}`$ is
$$\delta =M_{sc}mod\mathrm{\hspace{0.17em}\hspace{0.17em}2},$$
where $`M_{sc}`$ is the number of self-conjugate roots in the source corresponding to the state. This is exactly the same rule as the one established for pure sine-Gordon theory in . In the presence of the twist, such a rule of course has meaning only together with a definite convention for the choice of $`\alpha `$.
#### 4.9.2 One-breather states in the $`Vir(3,7)`$ case
It is interesting to note that in the case of $`Vir(3,n)`$ models, all the neutral states must come in two copies, since they can be built on top of either of the two ground states. We take the example of the $`Vir(3,7)`$ model and the states corresponding to a breather at rest. We have
$$p=\frac{3}{4}$$
but now there are two inequivalent values for the twist
$$\alpha =\frac{\pi }{3},\frac{2\pi }{3}.$$
When $`\alpha =0`$, we can calculate the critical value of $`l`$ to be $`l_{critical}=5.23`$ using (4.54). In this case, the twist helps a bit, because it makes the self-conjugate root a left mover; it is intuitively clear that the bigger the twist, the more it lowers the eventual value of $`l_{crtical}`$, which is in accord with the numerical results of table 4.10. From the TCS data one can identify that breather #1 is really the one-particle state in the sector of ground state #1 ($`\alpha =\pi /3`$), while breather #2 is in the sector over ground state #2 ($`\alpha =2\pi /3`$).
A direct calculation of the conformal weights gives the following results:$`\mathrm{\Delta }^+=\mathrm{\Delta }^{}=\frac{3}{28}`$for breather #1 and $`\mathrm{\Delta }^+=\mathrm{\Delta }^{}=0`$for breather #2, which are in complete agreement with the TCS data. We also checked the two different states containing two breathers with Bethe quantum numbers$`I_1=\frac{1}{2},I_2=\frac{1}{2},`$and found an equally excellent numerical agreement with TCS. Just like in the case of sine-Gordon and scaling Lee-Yang model, for these states one can continue the iteration of the NLIE down to any small value of $`l`$, although at the expense of a growing number of necessary iterations to achieve the prescribed precision.
### 4.10 Conclusions
In this thesis it is studied how the nonlinear integral equation deduced from the light cone lattice model of describes the excited states of the sine-Gordon/massive Thirring theory. The most important results are summarized as follows:
1. A derivation of the fundamental NLIE is presented from the light cone lattice which correctly takes into account the behaviour of the multivalued complex logarithm function.
2. By examining the infrared limits of the equation it has been shown that (1) it leads to the correct two-particle S-matrices for both scattering states and bounded states; (2) it is in agreement with the predictions of the TCS method if one chooses the correct quantization conditions for the source terms.
3. By computing the UV conformal weights from the NLIE we have shown that it is consistent with the UV spectrum of sG/mTh theory only if we choose the parameter $`\delta `$ (i.e. the quantization rule) as indicated in (4.31).
4. The predictions of the NLIE have been verified by comparing them to results coming from the TCS approach (for sG/mTh).
5. The framework required to deal with minimal models perturbed by $`\mathrm{\Phi }_{(1,3)}`$ is built up. Many examples and numerical/analytical checks are given for the ground state in the unitary and non unitary cases.
6. All the conformal dimensions can be reproduced (Kac table) with the convenient choice of the twist (4.59). The particular relation suggested in is not enough to describe the whole spectrum.
7. The IR computations given in section 4.6.4 reproduce correctly the S-matrix of perturbed minimal models in the attractive regime (at list in the cases involving only solitons), because the attractive second determination drops the twist.
8. Numerical calculations of concrete examples give a strong evidence for the correctness of the energy levels derived from the twisted NLIE for excited states.
The understanding of “(twisted) sine-Gordon NLIE” and of the finite size behaviour of the continuum theory defined from the NLIE (3.54) is not quite complete. Indeed some open questions have not an answer, until now and also further interesting developments can take place:
1. Is the set of scaling functions provided by the NLIE complete i.e. can we find to every sG/mTh or minimal model state a solution of the NLIE describing its finite volume behaviour? This can be called a “counting problem”. The main difficulty is that the structure of the solutions is highly dependent on the value of the coupling constant – to see that it is enough to consider e.g. the appearance of special sources.
2. The multi-kink states characteristic of the perturbed minimal models (except for the series $`Vir(2,2n+1)`$) have been omitted in the analysis so far performed. Although the general treatment of the IR spectrum in section 4.6.4 is valid for those states too, a detailed description is far more complicated than for states which contain only breathers and is left open to further studies.
3. There is also an unresolved technical difficulty, namely that the source configuration of the NLIE may change as we vary the volume parameter $`l`$. Typically what happens is that while the counting function $`Z`$ is monotonic on the real axis for large volume, this may change as we lower the value of $`l`$ and so-called special sources (and accompanying holes) may appear. We do not as yet have any consistent and tractable numerical iteration scheme to handle this situation, although the analytic UV calculations and intuitive arguments show that the appearance of these terms in the NLIE is consistent with all expectations coming from the known properties of perturbed CFT. In addition, in the range of $`l`$ where we can iterate the NLIE without difficulty, our numerical results show perfect agreement with TCS. We want to emphasize that these transitions are not physical: the counting function $`Z`$ and the energy of the state are expected to vary analytically with the volume. It is just their description by the NLIE with requires a modification of source terms. As it was pointed out also in , the whole issue is related to the choice of the branch of the logarithmic term in the NLIE (3.54).
4. The problem pointed out at the previous point is very similar to the behaviour of singularities encountered in the study of the analytic continuation of the TBA equation and we can hope that establishing a closer link between the two approaches can help to clarify the situation. From the form of the source terms in the NLIE it seems likely that the excited state equations can be obtained by an analytic continuation procedure analogous to the one used in TBA to obtain the excited state TBA equations. Certain features of the arrangement of the complex roots in the attractive regime and their behaviour at breather thresholds also point into this direction. This an interesting question to investigate because it can shed light on the organization of the space of states and can lead closer to solving the counting problem described above.
5. Even if on the lattice the Bethe vectors given in (2.37) completely describe the Hilbert space of the XXZ chain, the form of eigenvectors for continuum energy and momentum is completely unknown. This fact reflects the unresolved question of the form of energy and momentum eigenvectors for the minkowskian sine-Gordon theory (the Faddeev-Zamolodchikov algebra is only a phenomenological picture). Also the determination of correlation functions is in an early stage, but in recent publications explicit expressions for the field $`e^{ia\phi }`$ and its descendent are given.
6. The extension of light-cone approach and NLIE to other QFT models is in progress. An interesting extension is the description of finite volume spectrum of Ziber-Mikhailov-Shabat model (imaginary Bullogh-Dodd). A first suggestion in this direction, even if is obtained in a completely different approach, recently appeared in the , for the ground state.
The second and fourth points are important because their investigation may lead closer to understanding the relation between the TBA and the NLIE approaches. It is quite likely that establishing a connection between the two methods would facilitate the development of both and may point to some common underlying structure.
As a final remark, I want to speak about an application of the NLIE to a chemical compound, i.e. “copper benzoate” ($`\text{Cu}(\text{C}_6\text{D}_5\text{COO})_23\text{D}_2\text{O}`$). Its specific heat has been computed in using the vacuum untwisted ($`\alpha =0`$) NLIE, where the finite size $`L`$ is the inverse temperature $`L=(k_BT)^1`$. This corresponds to do equilibrium thermodynamics of sine-Gordon model.
A curiosity: in that contest the NLIE is called Thermal Bethe Ansatz, emphasizing the well known relation between statistical mechanics and QFT. Exactly in a similar contest (Heisenberg ferromagnet model) was born the Bethe Ansatz .
## Appendix A Fourier transformation: some conventions
Given a bounded continuous function $`f(x)`$ defined on the whole real axis, its Fourier transform is defined as (all the integrals must be taken on the whole real axis):
$$\stackrel{~}{f}(k)=𝑑xe^{ikx{\displaystyle \frac{1}{p+1}}}f(x)$$
(A.1)
and the inverse Fourier transform is
$$f(x)=\frac{1}{p+1}\frac{dk}{2\pi }e^{ikx{\displaystyle \frac{1}{p+1}}}\stackrel{~}{f}(k).$$
(A.2)
The convolution of two functions, defined as follows, can be expressed in terms of the Fourier transforms of the two functions:
$$\left(fg\right)_\theta =𝑑xf(\theta x)g(x)=\frac{1}{p+1}\frac{dk}{2\pi }e^{ik\theta {\displaystyle \frac{1}{p+1}}}\stackrel{~}{f}(k)\stackrel{~}{g}(k)$$
(A.3)
The following expression holds between the Fourier transforms of a function and its derivative:
$$\stackrel{~}{f}(k)=\frac{\stackrel{~}{f}^{}(k)}{ik}(p+1)$$
## Appendix B The function $`\varphi (\lambda ,\nu )`$
In this appendix all the important properties of the function $`\varphi `$ will be clarified. Let
$$\varphi (\vartheta ,\nu )=i\mathrm{log}\frac{\mathrm{sinh}{\displaystyle \frac{1}{p+1}}(i\pi \nu +\vartheta )}{\mathrm{sinh}{\displaystyle \frac{1}{p+1}}(i\pi \nu \vartheta )},\varphi (\vartheta ,\nu )=\varphi (\vartheta ,\nu )\vartheta $$
(B.1)
and continuity is required around the real axis. The interest is for $`\nu =1/2,\mathrm{\hspace{0.33em}1}`$ and $`p>0.`$ The requirement of oddity implies, for real values of $`\vartheta ,`$ that $`\varphi (0,\nu )=0.`$ Oddity and continuity can be implemented only if the fundamental determination (FD) of the logarithm is assumed, in a strip containing the real axis:
$$\pi <\mathrm{}m(\mathrm{log}_{FD}w)\pi .$$
This choice implies that $`\varphi `$ is real on the real axis. The argument of the log has poles and zeros. They are essential singularities for the $`\varphi `$ function. Their position is given by:
$$\begin{array}{c}\mathrm{}e\vartheta =0\\ \mathrm{}m\vartheta =\pm \pi \left(k(p+1)\nu \right),k.\end{array}$$
(B.2)
The fundamental strip around the real axis (FD) is bounded by the first encountered singularity:
$$\vartheta \times ]\mathrm{min}\{\nu \pi ,\pi \left(p+1\nu \right)\},\mathrm{min}\{\nu \pi ,\pi \left(p+1\nu \right)\}[.$$
(B.3)
To extend the definition of (B.1) to the whole plane it is necessary to give a prescription for the position of the cuts (because of the logarithm, in a small closed trip around one of the singularities (B.2), the value of $`\varphi `$ changes by $`2\pi `$). This is done as in figure B.1.
The derivative of $`\varphi `$ is the single value meromorphic function:
$$\varphi ^{}(\vartheta ,\nu )=\frac{1}{p+1}\frac{2\mathrm{sin}{\displaystyle \frac{2\nu \pi }{p+1}}}{\mathrm{cosh}{\displaystyle \frac{2\vartheta }{p+1}}\mathrm{cos}{\displaystyle \frac{2\nu \pi }{p+1}}}.$$
(B.4)
This shows that, on the real axis, $`\varphi `$ is monotonically increasing if $`0<{\displaystyle \frac{\nu \pi }{p+1}}<\pi /2`$ and decreasing if $`\pi /2<{\displaystyle \frac{\nu \pi }{p+1}}<\pi .`$ The asymptotic values, for $`\vartheta `$ in the fundamental strip, are:
$$\underset{\vartheta \pm \mathrm{}}{lim}\varphi (\vartheta ,\nu )=\pm \pi (1\frac{2\nu }{p+1}).$$
(B.5)
Out of the fundamental strip, the prescription indicated in figure B.1 must be used.
The Fourier transformation of $`\varphi ^{}`$ can be obtained from the definition (A.1), using the theorem of residues:
$$\stackrel{~}{\varphi ^{}}(k,\nu )=2\pi \frac{\mathrm{sinh}{\displaystyle \frac{\pi }{2}}\left({\displaystyle \frac{p+12\nu }{p+1}}\right)k}{\mathrm{sinh}{\displaystyle \frac{\pi }{2}}k}$$
(B.6)
This formula is the Fourier transform of $`\varphi ^{}(\vartheta ,\nu )`$ only for $`\vartheta `$ in the fundamental strip (B.3), because out of this strip new singularities of (B.4) appear in the computation of Fourier transform.
## Appendix C A lemma for UV computations
In the following lemma has been proved. Assume that $`f(x)`$ satisfies the non-linear integral equation
$$i\mathrm{log}f(x)=\phi (x)+2\mathrm{}m\frac{dy}{i}G(xyiϵ)\mathrm{log}(1+f(x+iϵ))$$
where $`\phi (x)`$ is real on the real axis and $`G(x)=G(x)`$ is real too, with bounded integral and peaked around the origin. From this equation follows that $`f(x)`$ has unit modulus on the real axis. To avoid crossing of the branch cut of the logarithm, assume that
$$\text{if }f(x+iϵ)\text{then}f(x+iϵ)>1.$$
(C.1)
Then the following expression holds:
$$\begin{array}{c}2\mathrm{}m\frac{dx}{i}\phi ^{}(x+iϵ)\mathrm{log}(1+f(x+iϵ))=\\ \\ =2\mathrm{}e_\mathrm{\Gamma }\frac{du}{u}\mathrm{log}(1+u)\frac{1}{2}𝑑xG(x)\left(F^2(+\mathrm{})F^2(\mathrm{})\right)\end{array}$$
(C.2)
where the curve $`\mathrm{\Gamma }`$ is any path in the complex plane that goes from $`f_{}=f(\mathrm{}+iϵ)`$ to $`f_+=f(+\mathrm{}+iϵ)`$ avoiding the branch cut, i.e. respecting the condition (C.1) and
$$F(x)=2\mathrm{}m\mathrm{log}(1+f(x+iϵ)).$$
In the 6-vertex case all the terms can be made manifest by the identification of $`G`$ with that defined in (3.21), $`f(x)=(1)^\delta e^{iZ_\pm (x)}`$ and $`\phi (x)`$ with (4.25).
Consider the kink “+”. The integral on the path $`\mathrm{\Gamma }`$ goes from the points
$$\begin{array}{c}f_{}=f(\mathrm{}+iϵ)=(1)^\delta e^{iZ_+(\mathrm{})}=e^{i𝒬_+(\mathrm{})}\\ f_+=f(+\mathrm{}+iϵ)=(1)^\delta e^{iZ_+(+\mathrm{}+iϵ)}=0\end{array}$$
(the first one can be computed with (4.17)). Observe that $`|f_{}|=1`$ then the path $`\mathrm{\Gamma }`$ can be composed by an unit radius arc from $`f_{}`$ to the point $`1`$ and a segment from $`1`$ to $`0`$. On the arc the integration variable has unit modulus then it is convenient to do the change of variables $`u=e^{i\alpha }`$. This simple computation gives
$$2\mathrm{}e_{arc}\frac{du}{u}\mathrm{log}(1+u)=𝒬_+^2(\mathrm{}).$$
The integration on the segment gives a well known dilogarithmic expression:
$$2\mathrm{}e_1^0\frac{du}{u}\mathrm{log}(1+u)=2_0^1\frac{du}{u}\mathrm{log}(1+u)=\frac{\pi ^2}{6}.$$
The quantities $`F(\pm \mathrm{})`$, from their definition, give:
$$\begin{array}{c}F(+\mathrm{})=0\\ F(\mathrm{})=2\mathrm{}m\mathrm{log}(1+(1)^\delta e^{iZ_+(\mathrm{})})=𝒬_+(\mathrm{}).\end{array}$$
Remembering now that the integral of $`G`$ appearing in (C.2) is expressed in (3.33) yields (the computation for the “–” kink is completely analogous):
$$\begin{array}{c}2\mathrm{}m\frac{dx}{i}\phi _\pm ^{}(x+i\eta )\mathrm{log}\left(1+(1)^\delta e^{iZ_\pm (x+i\eta )}\right)=\\ \\ =\pm \frac{\pi ^2}{6}\frac{𝒬_\pm ^2(\mathrm{})}{4}\frac{p+1}{p}\end{array}$$
(C.3) |
warning/0001/astro-ph0001289.html | ar5iv | text | # Optical Identification of the ASCA Large Sky Survey
## 1 Introduction
Since the discovery of the cosmic X-ray background (CXB) by Giacconi et al. (1962) in the 2–6 keV band, many efforts have been made to understand the origin of the CXB. Recently in ROSAT deep surveys, 70 – 80% of the CXB in the 0.5–2 keV band has been resolved into discrete sources at a flux limit of $`1\times 10^{15}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$(Hasinger et al. (1998)). Within the flux level of the deep survey, broad-line AGNs are the dominant population and the main contributor to the CXB in the 0.5–2 keV band. On the other hand, in the harder 2–10 keV band, only $`3`$% of the CXB was resolved into discrete sources (Piccinotti et al. (1982)) before ASCA surveys. Broad-line AGNs have X-ray power-law spectra with a photon index of $`\mathrm{\Gamma }=1.7`$ in the 2–10 keV band (Turner & Pounds (1989)) which are significantly softer than that of the CXB in that band ($`\mathrm{\Gamma }=1.41.5`$; Gendreau et al. (1995); Ishisaki et al. (1999)), thus there must be objects which have harder X-ray spectra than nearby broad-line AGNs and contribute significantly to the CXB in the hard band.
Resolving the hard X-ray sky is a direct way to reveal the nature of X-ray sources in the hard band. ASCA GIS observations of three ROSAT Deep PSPC Fields down to the $`5\times 10^{14}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$have been done so far (Georgantopoulos et al. (1997); Boyle et al. 1998a ). However, because of the large positional uncertainties of the GIS-selected X-ray sources, a large fraction of hard X-ray selected sources is still unidentified; the nature of hard X-ray sources and the difference from the ROSAT selected objects are still unclear.
Based on the unified scheme of AGNs, absorbed AGNs are proposed as candidates for the hard X-ray sources (Comastri et al. (1995); Madau, Ghisellini, & Fabian 1994). Because of absorption of soft X-ray photons by obscuring material, they have harder X-ray spectra than type 1 AGNs. To reproduce the CXB spectrum, it is argued that there are around three times more absorbed AGNs than non-absorbed AGNs in the universe and absorbed AGNs dominate the CXB above 2 keV. In consequence of the assumption, the existence of absorbed narrow-line QSO (so called type 2 QSO) is expected. From optical follow-up observations of X-ray surveys in the soft and hard bands, several luminous narrow-line AGNs have been found at intermediate to high redshift universe (e.g., Stocke et al. (1982); Almaini et al. (1995); Ohta et al. (1996); Boyle et al. 1998b ; Barcons et al. (1998); Schmidt et al. (1998)). On the other hand, existence of red broad-line QSOs whose red color suggests absorption to its nucleus was reported from a radio survey (Webster et al. (1995)). Such a population was also identified in ROSAT surveys (Kim & Elvis (1999)) and Beppo-SAX surveys (Fiore et al. (1999)), but X-ray spectra, number densities and contributions to the CXB of both of these narrow-line and red broad-line QSOs are not clear.
To reveal the nature of X-ray sources in the hard band, studies on a well-defined sample are important. We are now conducting an unbiased large and deep survey with ASCA in the region near the north Galactic pole, i.e., ASCA Large Sky Survey (hereafter LSS; Inoue et al. (1996); Ueda (1996); Ueda et al. (1998); Ueda et al. 1999a (Paper I)). The flux limit of the LSS ($``$ 1 $`\times 10^{13}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the 2–10 keV band) is 100 times deeper than the HEAO1 A2 survey, which was the deepest systematic survey in the hard band (Piccinotti et al. (1982)) before ASCA. We have surveyed 7.0 deg<sup>2</sup> and 5.4 deg<sup>2</sup> with the GIS and the SIS detectors, respectively. Combining the data from the GIS and the SIS, we detected 44 sources in the 2–10 keV band with the following criteria: 1) the significance of summed count rate of the GIS and the SIS should exceed 4.5, and 2) the significance of either the GIS or the SIS should also exceed 3.5 (Paper I ). They correspond to 20 – 30% of the CXB in this band. The advantages of the sample are 1) the survey area as a function of limiting count rates was determined well by simulations, 2) accurate X-ray positions were determined by using the SIS whose resolution is higher than the GIS and by correcting the temperature-dependent misalignment between the focal plane detectors and the attitude sensors, and 3) the X-ray spectrum of each source was determined by fitting power-law model with the Galactic absorption to the GIS and the SIS data simultaneously in the 0.7–10 keV band (Paper I ). Determining their X-ray spectra is important not only to know their X-ray properties but also to know the flux limit of the survey, because the limit varies with the X-ray spectrum of each source. The average of the apparent photon index of the 36 X-ray sources detected in the flux range between 0.8 $`\times 10^{13}`$ and 4 $`\times 10^{13}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ is $`\mathrm{\Gamma }=1.49\pm 0.1`$ (Paper I ), which is significantly harder than the spectra of X-ray sources detected in shallower surveys in the 2–10 keV band and close to that of the CXB. Identification of these sources is clearly important to understand the nature of hard X-ray sources and the origin of the CXB.
In this paper, we report results of optical identifications of X-ray sources detected in the hard band with the SIS in the ASCA LSS. The sample definition and selections of optical candidates are discussed in Section 2, results of optical spectroscopy for the selected candidates and reliability of the identifications are presented in Section 3, and optical and X-ray spectroscopic properties of identified objects are described in Section 4. The contribution of the population to the CXB and the redshift distribution of the AGNs are discussed in Section 5 and 6, respectively. In Section 7, we present multi-wavelength properties of the identified AGNs. Throughout this paper, we use $`q_0=0.5`$ and $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup>. We call each X-ray source with the exact name and the identification number, like AX J132032+3326(227), for convenience.
## 2 Observations
### 2.1 The X-ray Survey Observations and the Sample Definition
To minimize the effect of the galactic absorption and contamination from bright X-ray sources, the survey area was defined as a continuous region of $``$ 5 degree<sup>2</sup> near the north Galactic pole, centered at $`\alpha `$=13<sup>h</sup>14<sup>m</sup>, $`\delta `$=3130 (J2000). Seventy-six pointings of the survey-observations were done from Dec. 1993 to Jul. 1995. These exposures were designed to evenly cover the whole survey region with a 20 ksec effective exposure of the SIS. The source detections were done with the SIS data in the 0.7–7 keV, 0.7–2 keV, and 2–7 keV bands and the GIS data in the 0.7–7 keV, 0.7–2 keV, and 2–10 keV bands. For details of the survey observation, source extraction, and spectral fitting, including the survey region and the position of the sources on the sky, see Paper I .
In this paper, we concentrate on the 34 X-ray sources detected with the SIS in the 2–7 keV band above 3.5$`\sigma `$ (hereafter, the SIS 2–7 keV 3.5$`\sigma `$ sample). Table 1 shows the survey area as a function of limiting count rates for the sample. The typical and the deepest limiting count rates of the sample are 2 counts ksec<sup>-1</sup> and 1.2 counts ksec<sup>-1</sup>, respectively. They correspond to $`1.8\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$and $`1.1\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$in the 2–10 keV band for an X-ray source with a power-law spectrum with a photon index of 1.7. All X-ray sources in the SIS 2–7 keV 3.5$`\sigma `$ sample have the significance level larger than 4.5 in summed count rate of the GIS and the SIS in both of the 0.7–7 keV and 2–7 keV bands and the positional uncertainties of such sources were estimated to be $`0.^{}6`$ in radius with the 90% confidence level (Paper I ). The number of spurious sources was estimated to be at most a few percent (Paper I ), thus less than 1 spurious source was expected to be in the SIS 2–7 keV 3.5$`\sigma `$ sample.
### 2.2 Optical Imaging Observations and Selections of Optical Counterpart Candidates
AGNs are the most plausible optical counterparts for the majority of the X-ray sources, thus at least we have to reach the optical flux limit which is converted from the deepest X-ray flux limit based on X-ray-to-optical flux ratio of AGN. If we assume power-law spectra with an X-ray photon index of 1.7 and an optical energy index of $`0.5`$ together with the X-ray-to-optical flux ratio of AGNs identified in the ROSAT Deep Survey (Schmidt et al. (1998)), the expected optical magnitude for the optical counterpart of an X-ray source with $`1.1\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$in the 2–10 keV band is $`R=1621`$ mag.
Candidates of optical counterparts were mainly selected from the APM catalog which is obtained from scans of glass copies of the Palomar Observatory Sky Survey plate (McMahon et al. (1992)). However, the limiting magnitude of the catalog is $`R20`$ mag and is not deep enough to pick up optical counterparts for faint X-ray sources. To complement the depth of the data, we made imaging observations of the LSS region at the KISO 1.05m Schmidt telescope in March and April 1994. In these observations, we used the mosaic CCD camera (Sekiguchi et al. (1992)) which is made up of 15 1024$`\times `$1024 CCD chips and covers 2$`\times `$5 degree<sup>2</sup> with 15 shots. The spatial resolution in the setup was $`0.^{\prime \prime }75`$ pix<sup>-1</sup>. Images were taken in the $`R`$ band with an exposure time of 20 minutes. The reduction of the data was done by the usual method for optical imaging data. The weather condition during the observation was neither stable nor photometric and the limiting magnitude changed from field to field. The typical seeing was $`4^{\prime \prime }`$ (FWHM) and the typical limiting magnitude was $`R21`$ mag, about one magnitude deeper than the APM data.
There are several optical objects within each error circle of X-ray source above $`R=21`$ mag. Two X-ray sources (AX J132032+3326(227) and AX J130748+2925(002)) show clear excesses of galaxies in and around their error circles (see Figure Optical Identification of the ASCA Large Sky Survey). These two objects have been cataloged as candidate clusters of galaxies, Abell 1714 (Abell, Corwin, & Olowin 1989) and Zwcl 1305.4+2941 (Zwicky, Herzog, & Wild 1961), respectively. Zwcl 1305.4+2941 was also detected in the Einstein Medium Sensitivity Survey and its redshift was determined to be z=0.241 (Stocke et al. (1991)). We took deeper and better resolution images with a Tektronix 2048$`\times `$ 2048 CCD on the University of Hawaii $`88^{\prime \prime }`$ telescope on 1995 March and 1996 April in the $`R`$ and $`I`$ band and confirmed the excess of galaxies. Thus, we identified the two sources with clusters of galaxies. For other sources, active galaxies are the most plausible counterparts. From optical objects within $`0.^{}8`$, which is slightly larger than the estimated 90% confidence error radius ($`0.^{}6`$), of each X-ray source, we selected targets for spectroscopy, using 1) deeper X-ray follow-up data in the soft and hard band, 2) radio emission, and 3) blueness of UV or optical color which indicate existence of an activity in a object.
#### 2.2.1 Deep Follow-up Observations in X-ray band
To pinpoint optical counterparts of X-ray sources, 20 ksec follow-up observations were made with the ROSAT HRI in 2 fields in December 1997. We selected two fields centered at $`\alpha `$=13<sup>h</sup>14$`{}_{}{}^{\mathrm{m}}36_{}^{\mathrm{s}}`$, $`\delta `$=320112<sup>′′</sup> and $`\alpha `$=13<sup>h</sup>12$`{}_{}{}^{\mathrm{m}}29_{}^{\mathrm{s}}`$, $`\delta `$=311312<sup>′′</sup> (J2000) with a radius of $`19^{}`$, to cover as many ASCA sources as possible. We summarize the data reduction and the source detection in Appendix A. All of the X-ray sources in the SIS 2–7 keV 3.5$`\sigma `$ sample (AX J131521+3159(136), AX J131407+3158(127), and AX J131327+3155(121) in the former field, and AX J131249+3112(096), AX J131128+3105(080), and AX J131321+3119(103) in the latter field) except one (AX J131345+3118(104)) in these fields were detected by the HRI and their optical counterparts were pinpointed thanks to the $`5^{\prime \prime }`$ positional accuracy of the HRI. On the position of the missed source (AX J131345+3118(104)), there is an X-ray peak in the HRI image, though the significance level is slightly lower than the detection limit. Since there is also an optical object at the position, the X-ray source is pinpointed. It should be noted that there is no hard X-ray source which has an X-ray spectrum with an apparent photon index smaller than 1.0 in these HRI observed fields.
Deep 40 $``$ 100 ksec pointing observations for 4 sources, which have hard X-ray spectra, (AX J131501+3141(119), AX J131551+3237(171), AX J131210+3048(072), and AX J130926+2952(016)) were made by ASCA and more precise positions and X-ray spectra were obtained (Sakano et al. (1998); Sakano et al. (1999); Ueda et al. 1999b ). The error circles for these sources were estimated to be less than $`0.^{}6`$ at a 90% confidence level.
By the ROSAT PSPC, 15 and 13 ASCA LSS sources were detected in the ROSAT PSPC All-Sky Survey (Voges et al. (1999)) and pointing observations (Voges, private communication), respectively. 9 of them were detected in both of the observations, thus the X-ray positions of 19 ASCA sources were determined accurately. In details of the cross-identification, see Appendix B.
#### 2.2.2 Cross-correlation with FIRST radio source catalog
In the course of the optical counterpart selection, we examined the distribution of FIRST radio sources around LSS X-ray sources to evaluate the cross-correlation between radio and X-ray sources. The FIRST survey is a radio source survey conducted with the Very Large Array in the 1.4 GHz band with a $`5\sigma `$ limiting flux of 1 mJy (Becker, White, & Helfand (1995)). There are 17 radio sources within $`0.^{}5`$ from the centers of the X-ray sources; by contrast no radio source exists between $`0.^{}5`$ and $`1^{}`$ from them. Based on the surface number density of detected radio sources in the LSS field, the contamination of a radio source which is not an counterpart of an X-ray source is expected to be less than 1 source for the whole 34 X-ray sources within $`0.^{}5`$. Thus, the 17 radio sources are likely radio-counterparts of X-ray sources. In details of the cross-identification, see Appendix B. Three X-ray sources have two to four radio sources within $`0.^{}5`$. They are thought to be radio-loud objects with radio lobes or clusters of galaxies. One of the three X-ray sources is already identified with a cluster of galaxies (AX J132032+3326(227)). In summary, 35% (12/34) of the SIS 3.5$`\sigma `$ sample are detected in the FIRST survey and an optical counterpart is pinpointed thanks to the good positional accuracy ($`1^{\prime \prime }`$) of the FIRST catalog, except for AX J131639+3149(137). For the X-ray source, there is no optical object within an error circle of the FIRST radio source (see notes in Section 3.1). In the error circle of AX J131529+3117(110), there is a radio peak with a flux of 0.63 mJy which is slightly lower than the detection limit of the FIRST. Because there is an optical object at the position, we identified the optical object with the counterpart of the radio peak and the X-ray source.
The high fraction of radio-detected X-ray sources in comparison with other results (e.g., 10% in De Ruiter et al. (1997)) is due to a match between the flux limit of the VLA FIRST 1.4 GHz survey and that of the LSS; the ratio of the flux limits is one order of magnitude larger than the typical radio-to-X-ray flux ratio of radio-quiet broad-line AGN (Elvis et al. (1994)) but within the scatter of the ratio of the radio-quiet population. Thus, not only radio-loud AGNs but also some radio-quiet AGNs are included in the LSS-FIRST sample (see Section 7.2). It is worth noting that radio emission is transparent against obscuring material, thus obscured AGNs are detected in the radio wavelength as well as in the hard X-ray band, in an unbiased way.
#### 2.2.3 UV and Optical Color
Based on the results of the optical identification in soft X-ray surveys, a broad-line AGN is a plausible optical counterpart for a part of X-ray sources detected in the hard band and they can be selected by a blueness of their UV or optical color if they are at redshift smaller than 3. There are two UV-excess surveys which cover the southern half ($`\delta <32^{}10^{}`$) of the LSS. These are Usher (1981) and Moreau & Reboul (1995). Objects with an $`UV`$ color bluer than 0.0 down to $`B=20`$ mag and 0.1 down to $`V=20`$ mag are cataloged in the former and the latter, respectively. 8 out of 20 X-ray sources in the area have one or two UV-excess object(s) within $`0.^{}8`$ in total. 6 are also detected in the X-ray follow-ups with ROSAT HRI or in the FIRST survey. Considering the number density of UV-excess objects in Moreau & Reboul (1995), the expected number of contamination within $`0.^{}8`$ from the center of the 20 X-ray sources is calculated to be 0.5, thus all 8 correlated objects can be considered as plausible optical counterparts of X-ray sources. Such a low probability of contamination indicates that there is no serious bias to selecting an UV excess object which is not a true optical counterpart of an X-ray source. Additionally, two galaxies in the $`0.^{}8`$ radii of AX J131805+3349(233) are cataloged as UV-excess galaxies in the KUG catalog (Takase & Miyauchi-Isobe (1986)).
To select candidates of broad-line AGNs with fainter magnitude and to cover the whole LSS area, we also picked up objects with $`BR<1.0`$ mag as a broad-line AGN candidates in the whole area, assuming a typical color of broad-line QSOs ($`BR=0.52`$ mag; Richstone & Schmidt (1980)) with a dispersion of distribution ($`\mathrm{\Delta }BR=0.2`$ mag) and an error of photometry ($`\mathrm{\Delta }BR=0.3`$ mag). The APM $`OE`$ color was converted to the $`BR`$ color with $`(OE)_{\mathrm{APM}}=1.135\times (BR)`$ (Evans (1989)). 19 X-ray sources have one or two optical object(s) which meet the above criteria within $`0.^{}8`$. 5 of them are also detected in either the X-ray follow-ups with ROSAT HRI or in the FIRST survey. Within $`0.^{}6`$ and between $`0.^{}6`$ and $`0.^{}8`$ from the 34 X-ray sources, there are 19 and 6 blue objects, respectively. Based on the number density of objects with $`BR`$ color bluer than 1.0 in the APM catalog, the expected number of contamination of blue objects which are not X-ray source is estimated to be 6 and 5 in the inner circles and outer annuli, respectively. Thus, two thirds of the blue objects in the inner circles are expected to be true optical counterpart of X-ray sources, on the other hand, most of the blue objects in the outer annuli are thought to be contamination of non-X-ray objects.
Almost all (5 out of 6) the hard X-ray sources which have a photon index smaller than 1.0 have no such blue object within their error circles; therefore an optical counterpart different from a broad-line AGN is suggested. This tendency is the same as the previous results reported in Akiyama et al. (1998b), but the new result is more significant than the previous one, because we concentrate on brighter X-ray sources whose positions and X-ray photon indices were determined well and more accurate X-ray positions were obtained for each X-ray source by correcting for the temperature effect in the satellite acquisition as mentioned in Section 2.1.
In summary, except for 2 clusters of galaxies, 24 sources were pinpointed by X-ray follow-ups with the ROSAT HRI or PSPC or radio detections, 5 more sources have only one object which is selected by its blue optical color, and the remaining 3 sources have no optical counterpart which meets the above criteria. For these 3 sources, we picked up the brightest optical object within each error circle as a primary target. The list of the selected objects is shown in Table Optical Identification of the ASCA Large Sky Survey and, in finding charts of Figure Optical Identification of the ASCA Large Sky Survey, they are labeled with “A”, “B”, and “C” in order of distance from the X-ray centroid. Additionally, we listed some objects which were not selected by the above methods but observed in the optical spectroscopy. These objects were indicated with a label “Z”.
### 2.3 Optical Spectroscopy
We made spectroscopic observations for the optical-counterpart candidates which are selected based on the above criteria with the highest priority. Spectroscopic observations were made with the University of Hawaii $`88^{\prime \prime }`$ telescope on March 1998, except for the hardest source (AX J131501+3141(119)), a candidate of galactic star (AX J131850+3326(219)), and the two clusters of galaxies. We used the Wide Field Grism Spectrograph with a grating of 420 rulings mm<sup>-1</sup> and the blaze wavelength of 6400Å. The spatial resolution was $`0.^{\prime \prime }354`$ pixel<sup>-1</sup> and the seeing during the observation was $`0.^{\prime \prime }81^{\prime \prime }`$. A slit width of $`1.^{\prime \prime }2`$ was used. The spectral coverage ranges from 4000Å to 9000Å and the spectral resolution, which is measured by the HgAr lines in comparison frames and night-sky lines in object frames, is 12Å (FWHM). We also took imaging data of each X-ray source with a better angular resolution without filter for finding charts.
We have also done spectroscopic observations with the 3.5m telescope at Calar Alto observatory for a candidate galactic star (AX J131850+3326(219)) and three objects (AX J131831+3341(228), AX J130840+2955(014), and AX J131639+3149(137)) which were already observed in the previous run. We used the MOSCA instrument in a single-slit mode with a g250 grating which has 250 rulings mm<sup>-1</sup> and a blaze wavelength of 5700Å. The spectral coverage ranges from 4000Å to 8000Å. In the configuration, the sampling was 5.95 Å pixel<sup>-1</sup>. A slit width of $`1.^{\prime \prime }5`$, which was the same as the FWHM of the seeing, was used. Thus, the spectral resolution was 22Å (FWHM), which was estimated by measuring widths of night sky lines in object frames. The spatial resolution was $`0.^{\prime \prime }32`$ pixel<sup>-1</sup>.
For the hardest X-ray source (AX J131501+3141(119)), we made spectroscopic observations with the Kitt Peak National Observatories Mayall 4m and 2.1m telescopes. For details of the observations, see Akiyama et al. (1998a).
The data were analyzed using IRAF.<sup>1</sup><sup>1</sup>1 IRAF is distributed by the National Optical Astronomy Observatories, which is operated by the Association of Universities for Research in Astronomy, Inc. (AURA) under cooperative agreement with the National Science Foundation. After bias subtraction, flat-fielding, and wavelength calibration, optimum extraction method by apextract package was used to extract one dimensional spectral data from the two dimensional original data. For the UH data, flux calibrations were done with Feige 34. The flux calibrations did not work well in the wavelength larger than 7000Å and for some data in the wavelength shorter than 5000Å taken at very large zenith distances. For the Calar Alto data, flux calibrations were done with HD84937. The data were affected by fringes in the wavelength longer than 7000Å. Spectral fitting for emission lines was made by the $`\chi ^2`$ minimization method with specfit command in spfitpkg package in the IRAF. FWHMs of line widths were deconvolved by the spectral resolution shown above.
## 3 Results and Reliability
### 3.1 Results of Optical Identification
Many (25 out of 34) X-ray sources have one optical-counterpart candidate which shows broad permitted-emission lines whose FWHM is larger than those of forbidden-emission lines or 1000 km s<sup>-1</sup> in their optical spectra. The detected broad lines and line widths are summarized in Table Optical Identification of the ASCA Large Sky Survey. This suggests that most of the X-ray sources originate from AGNs.
In the remaining sources, 5 X-ray sources have only an object with narrow-emission lines. For 4 of them (AX J131758+3257(195), AX J131551+3237(171), AX J131501+3141(119), AX J131210+3048(072)), both the H$`\alpha `$ and H$`\beta `$ regions were observed. We classified these objects, using the \[NII\] 6583Å-to-H$`\alpha `$ and \[OIII\] 5007Å-to-H$`\beta `$ line ratios (Osterbrock (1989)). All of them show strong \[NII\] 6583Å line as well as an H$`\alpha `$ line and they fall in the region occupied by AGNs. For the remaining one object, AX J130840+2955(014), spectroscopic observations cover only the H$`\beta `$ emission line region, because of atmospheric absorption and fringes. It has a large \[OIII\] 5007Å-to-H$`\beta `$ flux ratio and the existence of an AGN is suggested, but the ratio is not large enough to exclude possibility of star-forming galaxy. However, the spectrum shows optical continuum which is dominated by old stellar populations and there is no indication of star-formation activity. Thus, the object is also identified with AGN.
We list notes on problematic cases, below.
The FIRST source (A in Figure Optical Identification of the ASCA Large Sky Survey) in the error circle was identified with a narrow-line AGN at a redshift of 0.245. There is also a blue object within the error circle (B) and it was identified with a broad-line AGN at a redshift of 1.577. Because most other X-ray sources whose X-ray spectra are similar to this X-ray source were identified with broad-line AGNs, we assigned the broad-line AGN as the optical counterpart of the X-ray source. It is possible that the narrow-line AGN or both is the origin of the X-ray source.
At the position of the FIRST source in the error circle ($`\alpha `$=13<sup>h</sup>16$`{}_{}{}^{\mathrm{m}}38.9_{}^{\mathrm{s}}`$, $`\delta `$=314957.8<sup>′′</sup>), no optical object brighter than $`R=22.5`$ mag can be seen. The nearest object, which is separated by 14<sup>′′</sup> from the radio position and has a blue color, was identified with a broad-line AGN at z=0.622. The projected distance between radio and optical source corresponds to 108 kpc at the redshift. If we assume that the radio emission is originated from the object, the radio-to-X-ray flux ratio corresponds to that of radio-loud objects (see Section 7.2).
The brightest stellar object (A in Figure Optical Identification of the ASCA Large Sky Survey) in the $`0.^{}8`$ error circle was identified with a G-type star. Based on an X-ray-to-optical flux ratio of a G-type star, this object cannot be the origin of the X-ray source. This source is not identified yet.
As a result, 34 X-ray sources in the SIS 2–7 keV 3.5$`\sigma `$ sample were identified with 30 AGNs including 5 objects with only narrow-emission lines, 2 clusters of galaxies, and 1 galactic star. Only 1 source is still unidentified. There are 2 X-ray sources which have HII-region like galaxies in their error circles. But they also have an optical counterpart candidate with AGN activity, thus we identified the X-ray sources with the AGNs. The resulting identification and classification for each X-ray source is indicated with bold characters in Table Optical Identification of the ASCA Large Sky Survey. The optical spectra of the identified objects are shown in Figure Optical Identification of the ASCA Large Sky Survey. Detected broad-lines in each object are listed in Table Optical Identification of the ASCA Large Sky Survey. Physical parameters of the identified objects are summarized in Table Optical Identification of the ASCA Large Sky Survey. In the next 2 subsections, we discuss the reliability of the identifications.
### 3.2 Contamination
Figure Optical Identification of the ASCA Large Sky Survey shows the distribution of optically identified objects in the X-ray error circles. Almost all X-ray sources were identified with optical objects within $`0.^{}6`$, though we picked up the candidates within $`0.^{}8`$. This result confirms the estimated 90% confidence error radius of $`0.^{}6`$.
We have checked the reliability of our identification by estimating chance contaminations in the error circle. Based on number counts of optically-selected broad-line AGN brighter than $`B=21.0`$ mag with a redshift smaller than 3 (Hartwick and Schade (1990)), the expected number of chance contaminations of broad-line AGN within 34 error circles is 0.29 and is negligible. For narrow-line AGNs, 3/5 of them, are pinpointed by the FIRST radio survey and chance contamination is very small for the sample as mentioned above. Therefore contamination of an object which is not an X-ray source is very low.
### 3.3 X-ray-to-optical Flux Ratio of the Identified Objects
An X-ray-to-optical flux ratio is one of useful tools to check the reliability of identification. To compare the X-ray-to-optical flux ratios of identified objects with $`\mathrm{log}f_\mathrm{X}/f_V`$ of the Einstein Medium Sensitivity Survey (EMSS) sample (Stocke et al. (1991)), we plot the relation between optical magnitude and hard X-ray flux in Figure Optical Identification of the ASCA Large Sky Survey. To plot the equal $`\mathrm{log}f_\mathrm{X}/f_V`$ lines, we converted the $`V`$ band magnitude to the $`R`$ band magnitude and 0.3–3.5 keV flux to 2–10 keV flux for the EMSS sample, using a $`VR`$ color of $`0.22`$ mag, which is equivalent to a typical optical energy index of broad line AGNs ($`0.5`$) and a photon index of 1.7 in the 0.3–10 keV band, respectively. The $`\mathrm{log}f_\mathrm{X}/f_V`$ values of the EMSS AGNs, clusters of galaxies, and galactic G-type stars are distributed $`1.0`$ to $`+1.2`$, $`0.5`$ to $`+1.5`$, and $`4.3`$ to $`2.4`$ (Stocke et al. (1991)), respectively. In the fainter flux range, the AGNs detected in the ROSAT HRI Lockman Hole survey also occupy the same value of $`\mathrm{log}f_\mathrm{X}/f_V`$ (Schmidt et al. (1998)). The $`\mathrm{log}f_\mathrm{X}/f_V`$ values of the identified AGNs, clusters of galaxies, and a galactic star in our sample are consistent with the distribution of the EMSS sample. The consistency of $`\mathrm{log}f_\mathrm{X}/f_V`$ with the EMSS objects supports the reliability of the identification in the LSS sample.
However, there are two AGNs which have values of $`\mathrm{log}f_\mathrm{X}/f_V`$ slightly outside of the EMSS sample. One AGN (AX J131831+3341(228)) has a $`\mathrm{log}f_\mathrm{X}/f_V`$ value of $`+1.3`$, which is X-ray louder than the EMSS AGNs. The AGN has a strong narrow \[OIII\] 5007Å line compared with the broad H$`\beta `$ line. Thus, the large $`\mathrm{log}f_\mathrm{X}/f_V`$ value might be explained by an optical absorption of its nucleus. The other AGN (AX J131805+3349(233)), whose $`\mathrm{log}f_\mathrm{X}/f_V`$ is smaller than $`1.0`$, is dominated by its host galaxy in the optical light (see Figure Optical Identification of the ASCA Large Sky Survey) and this component makes the object optically brighter in comparison with other AGNs. The lower limit of $`\mathrm{log}f_\mathrm{X}/f_V`$ on the optical counterpart of the unidentified source, AX J131832+3259(199), is $`+1`$. This value is consistent with that of the EMSS AGNs.
## 4 Optical and X-ray Spectral Properties of the Identified AGNs
### 4.1 Strength of Broad Emission Lines
To examine the strengths of broad-lines in comparison with narrow-lines, we measured the equivalent width ratios of broad plus narrow permitted lines-to-a narrow forbidden line. For objects at low redshifts, we measured the equivalent width ratios of H$`\alpha `$-to-\[NII\] 6583Å and H$`\beta `$-to-\[OIII\] 5007Å. For objects at redshifts between 0.3 and 0.75, we could measure only H$`\beta `$-to-\[OIII\] 5007Å ratio (we could not measure the ratio of AX J131639+3149(137), because of the fringe effect, thus the object is included in the higher redshift sample). The results are summarized in Table Optical Identification of the ASCA Large Sky Survey and shown in Figure Optical Identification of the ASCA Large Sky Survey. The objects with strong broad H$`\alpha `$ and H$`\beta `$ lines occupy the upper-right region of the figure. Most objects at redshifts lower than 0.3 do not show significant broad component in the H$`\beta `$ line and are placed in the lower part. About half of them show a broad component in the H$`\alpha `$ line and are distributed in the lower-right region. In the right box, among objects with redshifts between 0.3 and 0.75, most of the objects show significant broad component in the H$`\beta `$ line. In this redshift range, the fraction of objects with large H$`\beta `$-to-\[OIII\] 5007Å ratio is larger than the lower redshift sample in the left box. Based on the criteria used in Winkler (1992), the upper and lower boundary of the H$`\beta `$-to-\[OIII\] 5007Å ratio distribution in our sample correspond to type 1.2 and type 1.8 $``$ 2 Seyfert, respectively.
All 11 objects at redshift above 0.75 and AX J131639+3149(137) show broad MgII 2800Å or MgII 2800Å and CIII\] 1909Å lines. The equivalent widths and the line widths of their MgII 2800Å lines are summarized in Table Optical Identification of the ASCA Large Sky Survey. Their equivalent widths are consistent with those of a composite spectrum of optically-selected QSOs ($`50\pm 29`$Å; Francis et al. (1991)). AX J131707+3237(175) and AX J131021+3019(039) have a narrow-emission component in the MgII 2800Å and CIII\] 1909Å, respectively. Their line widths are as large as 1400 km s<sup>-1</sup> in FWHM and significantly larger than line widths of typical narrow-lines. Thus, there is no high-redshift luminous cousin of a narrow-line AGN in our sample. Such a deficiency of luminous narrow-line AGN is also reported in a radio selected sample of AGN (e.g., Lawrence (1991)) and a far-infrared selected sample (e.g., Barcons et al. (1995)).
### 4.2 X-ray Spectral Properties and the Effect of Absorption
The distribution of the apparent photon indices of the identified objects determined in the 0.7–10 keV band as a function of redshift is shown in Figure Optical Identification of the ASCA Large Sky Survey. All X-ray sources which have an apparent photon index smaller than 1 are identified with AGNs at redshifts smaller than 0.5. In the previous Figure Optical Identification of the ASCA Large Sky Survey, we mark the X-ray sources which have photon indices smaller than 1 with dots. AX J131551+3237(171), AX J131501+3141(119), AX J131210+3048(072), and AX J130840+2955(014) were identified with AGNs which show no significant broad line in H$`\alpha `$ and H$`\beta `$. Among them, the hardest source, AX J131501+3141(119), is identified with a type 2 Seyfert at z=0.072 (Akiyama et al. 1998a ; Sakano et al. (1998)). The remaining hard X-ray source (AX J130926+2952(016)) was identified with an AGN which shows a weak broad H$`\beta `$ line and has a small H$`\beta `$-to-\[OIII\]5007Å ratio. Thus, the object corresponds to a type 1.8 Seyfert (Winkler (1992)). These identifications support the picture (e.g., Comastri et al. (1995)) that the absorbed X-ray spectrum of narrow-line and weak-broad-line AGNs make the CXB spectrum harder than X-ray spectrum of broad-line AGN in the hard band. There still remains one hard source with a photon index of 0.58 unidentified.
Photon indices of other AGNs are distributed in the range from 1.2 to 2.1. There are three broad-line AGNs with an apparent photon index as small as 1.4 at intermediate to high redshifts. A dashed line in the figure shows the expected change of the apparent photon index in the 0.7–10 keV band with redshift for a type 1 Seyfert. It is derived by fitting a power-law model to the redshifted average X-ray spectrum of type 1 Seyferts (Gondek et al. (1996)) in the 0.7–10 keV band. Because of the reflection component, the observed photon index of type 1 Seyfert in the 0.7–10 keV band is expected to get harder to $`z2`$ and to reach a minimum photon index of 1.4 at a redshift of 2.5. The photon index distribution of the broad-line AGN sample is consistent with this line. Thus, the hardness of the apparent photon indices of the three high-redshift broad-line AGNs may be explained by the existence of the reflection component. It should, however, be noted that the existence and the strength of a reflection component in AGNs with a luminosity of 10<sup>45</sup> erg s<sup>-1</sup> is not confirmed, so far.
For the AGNs, we also fitted their X-ray spectra with an absorbed power-law model, assuming absorbing matter at the object’s redshift and an intrinsic photon index of 1.7. The results are shown in column 8 of Table Optical Identification of the ASCA Large Sky Survey and shown in Figure Optical Identification of the ASCA Large Sky Survey as a function of X-ray luminosities. The upper boundary of the distribution of the column density is consistent with an estimated limit on the column density in the LSS sample (see, Section 6). Hard X-ray sources which were identified with narrow-line or weak-broad-line AGNs at redshift smaller than 0.5 are fitted with column densities of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$. Most of the other AGNs are fitted with column densities less than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$. However, four high-redshift ($`2>z>0.5`$) luminous ($`L_{210\mathrm{keV}}=10^{44.545.5}`$ erg s<sup>-1</sup>) broad-line AGNs (AX J131816+3240(183), AX J131054+3004(037), AX J131724+3203(152), and AX J131021+3019(039)) are fitted with large column densities. Such large neutral hydrogen column densities conflict with the existences of broad MgII lines and strong UV continua in their optical spectra. However, there is a possibility that optical extinction is not always strongly correlated with X-ray absorption (e.g., NGC4151), as in the case where the X-ray absorbing gas is located within the dust sublimation radius, the absorbing column density is variable, or if the gas-to-dust ratio or the composition is different from the Galactic interstellar gas. These hardness may also be explained by the effect of the reflection component. The presence of the reflection component makes observed type 1 AGN spectrum harder than a photon index of 1.7 in the 0.7–10 keV band at high-redshift, as shown in Figure Optical Identification of the ASCA Large Sky Survey. The slightly hardened X-ray spectrum in the observed band is spuriously fitted with large column density in the rest-frame band (1.8–25 keV at a redshift of 1.5) at a high redshift. To distinguish the effect of a real X-ray absorption and the reflection component, X-ray spectroscopic observations for such broad-line AGNs below 0.7 keV in the observed frame are critical.
To examine the correlation between the strength of optical broad-line and the X-ray absorption column density, we plotted the AGNs with redshift smaller than 0.75 in the H$`\beta `$-to-\[OIII\] 5007Å versus hydrogen column density diagram in Figure Optical Identification of the ASCA Large Sky Survey. The narrow- or weak-broad-line AGNs with $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$ have 10 times smaller values of H$`\beta `$/\[OIII\]5007Å than typical QSOs. The small ratio suggests existence of absorption with $`A_V`$ of larger than 2$``$3 mag to the broad-line region. The optical extinction corresponds to a hydrogen column density of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=21.521.7`$, according to the relation $`N_\mathrm{H}/A_V=1.79\times 10^{21}`$ cm<sup>-2</sup> mag<sup>-1</sup>, which is determined from observations of Galactic objects (e.g., Predehl & Schmitt (1995)). Thus, the optical extinction derived from the H$`\beta `$ region is consistent with the X-ray spectra for these objects. One narrow-AGN and 5 weak-broad-line AGNs with $`\mathrm{log}(\mathrm{H}\beta /[\mathrm{OIII}]5007\mathrm{\AA })`$ of less than 0 are fitted with smaller hydrogen column densities than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$. In the H$`\alpha `$ region, however, 3/6 of them have a significant broad H$`\alpha `$ line and are plotted with open rectangles or an open star in Figure Optical Identification of the ASCA Large Sky Survey. They fall in the lower-right region in Figure Optical Identification of the ASCA Large Sky Survey. The small hydrogen column density is consistent with the existence of a broad H$`\alpha `$ line, which suggests smaller absorption to the broad-line region than AGNs with $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$. All strong-broad-line AGNs with $`\mathrm{log}(\mathrm{H}\beta /[\mathrm{OIII}]5007\mathrm{\AA })`$ larger than 0 are fitted with hydrogen column density of less than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$. These facts suggest the critical column density which divides the narrow-line and broad-line AGN is around $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$. The critical column density is consistent with the lower boundary of the distribution of absorbing column densities of Seyfert 2 galaxies (Risaliti, Maiolino, & Salvati 1999). Therefore, for AGNs with redshift smaller than 0.75, there is no broad-line AGN with a large absorption column density and the strengths of the broad Balmer lines are consistent with the hydrogen column density derived from the X-ray spectra.
## 5 Number Counts and Contribution to the CXB
We hereafter divide the AGN sample at the intrinsic column density of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$ : we refer the AGNs with the column density larger and smaller than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$ as “the absorbed AGN” and “the less-absorbed AGN”, respectively. At first, we include the four high-redshift broad-line AGNs fitted with the large column densities in the less-absorbed sample, because such large column densities could be only apparent (Section 4.2).
Figure Optical Identification of the ASCA Large Sky Survey shows the cumulated logN-logS relations for the absorbed and for the less-absorbed samples in the SIS $`3.5\sigma `$ sample. Since the sensitivity limit is given in count rates, the actual flux limit depends on X-ray spectrum of sources. For an X-ray spectrum with the canonical photon index of type 1 AGNs ($`\mathrm{\Gamma }=1.7`$), the deepest source-detection limit of the SIS $`3.5\sigma `$ sample (1.2 cts ksec<sup>-1</sup>; see Table 1) is equivalent for $`1.1\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$in the 2–10 keV band. If we assume an object with hydrogen column density of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=23`$ at redshift of 0.1, which corresponds to the hardest object in the LSS sample, the flux limit is estimated to be $`1.6\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$in the 2–10 keV band. In the calculation of Figure Optical Identification of the ASCA Large Sky Survey, the SIS count rate was converted into the flux with the best-fit photon index determined in the 0.7–10 keV band for each source to take into account this effect. As noted from Figure Optical Identification of the ASCA Large Sky Survey, the surface number density of the absorbed AGNs is comparable to that of the less absorbed AGNs at flux larger than $`2\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. The large fraction of the less-absorbed AGNs in the identified sample is due to different limiting flux for different X-ray spectra.
At the first step, we estimate the contributions to the CXB, integrating the logN-logS relations from the brightest to the faintest object in each population of the LSS sample. For less-absorbed AGNs, absorbed AGNs, and clusters of galaxies, the summed fluxes of sources per unit area are $`3.3\times 10^{12}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$degree<sup>-2</sup>, $`6.3\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$degree<sup>-2</sup>, and $`1.4\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$degree<sup>-2</sup> and contributions to the CXB are 17%, 3%, and 0.7% , respectively, if we use the CXB flux determined in the LSS field (Ishisaki et al. (1999), $`2.0\times 10^{11}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$degree<sup>-2</sup> in the 2–10 keV band).
As the second step, to determine contributions of each population in a wider flux range, we combined results of the HEAO1 A2 sample (Piccinotti et al. (1982)). The sample consists of 30 AGNs, 4 BL Lac objects, 30 clusters of galaxies, 1 starburst galaxy, and 3 unidentified sources with the flux limit of $`3.1\times 10^{11}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$ in the 2–10 keV band. In the 30 AGNs, 6 AGNs are fitted with hydrogen column density larger than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$ (Schartel et al. (1997)). The source number density at the flux limit is $`2.4\times 10^3\mathrm{degree}^2`$. If we assume that there is no break in the logN-logS relations, the slopes of the relations between the HEAO1 A2 and the LSS limits can be determined. The results are shown in Table 5. In this flux range, the total number count in the 2–10 keV band has a slope ($`\alpha `$ of $`N(>S)=kS^\alpha `$) of $`1.5`$ and is consistent with that of smoothly distributed sources in a simple Euclidean geometry as reported in Ueda et al. (1998). However, the slope is a summation of steep slopes for AGNs and a shallow slope for clusters of galaxies. Thus the Euclidean distribution is an apparent effect. The slope of the logN-logS relation of clusters of galaxies in the flux range is consistent with that determined in the 0.5–2 keV band ($`1.15`$, Vikhlinin et al. (1998)). Based on the slope of the logN-logS relation, the contributions from less-absorbed AGNs, absorbed AGNs, and clusters of galaxies are estimated to be 9%, 4%, and 1%, respectively, integrating the logN-logS relations from $`3\times 10^{11}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$down to the flux limit of $`2\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, which is the limit for absorbed AGNs and at which 16% of the CXB was resolved into discrete sources. The logN-logS relation of the absorbed AGNs with the four high-redshift broad-line AGNs is also shown in Figure Optical Identification of the ASCA Large Sky Survey with a thin long-dashed line. Based on the logN-logS relation, the contribution of the absorbed AGNs with the four AGNs to the CXB is estimated to be 6%.
## 6 Redshift Distribution of Identified AGNs and Deficiency of Obscured QSO
To examine the redshift and luminosity distributions of the AGNs, we plot the X-ray luminosity versus redshift diagram of all the identified AGNs of the SIS $`3.5\sigma `$ sample in Figure Optical Identification of the ASCA Large Sky Surveya. Absorbed AGNs are marked with dots. At first, the high-redshift broad-line AGNs which are fitted with large column densities are included in the less-absorbed sample as in Section 5. The number ratio of absorbed AGNs-to-less-absorbed AGNs is clearly changing with redshift, luminosity, or both. If we limit the sample above the flux level of $`2\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$, which is the flux limit for the absorbed AGNs, the redshift and luminosity distribution of the limited sample is also significantly different between these two populations. The number count of the absorbed AGNs is dominated by nearby low-luminous objects, in contrast to that of less-absorbed AGNs, which are dominated by QSOs at intermediate to high redshifts. In the same figure, the AGN sample from the HEAO1 A2 survey (Piccinotti et al. (1982)) is also plotted with small marks. Objects with hydrogen column density larger than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$ are also marked with small dots (Schartel et al. (1997)). The same tendency of deficiency of absorbed AGNs with a large luminosity is seen in their sample. Such a deficiency of luminous absorbed AGNs is also found in the ROSAT, ASCA, and Beppo-SAX deep surveys in the Lockman Hole (Hasinger et al. (1999)). In Figure Optical Identification of the ASCA Large Sky Surveyb, detection limits of the SIS 2–7 keV 3.5$`\sigma `$ sample for intrinsic luminosities with various absorption column densities is shown as a function of redshift. In this calculation, we assume the intrinsic photon index of 1.7 and consider the response function of ASCA with the SIS. The difference between the survey limits of non-absorbed objects and that of objects with $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22.5`$ is less than 0.2 in the logarithmic-scale intrinsic luminosity at the redshift range from 0 to 1.5. Thus, we can detect objects with intrinsic column densities of up to $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22.5`$ at redshifts smaller than 1.5 in an unbiased fashion. At higher redshifts objects with column densities of up to $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=23`$ are detected without bias, thanks to the redshift effect. It should be noted that the peak of the distribution of absorption column densities of type 2 Seyferts in the nearby universe ($`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)>23`$; Risaliti et al. (1999)) is beyond the detection limit.
To examine the deficiency of luminous absorbed AGNs more quantitatively, we compare the redshift distribution of absorbed AGNs and that of less-absorbed AGNs, considering the difference of their X-ray spectra. In Figure Optical Identification of the ASCA Large Sky Surveya, we show the redshift distribution of less-absorbed AGNs as well as the expected redshift distribution of broad-line AGNs. The four high-redshift broad-line AGNs which are fitted with large column densities are included in the less-absorbed sample. The expected redshift distribution of broad-line AGNs is calculated based on a hard-band X-ray luminosity function of broad-line AGNs and its evolution (Boyle et al. 1998a ); they are determined by an AGN sample from ASCA observations of Deep ROSAT fields and optical identifications of the Large Area Sky Survey/Modulation Collimator catalog obtained from HEAO1 mission. The expected redshift distribution matches well to that of the LSS sample. Based on Kolmogorov-Smirnov test, the probability that the observed redshift distribution agrees with the model is calculated to be 64%. Thus, we use the hard-band X-ray luminosity function as a base model for less-absorbed AGNs. In Figure Optical Identification of the ASCA Large Sky Surveyb, we compare the redshift distribution of absorbed AGNs in the LSS sample with the expected distribution of them calculated by assuming that the shape of the intrinsic (which means absorption corrected) luminosity function of absorbed AGNs is the same as that of less-absorbed AGNs. The equal number density of AGNs with $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$ to that of less-absorbed AGN is the same as in the spectrum model of the CXB; e.g., Comastri et al. (1995) used the value of 1.23 for the number ratio of AGNs with absorption of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$ to those with absorption less than $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$. As seen in Figure Optical Identification of the ASCA Large Sky Surveyb, at redshift smaller than 0.4, in which the X-ray luminosity of the sample AGN is less than $`10^{44}`$ erg s<sup>-1</sup>, the number of absorbed AGNs is comparable to the expectation. However, at redshifts larger than 0.4, the total expected number is about 10, which contrasts to the non-detection. The observed redshift distribution is different from the model redshift distribution; Kolmogorov-Smirnov test gives a probability of only 5% for the null hypothesis that the redshift distribution of absorbed AGNs is the same as the model, thus the hypothesis is rejected. Therefore, our sample suggests the deficiency of the absorbed AGNs with the column density of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$, in the redshift range between 0.5 and 2.0, or in the X-ray luminosity range larger than $`10^{44}`$ erg s<sup>-1</sup>, or both. We plot the redshift distribution of the four high-redshift broad-line AGNs with hard X-ray spectra in Figure Optical Identification of the ASCA Large Sky Surveyb with a histogram shaded with dashed lines. By including the four AGNs, the observed redshift distribution agrees the model redshift distribution with probability of 60%. If their large column densities are real, they could complement the deficiency of absorbed high-redshift luminous AGNs. Revealing the column density distribution of absorbed AGN as a function of redshift and luminosity and constructing the model of the origin of the CXB more accurately will be an important objective for the next generation X-ray surveys.
## 7 Properties of identified AGNs in other wavelengths
### 7.1 Optical Photometric Properties
We have examined the optical appearance of each identified object. The results are indicated in the last column of the Table Optical Identification of the ASCA Large Sky Survey. All of the low-luminosity AGNs show their host galaxy component in the optical light. We calculate optical total absolute magnitudes of the identified objects, using $`\alpha =0.5`$, which corresponds to $`VR`$ of 0.22 mag. The results are shown in Table Optical Identification of the ASCA Large Sky Survey and plotted in Figure Optical Identification of the ASCA Large Sky Survey as a function of X-ray luminosity. For high-luminosity AGNs, hard X-ray luminosities correlate well with optical luminosities. On the other hand, for the low-luminosity AGNs whose optical light is affected and maybe dominated by their host galaxy components, the correlation is broken. Considering the break, we estimate the absolute magnitudes of the host galaxies to be around $`M_V=2123.5`$ mag and are similar to that of QSO host galaxies (Bahcall et al. (1997)).
The optical colors of luminous broad-line AGNs distribute from $`BR`$ of $`0.2`$ mag to $`1.1`$ mag, consistent with the expected color range for broad-line AGNs. There is no broad-line QSO which has a $`BR`$ color larger than $`2.0`$ mag, which is a criterion for red QSOs in Kim & Elvis (1999). Two low-luminosity broad-line AGNs (AX J131725+3300(192) and AX J131407+3158(127)) which show a broad H$`\alpha `$ emission have red optical colors. The red colors are probably affected by the host galaxy, because their optical images show the host galaxy components.
### 7.2 Radio and Far-infrared Properties
To examine radio properties of the identified AGNs, we show a radio versus hard X-ray luminosity diagram in Figure Optical Identification of the ASCA Large Sky Survey. There are two sequences in the diagram; three objects which have large radio-to-hard X-ray luminosity ratios correspond to radio-loud AGNs and others which have smaller ratios correspond to radio-quiet AGNs. The radio-to-hard X-ray ratios of the radio-loud sequence are consistent with the average SED of radio-loud QSO (Elvis et al. (1994)). On the other hand, the radio-quiet sequence has one order of magnitude larger radio-to-hard X-ray luminosity ratio than the average SED of radio-quiet QSO (Elvis et al. (1994)). Based on the dispersion of the radio-to-optical luminosity ratios of radio-quiet QSOs (Visnovsky et al. (1992)), the large radio-to-hard X-ray luminosity ratio is accounted by a scatter and a selection bias for radio louder objects. The fraction (3/30) of radio-loud AGNs in the total AGN sample is consistent with that in the low-luminosity sample of optically selected QSOs (13%) (Visnovsky et al. (1992)). In radio images, all radio-quiet AGNs show only one point-like component. Two radio-loud AGNs consist of two components; one dominates the total flux located on the optical position and the other fainter source is $`15`$ arcsec away from the central component.
In the LSS sample, only the galactic star is detected in the IRAS survey (Moshir et al. (1992)). None of the AGN is detected in the IRAS Faint Source Survey. Based on the flux limit of the IRAS survey in the LSS region (0.1 Jy in 60 $`\mu `$m), most of the AGNs have lower limits to the logarithmic ratio of 2–10 keV luminosity-to-far-infrared luminosity ($`\nu _{60\mu \mathrm{m}}L\nu _{60\mu \mathrm{m}}`$) of less than $`1`$, consistent with those known for QSOs ($`1`$) and only one object (AX J131822+3347 (235)) has a lower limit ($`0.6`$) to the logarithmic ratio larger than the typical value of QSOs. The lower limits to the logarithmic ratio are significantly larger than those of star-forming galaxies, which have logarithmic ratios around $`4`$.
## 8 Summary
We present results of optical identifications of the X-ray sources detected in the ASCA Large Sky Survey. Optical spectroscopic observations were done for 34 X-ray sources which were detected with the SIS in the 2–7 keV band above 3.5 $`\sigma `$. The flux limit corresponds to $`1\times 10^{13}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$in the 2–10 keV band. The sources are identified with 30 AGNs, 2 clusters of galaxies, and 1 galactic star. Only 1 source is still unidentified.
All of the X-ray sources that have a hard X-ray spectrum with an apparent photon index of smaller than 1 in the 0.7–10 keV band are identified with narrow-line or weak-broad-line AGNs at redshifts smaller than 0.5. This fact supports the idea that absorbed X-ray spectra of narrow-line and weak-broad-line AGNs make the Cosmic X-ray Background (CXB) spectrum harder in the hard X-ray band than that of a broad-line AGN, which is the main contributor in the soft X-ray band. Assuming their intrinsic spectra are same as a broad-line AGN (a power-law model with a photon index of 1.7), their X-ray spectra are fitted with hydrogen column densities of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$ at the object’s redshift. On the other hand, X-ray spectra of the other AGNs are consistent with that of a nearby type 1 Seyfert. In the sample, four high-redshift luminous broad-line AGNs show a hard X-ray spectrum with an apparent photon index of $`1.3\pm 0.3`$. The hardness may be explained by the reflection component of a type 1 Seyfert. The hard X-ray spectra may also be explained by absorption with $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=2223`$ at the object’s redshift, if we assume an intrinsic photon index of 1.7. The origin of the hardness is not clear yet.
Based on the logN-logS relations of each population, contributions to the CXB in the 2–10 keV band are estimated to be 9% for less-absorbed AGNs ($`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)<22`$) including the four high-redshift broad-line AGNs with a hard X-ray spectrum, 4% for absorbed AGNs ($`22<\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)<23`$, without the four hard broad-line AGNs), and 1% for clusters of galaxies in the flux range from $`3\times 10^{11}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$to $`2\times 10^{13}`$$`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. If the four hard broad-line AGNs are included in the absorbed AGNs, the contribution of the absorbed AGNs to the CXB is estimated to be 6%.
In optical spectra, there is no high-redshift luminous cousin of a narrow-line AGN in our sample. The redshift distribution of the absorbed AGNs are limited below $`z=0.5`$ excluding the four hard broad-line AGNs, in contrast to the existence of 15 less-absorbed AGNs above $`z=0.5`$. The redshift distribution of the absorbed AGNs suggests a deficiency of AGNs with column densities of $`\mathrm{log}N_\mathrm{H}(\mathrm{cm}^2)=22`$ to $`23`$ in the redshift range between 0.5 and 2, or in the X-ray luminosity range larger than $`10^{44}`$ erg s<sup>-1</sup>, or both. If the large column densities of the four hard broad-line AGNs are real, they could complement the deficiency of X-ray absorbed luminous high-redshift AGNs.
MA, KO, and TY would like to thank S.Okamura, M.Sekiguchi and the MOSAIC CCD camera team, and staff members of the KISO observatory for their support during the imaging observations. MA, KO, and TY appreciate the support from members of the University of Hawaii observatory during the spectroscopic observations. KO is grateful to the hospitality during his stay at the Institute for Astronomy, University of Hawaii, where a part of this work was done. MA, KO, YU, and IL wish to thank members of the Calar Alto observatory for their help during the spectroscopic observations. We are also grateful to the referee for his useful suggestions. This research has made use of NASA/IPAC Extragalactic Database (NED), which is operated by the Jet Propulsion Laboratory, Caltech, under contract with the National Aeronautics and Space Administration. MA, WK, and MS acknowledge support from a Research Fellowships of the Japan Society for the Promotion of Science for Young Scientists. The optical follow-up program is supported by grants-in-aid from the Ministry of Education, Science, Sports and Culture of Japan (06640351, 08740171, 09740173) and from the Sumitomo Foundation.
## Appendix A ROSAT HRI Data Reduction and Source Detection
ROSAT HRI observations were made in 2 fields with a 20 ksec exposure in December 1997. These fields were centered at $`\alpha `$=13<sup>h</sup>14$`{}_{}{}^{\mathrm{m}}36_{}^{\mathrm{s}}`$, $`\delta `$=320112<sup>′′</sup> (LSS-HRI1) and $`\alpha `$=13<sup>h</sup>12$`{}_{}{}^{\mathrm{m}}29_{}^{\mathrm{s}}`$, $`\delta `$=311312<sup>′′</sup> (LSS-HRI2) (J2000) with a radius of $`19^{}`$.
These data were reduced with EXSAS package (Zimmermann et al. (1994)) on the MIDAS. For the source detection, we applied the detect command selecting photons in the PHA channel range from 2 to 8, only. X-ray sources detected in the HRI follow-up observations with Maximum Likelihood larger than 10 were listed in Table Optical Identification of the ASCA Large Sky Survey.
The conversion factor from a count rate to a 0.5–2.0 keV flux is $`1.854\times 10^{14}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$cts<sup>-1</sup> ksec. Thus the flux of the faintest source corresponds to $`1.13\times 10^{14}`$ $`\mathrm{erg}\mathrm{cm}^2\mathrm{s}^1`$. Considering systematic shifts between candidates of optical counterparts and X-ray sources, we shifted the X-ray positions with $`\mathrm{\Delta }\alpha =3.^{\prime \prime }39`$, $`\mathrm{\Delta }\delta =+5.^{\prime \prime }08`$ and $`\mathrm{\Delta }\alpha =2.^{\prime \prime }37`$, $`\mathrm{\Delta }\delta =+3.^{\prime \prime }18`$ for regions LSS-HRI1 and LSS-HRI2, respectively.
## Appendix B Cross identification with the ROSAT PSPC catalogs and the VLA FIRST catalog
We summarize the identifications of the ASCA LSS sources with the ROSAT PSPC sources and the VLA FIRST radio sources in Table Optical Identification of the ASCA Large Sky Survey. |
warning/0001/gr-qc0001018.html | ar5iv | text | # On the Structure of Solutions to the Static Vacuum Einstein Equations
## 0. Introduction
The static vacuum Einstein equations are the equations
(0.1)
$$ur=D^2u,$$
$$\mathrm{\Delta }u=0,$$
on a Riemannian 3-manifold $`(M,g)`$, with $`u`$ a positive function on $`M`$. Here $`r`$ denotes the Ricci curvature, $`D^2`$ the Hessian, and $`\mathrm{\Delta }=trD^2`$ the Laplacian on $`(M,g)`$. Solutions of these equations define a Ricci-flat 4-manifold $`N`$, of the form $`N=M\times _uS^1`$ or $`N=M\times _u`$, with Riemannian or Lorentzian metric of the form
(0.2)
$$g_N=g_M\pm u^2dt^2.$$
These equations are the simplest equations for Ricci-flat 4-manifolds. They have been extensively studied in the physics literature on classical relativity, where the solutions represent space-times outside regions of matter which are translation and reflection invariant in the time direction $`t`$. However, with the exception of some notable instances, (c.f. Theorem 0.1 below), many of the global properties of solutions have not been rigorously examined, either from mathematical or physical points of view, c.f. \[Br\] for example.
This paper is also motivated by the fact that solutions of the static vacuum equations arise in the study of degenerations of Yamabe metrics (or metrics of constant scalar curvature) on 3-manifolds, c.f. \[A1\]. Because of this and other related applications of these equations to the geometry of 3-manifolds, we are interested in general mathematical aspects of the equations and their solutions which might not be physically relevant; for example, we allow solutions with negative mass.
In this paper, we will be mostly concerned with the geometry of the 3-manifold solutions $`(M,g,u)`$ of (0.1), (i.e. the space-like hypersurfaces), and not with the 4-manifold metric. Thus, the choice of Riemannian or Lorentzian geometry on $`N`$ in (0.2) will play no role. This considerably simplifies the discussion of singularities and boundary structure, but still allows for a large variety of behaviors; c.f. \[ES\] for a survey on singularities of space-times.
Obviously, there are no non-flat solutions to (0.1) on closed manifolds, and so it will be assumed that $`M`$ is an open, connected oriented 3-manifold. Let $`\overline{M}`$ be the metric (or Cauchy) completion of $`M`$ and $`M`$ the metric boundary, so that $`\overline{M}=MM`$ is complete as a metric space.
In order to avoid trivial ambiguities, we will only consider maximal solutions of the equations (0.1). For example any domain $`\mathrm{\Omega }`$ in $`^3`$ with the flat metric, and $`u`$ a positive constant, satisfies (0.1). In this case, the metric boundary $`\mathrm{\Omega }`$ is artificial, and has no intrinsic relation with the geometry of the solution. The solution obviously extends to a larger domain, i.e. $`^3.`$ Thus, we only consider maximal solutions $`(M,g,u)`$, in the sense that $`(M,g,u)`$ does not extend to a larger domain $`(M^{},g^{})(M,g)`$ with $`u>`$ 0 on $`M^{}.`$ It follows that at the metric boundary $`M`$ of $`M`$, either the metric or $`u`$ degenerates in some way or $`u`$ approaches 0 in some way, (or a combination of such).
A classical result of Lichnerowicz \[L1, p.137\] implies that if the metric $`(M,g)`$ satisfying (0.1) is complete, (i.e. $`M=\mathrm{}),`$ and $`u`$ 1 at infinity, then $`u`$ 1 and $`M`$ is flat, i.e. $`^3`$ or a quotient of $`^3.`$ More generally, it is proved in \[A1, Thm.3.2\] that if $`(M,g)`$ is a complete solution to (0.1), (hence $`u>`$ 0 everywhere), then $`(M,g)`$ is flat and $`u`$ is constant, i.e. the assumption on the asymptotic behavior of $`u`$ is not necessary, c.f. also Theorem 1.1 below.
Thus, there are no complete non-trivial solutions to (0.1) and hence $`M`$ must be non-empty. The set formally given by
$$\mathrm{\Sigma }=\{u=0\}\overline{M},$$
is called the horizon. It is closely related to the notion of event horizon in general relativity. More precisely, $`\mathrm{\Sigma }`$ may be defined as the set of limit points of Cauchy sequences on $`(M,g)`$ on which $`u`$ converges to $`0`$. Although $`M\mathrm{}`$, it is possible that $`\mathrm{\Sigma }=\mathrm{}`$. However, most solutions of physical interest do have $`\mathrm{\Sigma }\mathrm{}`$. In the framework of classical relativity, the non-triviality of a static vacuum solution, i.e. the non-vanishing of its curvature, is due to the presence of matter or field sources at $`M,`$ or ’inside‘ the horizon $`\mathrm{\Sigma }`$ in case $`(M,g)`$ extends as a vacuum solution past $`\mathrm{\Sigma }.`$
It is natural to consider the situation where $`(M,g)`$ is not complete and for which the metric boundary $`M`$ of $`M`$ coincides with the horizon $`\mathrm{\Sigma }`$. More precisely, we will say that $`(M,g)`$ is complete away from the horizon $`\mathrm{\Sigma }`$ if for any sequence $`p_ipM`$ in the metric topology on $`\overline{M}=MM`$ one has $`u(p_i)`$ 0. Conversely if $`\{p_i\}`$ is a bounded sequence in $`M`$ with $`u(p_i)`$ 0, then the definition of $`(M,g,u)`$ implies that a subsequence of $`\{p_i\}`$ converges to a point $`pM.`$ Thus, $`M=\mathrm{\Sigma }`$ is given by the Hausdorff limit of the $`\epsilon `$ -levels $`L^\epsilon `$ of $`u`$ as $`\epsilon `$ 0.
While most solutions $`(M,g)`$ of physical interest are complete away from $`\mathrm{\Sigma },`$ there are many solutions for which this is not the case, c.f. §2 for further discussion. In such examples, the curvature typically blows up within a finite distance to $`\mathrm{\Sigma }.`$ Among the solutions which are complete away from $`\mathrm{\Sigma },`$ most all are singular at $`\mathrm{\Sigma },`$ again in the sense that the curvature of the metric $`g`$ blows up on approach to $`\mathrm{\Sigma }.`$ This is closely related to the fact that the equations (0.1) are formally degenerate at $`\mathrm{\Sigma }.`$ Thus, in general, even when $`M`$ is complete away from $`\mathrm{\Sigma },`$ the metric completion $`M\mathrm{\Sigma }`$ need not be a smooth manifold with boundary.
We will say that $`(M,g,u)`$ extends smoothly to $`\mathrm{\Sigma },`$ if (i): the set $`\mathrm{\Sigma }`$ is a smooth surface and the partial completion $`M^{}=M\mathrm{\Sigma }\overline{M}`$ is a smooth manifold with boundary, (ii): the metric $`g`$ extends smoothly to a smooth Riemannian metric on $`M^{},`$ and (iii): the potential $`u`$ extends smoothly to $`\mathrm{\Sigma }`$ with $`\mathrm{\Sigma }=\{u=0\};`$ smoothness here means at least $`C^2.`$ Note that it might be possible that $`\mathrm{\Sigma }`$ has infinitely many components, or non-compact components of infinite topological type. In any case, one immediate strong consequence of (0.1) is that if $`g`$ extends smoothly to $`\mathrm{\Sigma },`$ then at $`\mathrm{\Sigma }`$ one has
(0.3)
$$D^2u|_\mathrm{\Sigma }=0.$$
Thus, $`\mathrm{\Sigma }`$ is a totally geodesic surface in $`\overline{M}`$ and $`|u|`$ is a non-zero constant on each component of $`\mathrm{\Sigma },`$ c.f. Remark 1.5 below. Observe that if $`(M,g)`$ is smooth up to $`\mathrm{\Sigma },`$ and complete away from $`\mathrm{\Sigma },`$ then the isometric double of $`M`$ across $`\mathrm{\Sigma }`$ is a smooth complete Riemannian manifold. The harmonic function $`u`$ extends smoothly across $`\mathrm{\Sigma }`$ as harmonic function, odd w.r.t. reflection in $`\mathrm{\Sigma }.`$
By far the most significant solution of the static vacuum equations is the Schwarzschild metric, of mass $`m`$, given on the space-like hypersurface $`M`$ by
(0.4)
$$g_S=(1\frac{2m}{r})^1dr^2+r^2ds_{S^2}^2,u=(1\frac{2m}{r})^{1/2},r>2m.$$
This metric models the vacuum exterior region of an isolated static star or black hole. It is a spherically symmetric metric on $`M`$ = $`(2m,\mathrm{})\times S^2`$ and has $`\mathrm{\Sigma }`$ given by a (totally geodesic) symmetric $`S^2,`$ of radius $`2m`$. The mass $`m`$ is usually assumed to be positive, but we will not make this assumption here. Thus, we allow $`m`$ 0. Of course if $`m=`$ 0, then $`g_S`$ is just the flat metric with $`u=`$ 1. If $`m<`$ 0, (0.4) is understood to be defined for $`r>`$ 0.
The Schwarzschild metric is asymptotically flat in the sense that there is a compact set $`K`$ in $`M`$ such that $`MK`$ is diffeomorphic to $`^3B(R),`$ and in a suitable chart on $`MK`$, the metric approaches the Euclidean metric at a rate of $`1/r,`$ i.e.
(0.5)
$$g_{ij}=(1+\frac{2m}{r})\delta _{ij}+O(1/r^2),$$
with curvature decay of order $`1/r^3,r=|x|,`$ and with $`m.`$ The function $`u`$ (up to a multiplicative constant) has the asymptotic form $`u=`$ 1 $`\frac{m}{r}+O(r^2)`$ with $`|u|=O(1/r^2).`$ A triple $`(M,g,u)`$ satisfying these conditions is called asymptotically flat.
We note the following remarkable characterization of the Schwarzschild metric.
###### Theorem 0.1.
(Black-hole uniqueness),\[I1\],\[Ro\],\[BM\]
Let $`(M,g,u)`$ be a solution of the static vacuum Einstein equations, which is smooth up to $`\mathrm{\Sigma }`$ and complete away from $`\mathrm{\Sigma }.`$ If $`\mathrm{\Sigma }`$ is a compact, (possibly disconnected) surface and $`(M,g,u)`$ is asymptotically flat, then $`(M,g,u)`$ is the Schwarzschild metric, with $`m>0`$.
The hypothesis that the space-like hypersurface $`(M,g,u)`$ is asymptotically flat is very common in physics. Namely, in modeling a static space-time outside an isolated, i.e. compact, field or matter source, it is natural to assume that in regions far away from the source the geometry of space approximates that of $`^3,`$ i.e. empty space.
Nevertheless, mathematically the asymptotically flat assumption is quite strong in that it severely restricts both the topology and geometry of $`(M,g)`$ outside a large compact set. Further, the physical reasoning above presupposes that there are no complete non-flat solutions to the vacuum equations (0.1), i.e. that a static gravitational field is non-empty solely due the presence of matter somewhere.
###### Remark 0.2.
This latter issue in fact led Einstein to hypothesize that space $`M`$ is compact, in order to avoid dealing with ’artificial‘ boundary value problems at infinity, c.f. \[E, p.98ff.\]. This issue, closely related to Mach’s Principle, is discussed in some detail in work of Lichnerowicz, (Propositions A and $`B`$ of \[L1, §31\] and \[L2, Ch.II\]); c.f. also \[MTW, §21.12,\],\[Ri, §9.12\] for instance.
As remarked above, there are in fact no complete non-trivial static vacuum solutions, so that the asymptotically flat assumption in Theorem 0.1 may not be unreasonable, c.f. however \[El\], \[G\]. In fact, the main result of this paper is that this hypothesis is not necessary in most circumstances; it follows from much weaker assumptions.
To explain this, we first need to consider a weakening of the condition that $`M`$ is compact. Let
(0.6)
$$t(x)=dist(x,M).$$
Define $`M`$ to be pseudo-compact if there is a tubular neighborhood $`U`$ of $`M`$ whose boundary $`U`$ in $`M`$, $`UM`$, is compact, i.e. $`\{t(x)=s_o\}`$ is compact, for some $`s_o>`$ 0, (and hence all $`0<s_o<\mathrm{}`$). As will be seen in §2, there are numerous examples of static vacuum solutions with $`M`$ pseudo-compact but not compact.
Let $`E`$ be an end, i.e. an unbounded component, of $`\overline{M}U.`$ The mass of $`E`$ may be defined by
(0.7)
$$m_E=lim_s\mathrm{}m_E(s)=lim_s\mathrm{}\frac{1}{4\pi }_{S_E(s)}<logu,t>dA,$$
where $`S_E(s)=t^1(s)E`$. Since $`u`$ is harmonic, log $`u`$ is superharmonic, so that the divergence theorem implies that $`m_E(s)`$ is monotone decreasing in $`s`$. Hence the limit (0.7) is well-defined, (possibly $`\mathrm{}`$). Note that the static vacuum equations are invariant under multiplication of $`u`$ by positive constants. We use the $`logu`$ term in (0.7) in place of $`u`$ so that the mass is independent of this rescaling in $`u`$.
The following is the main result of this paper.
###### Theorem 0.3.
Let $`(M,g,u)`$ be a solution of the static vacuum Einstein equations.
(i). Suppose $`M`$ is pseudo-compact.
Then $`\overline{M}U`$ has only finitely many ends $`\{E_i\},0iq<\mathrm{}.`$ Supposing $`q>`$ 0 and (i) holds, let $`E\{E_i\}`$ be any end of $`M`$ satisfying
(ii). $`u(x_i)u_o,`$ for some constant $`u_o>`$ 0 and some sequence $`x_iE`$ with $`t(x_i)\mathrm{}.`$
Then the end $`(E,g,u)`$ is either asymptotically flat, or small, in the sense that the area growth of geodesic spheres satisfies
(0.8)
$$^{\mathrm{}}\frac{1}{areaS_E(s)}𝑑s=\mathrm{}.$$
Further, if $`m_E`$ 0 and $`sup_Eu<\mathrm{},`$ then the end $`E`$ is asymptotically flat.
We make several remarks on this result. First, if $`M`$ is not pseudo-compact, then it is easy to construct static vacuum solutions for which $`\overline{M}U`$ has infinitely many ends, c.f. the end of §1 or §2(I). Further, there are examples of ends $`E`$ with compact boundary, on which (ii) does not hold and which are neither asymptotically flat nor small in the sense of (0.8), c.f. Example 2.11. Thus, both hypotheses (i) and (ii) are necessary in Theorem 0.3. There are also examples with $`\overline{M}`$ compact, c.f. §2(I), so that $`\overline{M}U`$ may have no ends, ($`q=0`$).
On the other hand, both alternatives in Theorem 0.3, namely asymptotically flat or small ends, do occur. Asymptotically flat ends satisfy $`areaS_E(s)4\pi s^2`$ as $`s\mathrm{}`$, while small ends have small area growth. For example, (0.8) implies, at least, that there is a sequence $`s_i\mathrm{}`$ such that
$$areaS_E(s_i)s_i(logs_i)^{1+\epsilon },$$
for any fixed $`\epsilon >`$ 0. All known examples of solutions satisfying (0.8) are topologically of the form $`(^2B)\times S^1`$ outside a compact set and the geodesic spheres have at most linear area growth. The main example of a static vacuum solution with a small end is the family of Kasner metrics, c.f. Example 2.11 below.
To illustrate the sharpness of the last statement in Theorem 0.3, we construct in Remark 3.8 a (dipole-type) static vacuum solution $`(M,g,u)`$ with a single end $`E`$ on which $`m_E=`$ 0, $`sup_E<\mathrm{},`$ and which satisfies (0.8). Similarly, Example 2.11 provides static vacuum solutions with $`m_E>0`$, $`sup_Eu=\mathrm{}`$ and satisfying (0.8). Thus, the last result is also sharp.
Assuming $`M`$ is pseudo-compact, it is easy to see that if $`m_E>0`$ then (ii) holds, so that (ii) may be replaced by the assumption $`m_E>0`$. Thus, for the physically very reasonable class of solutions such that $`M`$ is pseudo-compact, $`m_E>0`$ for all ends $`E`$, and $`u`$ is bounded, all ends $`E`$ of $`M`$ are asymptotically flat. We also point out that if $`M`$ is complete away from $`\mathrm{\Sigma }`$, then (ii) holds at least on some end $`E`$.
The proof of Theorem 0.3 gives some further information on the asymptotic structure of the small ends. For instance, the curvature in such ends decays at least quadratically, c.f. (1.3), and in the geodesic annuli $`A_E(\frac{1}{2}s_i,2s_i)`$ in $`E`$, the metric approaches in a natural sense that of a Weyl solution, i.e. a static axisymmetric solution, as $`s_i\mathrm{}`$. Thus, asymptotically, small ends have at least one non-trivial Killing field. However, it is not known for example if the metric is asymptotic to a unique Weyl solution, or even if small ends are necessarily of finite topological type.
It would also be of interest to prove that $`(M,g,u)`$ has a unique end if $`M`$ is pseudo-compact. However, we have not been able to do this without further assumptions, c.f. Remark 3.10.
The definition of asymptotically flat rules out the possibility that $`M`$ is smooth up to $`\mathrm{\Sigma }`$ and complete away from $`\mathrm{\Sigma },`$ with $`\mathrm{\Sigma }`$ non-compact, since for example $`u=0`$ on $`\mathrm{\Sigma }`$ while $`u1`$ on any asymptotically flat end. Consider however the following (B1) solution,
(0.9)
$$g_{B1}=(1\frac{2m}{r})^1dr^2+(1\frac{2m}{r})d\varphi ^2+r^2d\theta ^2,$$
where $`\varphi /4m[0,2\pi ],\theta [0,\pi ],r2m>`$ 0, with $`u=rsin(\theta ).`$ Note that the potential $`u`$ is unbounded.
The B1 metric is just a 3-dimensional slice of the 4-dimensional Schwarzschild metric $`(N,g_S+u^2d\theta ^2),(t`$ changed to $`\theta `$ in (0.2)), obtained by dividing $`N`$ by an $`S^1Isom(S^2)`$ orthogonal to $`d\theta ;`$ this is the slice ’orthogonal’ to the usual slice giving the Schwarzschild metric (0.4). This metric has $`\mathrm{\Sigma }`$ given by two disjoint, isometric copies of $`^2,`$ each of positive Gauss curvature and asymptotic to a flat cylinder. It is smooth up to $`\mathrm{\Sigma }`$ and complete away from $`\mathrm{\Sigma }`$. The metric is globally asymptotic to the flat metric on $`^2\times S^1,`$ again with curvature decay of order $`O(1/r^3)`$ and with $`u`$ of linear growth in distance to $`\mathrm{\Sigma }.`$ Such solutions will be called asymptotically cylindrical. In fact a large class of Weyl solutions have ’dual‘ solutions in this sense which are asymptotically cylindrical, c.f. Remark 2.9.
This paper is organized as follows. Following discussion of some general topics on static space-times in §1, we analyse in some detail the class of Weyl vacuum solutions in §2. Several new results on the structure of these solutions are given; for instance Proposition 2.2 gives a new characterization of Weyl metrics. In addition, some efforts have been made to give a reasonably clear and organized account of the breadth of possibilities and behavior of Weyl metrics, since their treatment in the literature is rather sketchy and since they serve as a large class of models on which to test Theorem 0.3. Theorem 0.3 is proved in §3, and the paper concludes with several remarks on generalizations, and some open questions.
I would like to thank the referee for suggesting a number of improvements in the exposition of the paper.
## 1. Background Discussion.
Let $`(M,g,u)`$ be an open, connected oriented Riemannian 3-manifold and $`N=M\times _u`$ or $`N=M\times _uS^1,`$ as in (0.2). Thus $`N`$ represents a static space-time and $`MN`$ is totally geodesic. The Einstein field equations on $`N`$ are
(1.1)
$$r_N\frac{s_N}{2}g_N=T,$$
where $`T`$ is the $`(`$ or $`S^1`$-invariant) stress-energy tensor. (We are ignoring physical constants here). These equations may be expressed on the space-like hypersurface $`M`$ as the system
(1.2)
$$ru^1D^2u+(u^1\mathrm{\Delta }u\frac{1}{2}s)g=T_H,$$
$$\frac{1}{2}s=T_V,$$
where $`T_H`$ is the horizontal or space-like part of $`T`$ and $`T_V`$ is the vertical or time-like part of $`T`$. These are the equations on $`M`$ for a Lorentzian space-time $`N`$; in case $`N`$ is Riemannian, the first equation is the same while the second is -$`\frac{1}{2}s=T_V`$. When $`T=`$ 0, one obtains the vacuum equations (0.1), which are the same for Lorentzian or Riemannian signature. A common example, with $`T0`$ is a static perfect fluid, c.f. \[Wd, Ch.4\], with $`T`$ given by
$$T=(\mu +\rho )dt^2+\rho g,$$
where $`\mu ,\rho `$ are time independent scalar fields representing the energy density and pressure respectively. The equations (1.2) imply that the full Riemann curvature $`R_N`$ of $`N`$ is determined by $`r`$, $`u`$ and $`T`$.
The horizon $`\mathrm{\Sigma }=\{u=0\}`$ corresponds formally to the fixed point set of the $`S^1`$ action on $`N`$ and requires special consideration. For example, even if $`M`$ is smooth up to $`\mathrm{\Sigma }`$ the Riemannian 4-manifold $`(N,g_N)`$ might not be smooth across $`\mathrm{\Sigma },`$ even though the curvature $`R_N`$ is smooth. Namely, assuming the $`S^1`$ parameter $`t[0,2\pi )`$, if $`|u|_\mathrm{\Sigma }`$ 1, then $`N`$ has cone singularities (with constant angle by (0.3)) along and normal to the totally geodesic submanifold $`\mathrm{\Sigma }N.`$ (This issue does not arise for Lorentzian metrics). By multiplying the potential function $`u`$ by a suitable constant, one can make the metric $`g_N`$ smooth across any given component of $`\mathrm{\Sigma };`$ one cannot expect however in general that this can be done simultaneously for all components of $`\mathrm{\Sigma },`$ if there are more than one. This issue will reappear in §2.
The following result is proved in \[A1, Cor.A.3\]. It implies, (by letting $`t\mathrm{}),`$ that if $`(M,g,u)`$ is a complete solution to the static vacuum equations with $`u>`$ 0 everywhere, then $`M`$ is flat, and $`u`$ is constant.
###### Theorem 1.1.
Let (M, g, u) be a solution to static vacuum equations (0.1). Let t(x) $`=`$ dist(x, $`M)`$ as in (0.6). Then there is a constant $`K<\mathrm{},`$ independent of (M, g, u), such that
(1.3)
$$|r|(x)\frac{K}{t(x)^2},|logu|(x)\frac{K}{t(x)}.$$
###### Remark 1.2.
The same result has recently been proved for stationary vacuum solutions, i.e. space-times admitting a complete time-like Killing field, c.f. \[A2\]. We will discuss elsewhere to what extent Theorem 0.3 generalizes to stationary vacuum solutions.
Recall from §0 that the potential function $`u`$ of a static vacuum solution may be freely renormalized by arbitrary positive constants; hence the appearance of log $`u`$ in (1.3), as in (0.7).
Theorem 1.1 implies that the curvature of $`(M,g)`$ is controlled away from $`M,`$ and hence the local geometry of solutions is controlled away from $`M`$ by lower bounds on the local volume or injectivity radius. More precisely, we have the following results which are essentially a standard application of the Cheeger-Gromov theory of convergence/collapse of Riemannian manifolds, c.f. \[CG\], \[A3, §2\], or also \[P, Ch.12\] for an introduction to these results. Further details of the proofs of these results are given in \[A2\], (for the more general class of stationary space-times), and also in \[A1,App.\].
###### Lemma 1.3.
(Non-Collapse). Let $`(M_i,g_i,u_i)`$ be a sequence of solutions to the static vacuum equations (0.1). Suppose
$$|r_i|\mathrm{\Lambda },diamM_iD,volM_i\nu _o,$$
and
$$dist(x_i,M_i)\delta ,$$
for some $`x_iM_i`$ and positive constants $`\nu _o,\mathrm{\Lambda },`$ D, $`\delta .`$ Assume also that $`u_i`$ is normalized so that $`u_i(x_i)`$ = 1.
Then, for any $`ϵ>`$ 0 sufficiently small, there are domains $`U_iM_i,`$ with $`ϵ/2dist(U_i,M_i)ϵ,`$ and $`x_iU_i`$ such that a subsequence of the Riemannian manifolds $`(U_i,g_i)`$ converges, in the $`C^{\mathrm{}}`$ topology, modulo diffeomorphisms, to a limit manifold (U, g), with limit function $`u`$ and base point $`x=`$ lim $`x_i.`$ The triple (U, g, u) is a smooth solution of the static equations with u(x) $`=`$ 1.
###### Lemma 1.4.
(Collapse). Let $`(M_i,g_i,u_i)`$ be a sequence of solutions to the static vacuum equations (0.1). Suppose
$$|r_i|\mathrm{\Lambda },diamM_iD,volM_i0$$
and
$$dist(x_i,M_i)\delta ,$$
for some $`x_iM_i`$ and constants $`\mathrm{\Lambda },`$ D, $`\delta .`$ Assume also that $`u_i`$ is normalized so that $`u_i(x_i)`$ = 1.
Then, for any $`ϵ>`$ 0 sufficiently small, there are domains $`U_iM_i,`$ with $`ϵ/2dist(U_i,M_i)ϵ`$ with $`x_iU_i,`$ such that $`U_i`$ is either a Seifert fibered space or a torus bundle over an interval. In both cases, the $`g_i`$-diameter of any fiber F, (necessarily a circle $`S^1`$ or torus $`T^2),`$ goes to 0 as $`i\mathrm{},`$ and $`\pi _1(F)`$ injects in $`\pi _1(U_i).`$
Consequently, there are infinite $``$ or $``$ covers $`\stackrel{~}{U}_i`$ of $`U_i`$, such that $`\{\stackrel{~}{U}_i,g_i,x_i\}`$ does not collapse and hence has a subsequence converging smoothly to a limit $`(\stackrel{~}{U},g,x)`$ of the static vacuum equations with $`x=`$ lim $`x_i^{},x_i^{}`$ a lift of $`x_i`$ to $`\stackrel{~}{U}_i.`$ The limit $`(\stackrel{~}{U},g,x)`$ admits a free isometric $``$ or $``$ action, (c.f. §2), which also leaves the potential function $`u`$ invariant, and $`u(x)=`$ 1.
In studying static solutions, it is often very useful to consider the conformally equivalent metric $`\stackrel{~}{g}=u^2g`$ on $`M`$. An easy calculation using the behavior of Ricci curvature under conformal deformations, c.f. \[Bes, p.59\], shows that the Ricci curvature $`\stackrel{~}{r}`$ of $`\stackrel{~}{g},`$ in the vacuum case (0.1), is given by
(1.4)
$$\stackrel{~}{r}=2(dlogu)^20.$$
Further, if $`\stackrel{~}{\mathrm{\Delta }}`$ denotes the Laplacian of $`\stackrel{~}{g},`$ then
(1.5)
$$\stackrel{~}{\mathrm{\Delta }}logu=0.$$
The equations (1.4)-(1.5) are equivalent to the static vacuum equations (0.1). Since these equations are invariant under the substitution $`uu,`$ it follows that if $`(M,g,u)`$ is a static vacuum solution, then so is $`(M,g^{},u^1),`$ with $`g^{}`$ given by
$$g^{}=u^4g.$$
Similarly, observe that if $`(N,g_N)`$ is the associated Ricci-flat 4-manifold (0.2), then
(1.6)
$$\mathrm{\Delta }_Nlogu=0.$$
Here and below, log always denotes the natural logarithm.
We discuss briefly some of the simplest static vacuum solutions:
Levi-Civita Solutions.
There are 7 classes of so-called degenerate static vacuum solutions, where the eigenvalues $`\lambda _i`$ of the Ricci curvature $`r`$ satisfy $`\lambda _1=\lambda _2=2\lambda _3,`$ called A1-A3, B1-B3, C, c.f. \[EK, §2-3.6\]. The $`B`$ metrics are dual to the $`A`$ metrics, as mentioned in §0, c.f. §2 for details. The A1 metric is the Schwarzschild metric. It is of interest to examine the A2 metric, given in standard cylindrical coordinates on $`^3`$ by
(1.7)
$$g_{A2}=z^2(dr^2+(sinh^2r)d\varphi ^2)+(\frac{2m}{z}1)^1dz^2,$$
with $`u=(\frac{2m}{z}1)^{1/2}`$ and $`z[0,2m],m>0`$. The horizon $`\mathrm{\Sigma }=\{u=0\}`$ is given by the set $`\{z=2m\}`$ and hence is the complete hyperbolic metric $`H^2,`$ with curvature $`(2m)^2.`$ It is easily verified that the A2 metric is smooth up to $`\mathrm{\Sigma }.`$ However, $`\mathrm{\Sigma }M;`$ the set $`\{z=`$ 0} is at finite distance to $`\mathrm{\Sigma },`$ and so $`M`$ has another (singular) component obtained by crushing (compact subsets of) the hyperbolic metric to a point.
Let $`\mathrm{\Gamma }`$ be any discrete group of hyperbolic isometries. Then $`\mathrm{\Gamma }`$ extends in an obvious way to a group of isometries of $`g_{A2}.`$ The uniformization theorem for surfaces implies that any orientable surface except $`S^2`$ and $`T^2,`$ including surfaces of infinite topological type and infinitely many ends, admits a complete hyperbolic metric, i.e. is the quotient $`H^2/\mathrm{\Gamma },`$ for some $`\mathrm{\Gamma }.`$ Hence, any such surface and hyperbolic metric can be realized as the horizon $`\mathrm{\Sigma }`$ of a static vacuum solution.
Topologically, for $`\mathrm{\Sigma }=H^2/\mathrm{\Gamma }`$, we have $`M`$ = $`\mathrm{\Sigma }\times I`$. Hence for example if $`\mathrm{\Sigma }`$ has infinitely many ends, then $`M`$ also has infinitely many ends; in particular, this shows that the hypothesis that $`M`$ is pseudo-compact in Theorem 0.3 is necessary.
The A3 metric is
(1.8)
$$g_{A3}=z^2(dr^2+r^2d\varphi ^2)+zdz^2,$$
with $`u=z^{1/2}>0,r>0.`$ Hence $`\mathrm{\Sigma }`$ is empty in this case - it occurs at infinity in the metric. This metric may be realized as a pointed limit of the A2 metric as $`m\mathrm{}`$ and also as a limit, in a certain sense, of the A1 metric, c.f. Example 2.11.
###### Remark 1.5.
The discussion above raises the natural question if any orientable connected Riemannian surface $`(\mathrm{\Sigma },g)`$ can be realized as the horizon of a static vacuum solution, smooth up to $`\mathrm{\Sigma },`$ which is defined at least in a neighborhood of $`\mathrm{\Sigma }.`$ In general, this appears to be unknown.
Observe that any complete constant curvature metric on an orientable surface can be realized in this way. The Schwarzschild metric gives the constant curvature metric on $`S^2,`$ the quotients of the A2 metric give all hyperbolic surfaces, and quotients of the flat metric, with $`u`$ a linear function, give all flat metrics on a surface, $`(T^2,S^1\times ,\mathrm{or}^2).`$
Geroch-Hartle show in \[GH\] that any rotation-invariant metric on $`S^2`$ or $`T^2`$ can be realized at the horizon. Except for the Schwarzschild metric, such solutions are not complete away from $`\mathrm{\Sigma }.`$
Observe that the full 1-jet of $`(M,g)`$ at $`\mathrm{\Sigma }`$ (assumed connected) is determined solely by the surface metric $`(\mathrm{\Sigma },g)`$, since $`\mathrm{\Sigma }`$ is totally geodesic and, renormalizing $`u`$ if necessary, $`|u|`$ 1 on $`\mathrm{\Sigma }`$ by (0.3). Observe also that one cannot have $`|u|`$ 0 on $`\mathrm{\Sigma },`$ since $`u`$ is harmonic and the divergence theorem applied to a small neighborhood $`U`$ of $`\mathrm{\Sigma }`$ would imply that $`u`$ 0 on $`U`$, which is ruled out.
On the other hand, the metric $`(M,g)`$ is not uniquely determined by its boundary values $`(\mathrm{\Sigma },g)`$. Namely, the flat metric on $`T^2`$ is realized by the flat vacuum solution $`M=T^2\times ^+,u=t=dist(\mathrm{\Sigma },)`$ and also (locally) by a non-flat metric, c.f. \[T\],\[P\]. Similar remarks hold for local perturbations of the Schwarzschild metric, c.f. \[GH\].
## 2. Weyl Solutions.
A large and very interesting class of explicit solutions of the static vacuum equations are given by the Weyl solutions \[W\], c.f. also \[EK,§2.3-9\] or \[Kr, Ch.16-18\]. In fact, it appears that essentially all known explicit solutions of the static vacuum equations are of this form. Since the literature on these solutions is not very organized or rigorous, especially regarding their global structure, we discuss these solutions in some detail. These metrics will also illustrate the necessity of the hypotheses in Theorem 0.3.
###### Definition 2.1.
A Weyl solution is a solution $`(M,g,u)`$ of the static vacuum equations (0.1) which admits an isometric $``$-action, i.e. a non-zero homomorphism $`Isom(M)`$, leaving $`u`$ invariant.
Apriori, the topology of a Weyl solution could be quite non-trivial; for example $`M`$ could be any Seifert fibered space. The first result shows that only the simplest topology (and geometry) is possible. For the moment, we exclude any possible fixed point set of the $``$-action from the discussion.
###### Proposition 2.2.
Let $`(M,g,u)`$ be a Weyl solution with $``$-action without fixed points, which does not admit a (local) free isometric $`\times `$ action. Then the universal cover $`(\stackrel{~}{M},g)`$ of $`(M,g)`$ is a warped product of the form
(2.1)
$$\stackrel{~}{M}=V\times _f,g=g_V+f^2d\varphi ^2,$$
with (V, $`g_V)`$ a Riemannian surface and $`f`$ a positive function on V. The $``$-action on $`\stackrel{~}{M}`$ is by translation on the second factor.
Proof: This result, whose proof is purely local, is a strengthening in this situation of a well-known result in general relativity, Papapetrou’s theorem, c.f. \[Wd,Thm.7.1.1\], which requires certain global assumptions, (e.g. smoothness up to $`\mathrm{\Sigma }).`$
Let $`K`$ denote the (complete) Killing field generated by the $``$-action on $`M`$, and $`f=|K|`$. We may assume that $`u`$ is not a constant function on $`M`$, since if $`u`$ is constant, the metric is flat, and so admits a local $`\times `$ action. Since $`(M,g,u)`$ is real-analytic, $`u`$ is not constant on any open set in $`M`$. We thus choose a neighborhood $`U`$ of any point $`p`$ where $`u(p)0`$ on which $`|u|>0`$. Define $`e_1`$ by $`e_1=u/|u|`$, and extend it to a local orthonormal frame $`e_1,e_2,e_3`$ for which $`e_3=K/|K|=K/f`$. Note that this is possible since $`u`$ is required to be invariant under the flow of $`K`$, so that
(2.2)
$$<u,K>=0.$$
In $`U`$, the metric $`g`$ may be written as
(2.3)
$$g=\pi ^{}g_V+f^2(d\varphi +\theta )^2,$$
where $`\pi :UV`$ is a Riemannian submersion onto a local surface $`(V,g_V)`$, $`\theta `$ is a connection 1-form, $`K=/\varphi `$ and $`f`$ is a function on the orbit space $`V`$. If $`\theta =`$ 0, then the result follows. Thus, we assume $`|\theta |>`$ 0 in $`U`$ and show this implies that $`g`$ has a free isometric local $`\times `$ action.
Consider the 1-parameter family of metrics
(2.4)
$$g_s=\pi ^{}g_V+s^2f^2(d\varphi +s^2\theta )^2,$$
for $`s>`$ 0, with $`g_1=g`$. Geometrically, this corresponds to rescaling the length of the fibers of $`\pi `$ and changing the horizontal distribution of $`\pi ,`$ (when $`\theta `$ 0). Now it is a standard fact that the 1-parameter family of 4-metrics
(2.5)
$$g_s^4=g_s\pm u^2dt^2$$
remains Ricci-flat for all $`s`$. This can be seen from standard formulas for Riemannian line bundles, c.f. \[Bes, 9.36, 9G\], \[Kr, 16.1-3\] or \[A2, §1.2\]. Thus, the metrics $`g_s`$ all satisfy the static vacuum equations
$$ur_s=D_s^2u$$
from (0.1), with the same potential $`u`$. Equivalently, the conformal metrics $`\stackrel{~}{g}_s=u^2g_s`$ satisfy
(2.6)
$$\stackrel{~}{r}_s=2(dlogu)^2,$$
for all $`s`$, c.f. (1.4) The right side of (2.6) is of course independent of $`s`$.
We claim that the metrics $`g_s`$ on $`U`$ are all locally isometric. While this could be proved by a lengthy direct computation, we argue more conceptually as follows. Let $`e_i^s`$ be a local orthonormal frame for $`g_s,`$ determined as above for $`g`$. We then have $`e_1^s=e_1,e_3^s=s^1e_3`$ while $`e_2^s`$ varies in the plane $`<e_2,e_3>.`$ The same relations hold w.r.t. $`\stackrel{~}{g}_s.`$ Recall also that the full curvature tensor is determined by the Ricci curvature in dimension 3. It then follows from these remarks and (2.6) that for each $`qU`$ and $`s>`$ 0, there is a sectional curvature preserving isomorphism $`F_s:T_qMT_qM,`$ i.e.
$$\stackrel{~}{K}_s(F_s(P))=\stackrel{~}{K}_1(P),$$
where $`P`$ is any 2-plane and $`\stackrel{~}{K}_s`$ is the sectional curvature w.r.t. $`\stackrel{~}{g}_s.`$ Clearly $`F_s`$ varies smoothly with $`q`$ and $`s`$. Using the expression (2.6), a result of Kulkarni \[Ku\] then implies that the metrics $`\stackrel{~}{g}_s`$ are locally isometric and hence so are the metrics $`g_s.`$
Let $`\mathrm{\Omega }=d\theta `$ be the curvature form of the line bundle $`\pi .`$ Then w.r.t. the metric $`g`$, $`|\mathrm{\Omega }|=|\mathrm{\Omega }(e_1,e_2)|=|<_{e_1}e_2,e_3>|.`$ The same equalities hold w.r.t. $`\mathrm{\Omega }_s=d\theta _s=s^2d\theta `$ and the $`g_s`$ metric. A short computation then gives
(2.7)
$$|\mathrm{\Omega }_s|_{g_s}0,\mathrm{as}s\mathrm{}.$$
Hence consider the behavior of the metrics $`g_s`$ as $`s\mathrm{}.`$ We are then expanding or blowing up the metric in the fiber direction, at a given base point. Since the metrics $`g_s`$ are isometric, there are (local) diffeomorphisms $`\psi _s`$ such that $`\psi _s^{}g_s`$ converges to a limit metric $`g_{\mathrm{}}.`$ At a given base point, the diffeomorphisms $`\psi _s`$ expand or blow up smaller and smaller intervals of the parameter $`\varphi `$ to unit size, giving rise to a limit parameter $`\varphi _{\mathrm{}}.`$ When $`\theta =`$ 0, this is the only change; the limit metric $`g_{\mathrm{}}`$ is the same as $`g`$ with the parameter $`\varphi `$ replaced by $`\varphi _{\mathrm{}}.`$ (This is completely analogous to passing from the flat metric on $`^2\{0\}`$ to the flat metric on its universal cover, i.e. unwrapping the circles to lines).
However, when $`\theta `$ 0, the $`e_2`$ direction is also being expanded or blown-up in a similar way. The function $`u`$ is left invariant under the family $`\{\psi _s\}.`$
It follows from (2.7) that the limit metric $`g_{\mathrm{}}`$ is a static vacuum solution of the form
(2.8)
$$g_{\mathrm{}}=\pi ^{}g_V_{\mathrm{}}+f_{\mathrm{}}^2(d\varphi _{\mathrm{}})^2,$$
i.e. the 1-form $`\theta _{\mathrm{}}=`$ 0 in the limit. Further, since the $`e_2`$ direction has been blown up, the function $`f_{\mathrm{}}`$ varies only in the $`e_1`$ direction, i.e. $`f_{\mathrm{}}=f_{\mathrm{}}(u).`$
Metrics of the form (2.8) are analysed in detail below. Referring to (2.10), let $`r=f_{\mathrm{}}u=h(u)`$. In a possibly smaller open subset of $`U`$, we may invert $`h`$ and write $`u=u(r)`$, where $`r`$ is a local coordinate on $`V_{\mathrm{}}.`$ It is easy to see, (c.f. (2.12)-(2.13) below for example), that $`g_{\mathrm{}}`$ admits a non-vanishing Killing field $`/z,`$ tangent to $`V_{\mathrm{}}`$ but orthogonal to $`/r,`$ and hence $`g_{\mathrm{}}`$ admits a free isometric local $`\times `$ action. The metrics $`g_s`$ are all locally isometric and so the metric $`g=g_1`$ also has a free isometric local $`\times `$ action.
Since the proof above is completely local, Proposition 2.2 holds locally, (in suitably modified form), even if $`(M,g)`$ admits only a local or partial $``$ -action.
Static vacuum solutions admitting a free isometric local $`\times `$ action are completely classified; they are either flat or belong to the family of Kasner metrics, c.f. Example 2.11 below or \[EK, Thm.2-3.12\]. Such solutions do have Killing fields $`K`$ which are not hypersurface orthogonal, i.e. $`d\theta `$ 0 in (2.3). For example, if $`/\psi `$ and $`/z`$ are standard generators of the (local) $`\times `$ action, then linear combinations such as $`K=/\psi +/z`$ are non-hypersurface orthogonal Killing fields.
Nevertheless, all such solutions do admit, of course, hypersurface orthogonal Killing fields and so may be written in the form (2.1). For the remainder of the paper, we thus assume that a Weyl solution has the form (2.1) locally. In addition, we will always work with the $``$-quotient of the metric (2.1) and so consider Weyl solutions as warped products of the form $`V\times _fS^1`$; it will not be assumed in general that $`V`$ is simply connected.
Duality. Observe that Weyl solutions naturally come in ’dual’ pairs. Namely the Ricci-flat 4-manifold $`(N,g_N)`$ has the form
(2.9)
$$g_N=g_V+f^2d\varphi ^2+u^2dt^2,$$
and so both $`M_u=V\times _fS^1`$ and $`M_f=V\times _uS^1`$ are static vacuum solutions on the 3-manifolds, with potentials $`u`$, resp. $`f`$. Consider the product of the lengths of circles, or equivalently, the area of the torus fiber in $`N`$,
(2.10)
$$r=fu.$$
This is a globally defined positive harmonic function on $`(V,g_V).`$ To see this, on $`M=M_u,`$ by (2.5), we have
$$0=\mathrm{\Delta }u=\mathrm{\Delta }_Vu+<logf,u>,$$
and the same formula, with $`u`$ and $`f`$ reversed, holds on $`M_f.`$ Hence
$$\mathrm{\Delta }_Vfu=f\mathrm{\Delta }_Vu+u\mathrm{\Delta }_Vf+2<u,f>=0.$$
Charts. We now describe a collection of preferred charts in which to express the Weyl solution $`(M,g)`$; this description is due to Weyl \[Wl\]. The surface $`V`$ may be partitioned into a collection of maximal domains $`V_i`$ on which the harmonic conjugate $`z`$ of $`r`$ is single-valued, so that $`F=r+iz`$ is a well-defined holomorphic function from $`V_i`$ into the right half-plane $`^+=`$ {$`(r,z)`$: $`r>`$ 0, $`z\}.`$ One might also pass to a suitable cover, for instance the universal cover, of $`V`$ to obtain a globally defined conjugate harmonic function, but it is preferable not to do so.
Now each $`V_i`$ may be further partitioned into a collection of domains $`V_{ij}`$ on which $`F`$ is a conformal embedding into $`^+,`$ so that $`g|_{V_{ij}}=F^{}(dr^2+dz^2).`$ We will thus simply view $`V_{ij}`$ as a domain in $`^+,`$ with $`g_V`$ a metric pointwise conformal to the flat metric $`dr^2+dz^2.`$
It follows that the corresponding domain $`M_{ij}=V_{ij}\times _fS^1`$ is embedded as a domain $`\mathrm{\Omega }=\mathrm{\Omega }_{ij}`$ in $`^3`$ endowed with cylindrical coordinates $`(r,z,\varphi )`$, $`\varphi [0,2\pi )`$ with the background (unphysical) complete flat metric $`dr^2+dz^2+r^2d\varphi ^2.`$ We note that all the data above are canonically determined by the two Killing fields on $`N`$ and thus the coordinates $`(r,z,\varphi )`$ are called canonical cylindrical or Weyl coordinates for $`(M,g)`$. Of course $`\mathrm{\Omega }`$ is axially symmetric, i.e. symmetric w.r.t. rotation about the $`z`$-axis.
To express the metric $`g|_\mathrm{\Omega }`$ in these coordinates, the field equations (0.1) imply that the function
(2.11)
$$\nu =logu$$
is an axially symmetric (independent of $`\varphi )`$ harmonic function on $`\mathrm{\Omega }^3;`$ this again follows in a straightforward way from computation of the Laplacian of $`u`$ and $`f`$ on $`M_f`$ and $`M_u`$ as above. A computation of the conformal factor for the metric $`g_V,`$ c.f. \[Wd,Ch.7.1\], then leads to the expression of $`g`$ in these coordinates:
(2.12)
$$g=u^2(e^{2\lambda }(dr^2+dz^2)+r^2d\varphi ^2),$$
where $`\lambda `$ is determined by $`\nu `$ as a solution to the integrability equations
(2.13)
$$\lambda _r=r(\nu _r^2\nu _z^2),\lambda _z=2r\nu _r\nu _z.$$
The equations (2.13) mean that the 1-form $`\omega =r(\nu _r^2\nu _z^2)dr+2r\nu _r\nu _zdz`$ is closed on $`\mathrm{\Omega }.`$
Conversely, given any axially symmetric harmonic function $`\nu `$ on a connected open set $`\mathrm{\Omega }`$ in $`^3,`$ if the closed 1-form $`\omega `$ is exact, (for example if $`\pi _1(\mathrm{\Omega }^+)=0),`$ the equations (2.13) determine $`\lambda `$ up to a constant and the metric (2.12) gives a solution to the static vacuum equations with $`S^1`$ symmetry. (The addition of constants to $`\nu `$ or $`\lambda `$ changes the metric at most by diffeomorphism or homothety).
It is remarkable that solutions to the non-linear vacuum equations (0.1) can be generated in this way by solutions to the linear Laplace equation on $`^3.`$
###### Remark 2.3.
(i). The Levi-Civita solutions in §1 are all Weyl solutions. However, the expressions for the A1-A3 metrics in (0.4),(1.7),(1.8) and the B1 metric in (0.9) are not in Weyl canonical coordinates. Note that the quotients of the A2 metric discussed above are no longer Weyl solutions, although they could be considered as local Weyl solutions; the $``$-action is only locally defined on the quotients.
(ii). There seem to be no known Weyl solutions which can not be expressed globally in the form (2.12).
Fixed Point Set. The behavior of solutions at the part of $`M`$ where either one of the two $`S^1`$ or $``$ actions on $`N`$ has fixed points requires special considerations. This is of course the locus where $`u=`$ 0 or $`f=`$ 0, and hence includes the part $`\overline{\mathrm{\Omega }}A`$ of the $`z`$-axis $`A=\{r=0\}`$ in any canonical coordinate chart. It is not necessarily the case however that this locus is contained in $`A`$, c.f. the end of Example 2.10.
Note that given any Weyl solution (2.12), any covering of $`^3A`$ induces another solution of the form (2.12), but with $`\varphi `$ parametrizing a circle of length $`2\pi k.`$ For the universal cover $`(k=\mathrm{}),`$ the $`\varphi `$-circle is replaced by a line. In fact, (2.12) is well-defined when $`\varphi `$ runs over any parameter interval \[0, $`2\pi \alpha ).`$ Observe however that any asymptotically flat Weyl solution must have $`\alpha =`$ 1, since the metric must be smooth near infinity. Thus, we will assume $`\alpha =`$ 1 in the following, unless stated otherwise.
Now suppose there is an open interval $`J`$ in $`A`$ such that the functions $`u`$ and $`\lambda ,`$ and hence the form (2.12) extend continuously to $`J`$. The form $`g`$ then represents a continuous metric in a neighborhood of $`J`$ if and only if the elementary flatness condition
(2.14)
$$\lambda =0,$$
is satisfied on $`J`$. On intervals where (2.14) does not hold, the metric $`g`$ has cone singularities, so that it is not locally Euclidean. From (2.13), it is clear that if $`\lambda `$ has a $`C^1`$ extension to $`J`$, then $`\lambda `$ is constant on $`J`$. However, such constants may vary over differing components of $`\overline{\mathrm{\Omega }}A.`$ This will be analysed further in Remark 2.8.
For the remainder of §2, we assume that
$$M=\mathrm{\Omega },$$
so that the Weyl solution is given globally in the form (2.12). Let $`I`$ be the set where $`\nu =\mathrm{},`$ i.e. the $`G_\delta `$ set in $`\mathrm{\Omega }^3`$ given by
(2.15)
$$I=\underset{n}{}\nu ^1(\mathrm{},n),$$
where $`n`$ runs over negative integers. It is usually assumed in physics that $`I`$ is non-empty, although this need not be the case; this will also be discussed further below. The set $`I`$ corresponds to the horizon $`\mathrm{\Sigma }`$ of the Weyl solution $`(M,g)`$, since $`u=`$ 0 on $`I`$. This correspondence is formal however, since the geometry and topology of $`I^3`$ is very different than that of $`\mathrm{\Sigma }(\overline{M},g)`$, c.f. most of the examples below. For the same reasons, although $`M=\mathrm{\Omega }`$ topologically, the metric boundary $`M`$ of $`(M,g)`$ is (most always) very different than the Euclidean boundary of $`\mathrm{\Omega }^3.`$
In the following, we discuss some of the most significant possible behaviors for the potential function $`u`$, and the associated Weyl solution, in order to illustrate the breadth of these solutions. The discussion is by no means complete or exhaustive.
(I). $`\overline{\mathrm{\Omega }}`$ compact.
Let $`M=\mathrm{\Omega }`$ be any bounded, $`C^{\mathrm{}}`$ smooth axisymmetric domain (i.e. connected open set with smooth compact boundary) in $`^3`$ and let $`\varphi `$ be any $`C^{k,\alpha }`$ function on $`\mathrm{\Omega },k`$ 0, $`\alpha `$ (0,1), which is axially symmetric about the $`z`$-axis. For simplicity, assume that $`\mathrm{\Omega }^+`$ is simply connected. Let $`\nu `$ be the solution to the Dirichlet problem
$$\mathrm{\Delta }\nu =0,\nu |_\mathrm{\Omega }=\varphi .$$
Then $`\nu `$ is also axi-symmetric about the $`z`$-axis, and hence $`\nu `$ generates a Weyl solution as in (2.12).
Suppose that for a given $`k`$ 1 $`\alpha `$ (0,1), $`\varphi `$ as above is $`C^{k,\alpha }`$ on $`\mathrm{\Omega },`$ but is nowhere $`C^{k+1}`$ on $`\mathrm{\Omega }.`$ Then $`\nu `$ extends to a $`C^{k,\alpha }`$ function on the Euclidean closure $`\overline{\mathrm{\Omega }}`$ and hence, from (2.13), the function $`\lambda `$ in (2.12) is also uniformly bounded. This means that the metric $`g`$ is quasi-isometric to the flat metric on $`\mathrm{\Omega },`$ and hence the metric boundary of $`\mathrm{\Omega }`$ w.r.t. the Weyl metric $`g`$ is the same as its Euclidean boundary. Since $`\nu `$ is not $`C^{k+1}`$ anywhere on $`\mathrm{\Omega },`$ this solution $`M=\mathrm{\Omega }`$ is maximal, i.e. admits no larger static vacuum extension; $`C^2`$ smooth solutions of the static vacuum equations are analytic. Further $`\nu `$ is bounded, so that $`u=e^\nu `$ is bounded away from 0, and hence the solution $`(\mathrm{\Omega },g)`$ has no horizon.
As noted in §0, the presence of the boundary $`M`$ is physically assumed due to the presence of matter or field sources. Thus, at least when $`k`$ 2, the vacuum solution $`(M,g)`$ can be extended to a larger space-like domain $`(M^{},g^{})(M,g)`$ with non-zero stress-energy $`T`$ in $`M^{}M.`$
On the other hand, if $`k=`$ 0 above, then the geometry of the metric boundary $`(M,g)`$ will in general be very different than the smooth geometry of $`\mathrm{\Omega }`$ in $`^3.`$ Further, one can of course consider non-smooth domains $`\mathrm{\Omega }^3`$ in this situation. These remarks indicate that the structure of the metric boundary $`M`$ seemingly can be quite arbitrary.
For the remainder of this section, we assume that $`\overline{\mathrm{\Omega }}`$ is non-compact in $`^3.`$ The same remarks as above hold for non-compact domains with smooth (surface) boundary. Thus for example $`M=\mathrm{\Omega }`$ might have infinitely many ends if $`\mathrm{\Omega }`$ is non-compact, showing that the assumption (i) in Theorem 0.3 is necessary. For simplicity, we only consider the following situation from now on.
(II). Suppose
(2.16)
$$dim_{}\mathrm{\Omega }1,$$
where the boundary is in the topology of $`^3`$ and $`dim_{}`$ is the Hausdorff dimension. Thus $`\mathrm{\Omega }`$ is a closed set of capacity 0, c.f. \[H, Thm.5.14\], and so in particular is a polar set. Clearly $`\overline{\mathrm{\Omega }}=^3.`$
(A).(Positive Case). Suppose that $`\nu `$ is locally bounded above, i.e.
(2.17)
$$\underset{B_x(1)}{sup}\nu <\mathrm{},x\mathrm{\Omega }.$$
It follows, c.f. \[H, Thm.5.18\] that $`\nu `$ extends uniquely to a globally defined subharmonic function on $`^3.`$ Hence, one may use the value distribution theory of subharmonic functions on $`^3`$ to analyse the geometry of Weyl solutions.
The Riesz representation theorem c.f. \[H,Thm.3.9\], implies that any subharmonic function $`\nu `$ on $`^3`$ may be represented semi-globally, i.e. on $`B(R)=B_0(R)^3`$ for any $`R<\mathrm{}`$, as
(2.18)
$$\nu (x)=_{B(R)}\frac{1}{|x\xi |}𝑑\mu _\xi +h(x),$$
where $`d\mu _\xi `$ is a positive Radon measure on $`B(R)`$ called the Riesz measure of $`\nu `$ and $`h`$ is a harmonic function on $`B(R)`$; both $`d\mu `$ and $`h`$ are axi-symmetric if $`\nu `$ is. (A Radon measure is a Borel measure which is finite on compact subsets). For the moment, we only consider the situation where there exists $`K<\mathrm{},`$ independent of $`R`$, s.t.
(2.19)
$$_{B(R)}\frac{1}{|x\xi |}𝑑\mu _\xi K.$$
In this case, one obtains a global representation of $`\nu `$ as
(2.20)
$$\nu (x)=_^3\frac{1}{|x\xi |}𝑑\mu _\xi +h(x),$$
where $`d\mu _\xi `$ is a positive measure and $`h`$ a harmonic function on $`^3.`$ (In (D) below, we briefly discuss the situation where (2.19) is not assumed). In particular, if $`\nu `$ is uniformly bounded above, say sup $`\nu =`$ 0, then the Liouville theorem for harmonic functions implies that $`h`$ 0, and one has the expression
(2.21)
$$\nu (x)=_^3\frac{1}{|x\xi |}𝑑\mu _\xi .$$
Note that since $`\nu `$ is harmonic on $`\mathrm{\Omega }`$,
(2.22)
$$\overline{I}suppd\mu \mathrm{\Omega },$$
but in many situations, as will be seen below, the first inclusion is strict.
(A)(i). Pure harmonic potentials.
Suppose that $`\nu `$ is a smooth harmonic function defined on all of $`^3,`$ so that $`\nu =h`$ in (2.20). It is clear that in this case $`\mathrm{\Sigma }=\mathrm{}`$ in the sense that $`I=\mathrm{}`$ in $`^3.`$ Since $`\nu `$ is axisymmetric, $`\nu `$ may be viewed as an expansion in Legendre polynomials, i.e.
$$\nu =\underset{k0}{}a_kR^kP_k(\frac{z}{R}),$$
where $`R^2=r^2+z^2.`$ For instance, $`RP_1(\frac{z}{R})=z`$, $`R^2P_2(\frac{z}{R})=3z^2R^2.`$
While these solutions are defined on all of $`^3,`$ no such solution gives a complete Weyl metric $`g`$ on $`M=^3,`$ by Theorem 1.1. For instance, for $`\nu =z`$, the Weyl metric is
$$g=e^{2z}(e^{r^2}(dr^2+dz^2)+r^2d\varphi ^2)),u=e^z.$$
Any straight ray in the $`(r,z)`$ half-plane has finite length in this metric, except a ray parallel to the negative $`z`$-axis. The horizon $`\mathrm{\Sigma }`$ occurs formally at $`\{z=\mathrm{}\},`$ of infinite $`g`$-distance to any point in $`^3`$.
(A)(ii). Newtonian potentials.
Suppose that $`h=0`$ in (2.20), so that $`\nu `$ is the Newtonian potential of an axisymmetric positive mass distribution $`d\mu `$ as in (2.21). This situation corresponds exactly to the Newtonian theory of gravity, (or equivalently the electrostatics of a positively charged distribution). While there is a vast classical literature on this subject, we will only consider the most interesting situation where
(2.23)
$$suppd\mu =\overline{I},$$
so that $`\nu `$ approaches $`\mathrm{}`$ on a dense set in supp $`d\mu .`$
The following Lemma characterizes this situation.
###### Lemma 2.4.
Let $`d\mu `$ be an axi-symmetric positive Radon measure on $`^3.`$ Then
(2.24)
$$suppd\mu =\overline{I}suppd\mu A.$$
Proof: Suppose first that supp $`d\mu `$ is not contained in $`A`$. Since $`d\mu `$ is axially symmetric, part of supp $`d\mu `$, namely the part not contained in $`A`$, is then given by a union of circles about the $`z`$-axis. Suppose first that there is a circle $`C`$ which is an isolated component of supp $`d\mu `$, so that $`d\mu |_C`$ is a multiple of Lebesgue measure on $`C`$. This case has been examined in \[Wl\],\[BW\], and we refer there for details. In particular in this case the potential $`\nu `$ is bounded below on and near $`C`$, and hence supp $`d\mu \overline{I}.`$ If $`C`$ is not isolated, then using (2.21), the same reasoning holds, since the measure $`d\mu `$ is then even less concentrated on the circles.
On the other hand, if supp $`d\mu A`$, then $`d\mu `$ is a positive Radon measure on $`A`$. Standard measure theory implies that the upper density of $`d\mu `$ w.r.t. Lebesgue measure $`dA`$ at $`aA,`$ i.e. $`lim\; sup_{r0}\frac{\mu (B_a(r))}{r},`$ is positive, for Lebesgue almost all $`asuppd\mu .`$ From the expression (2.21), it is clear that for any such $`a`$, $`\nu (x)\mathrm{}`$ as $`xa`$. This gives the converse. ∎
For the remainder of the discussion in (A), we assume (2.23) holds. From the theory of subharmonic functions on $`^3,`$ the set $`I`$ given by (2.15) may be an arbitrary $`G_\delta `$ set in $`A`$ $`^3,`$ i.e. a polar set in $`A`$. Since countable unions of polar sets are polar, note that $`I`$ is not necessarily closed in $`A^3.`$ (For example, let $`\{z_i\}A`$ be an increasing sequence converging to a limit point $`z`$, with say $`dist(z_i,z_{i+1})=i^2,`$ and let $`d\mu =2^i\delta _{z_i},`$ where $`\delta _{z_i}`$ is the Dirac measure based at $`z_i.`$ Then $`I=\{z_i\}`$ and supp $`d\mu =\{z_iz\}).`$
Given any $`x^3,`$ let $`m_x(r)`$ be the mass of the measure $`d\mu `$ in the ball $`B_x(r),`$ i.e.
(2.25)
$$m_x(r)=_{B_x(r)}𝑑\mu .$$
This is a non-negative increasing function on $`^+,`$ for any given $`x`$ and the limit
(2.26)
$$m=lim_r\mathrm{}m_x(r)>0,$$
is the total mass of $`d\mu .`$ This agrees, up to a universal constant factor, with the (ADM or Komar) mass in general relativity, when the latter is defined, and with (0.7) for solutions with pseudo-compact boundary. Note that one may have $`m=+\mathrm{}.`$
Lemma 2.4 and a standard result from potential theory, c.f. \[H, Thm.3.20\], characterize the possible Riesz measures satisfying (2.23).
###### Lemma 2.5.
A necessary and sufficient condition that a positive Radon measure $`d\mu `$ is the Riesz measure of an axi-symmetric subharmonic function $`\nu `$ on $`^3`$ with sup $`\nu =`$ 0 and supp $`d\mu =\overline{I}`$ is that supp $`d\mu `$ A, and, for any given $`xA`$,
(2.27)
$$_1^{\mathrm{}}\frac{m_x(r)}{r^2}𝑑r<\mathrm{}.$$
It is easy to see that a Weyl solution $`(M,g)`$ generated by a potential $`\nu `$ as in (2.21) for which supp $`d\mu =\overline{I}`$ is a compact subset of the axis $`A`$, is asymptotically flat, in the sense preceding Theorem 0.1. Further, the simplest or most natural surfaces enclosing any finite number of compact components of $`\overline{I},`$ and intersecting $`A`$ outside $`\overline{I},`$ are 2-spheres in $`M`$. Of course if $`suppd\mu A`$ is non-compact, then the solution cannot be asymptotically flat. A simple example is the solution generated by the measure
$$d\mu =\frac{1}{1+|\zeta |}dA_\zeta ,$$
where $`\zeta `$ parametrizes $`A`$ and $`dA`$ is Lebesgue measure on $`A`$. Observe also that such solutions do not have pseudo-compact boundary.
It is worthwhile to discuss some standard examples of Weyl solutions and their corresponding measures.
###### Example 2.6.
(i).(Curzon Solution). From the point of view of the Riesz measure, perhaps the simplest example is the measure $`d\mu `$ given by a multiple of the Dirac measure at some point on $`A`$, so that $`\nu =m/R,`$ $`R(x)=|x|`$, is a multiple of the Green’s function on $`^3.`$ This gives rise to the Curzon (or monopole) solution, c.f. \[Kr,(18.4)\],
(2.28)
$$g_C=e^{2m/R}[e^{m^2r^2/R^4}(dr^2+dz^2)+r^2d\varphi ^2],$$
with $`u=e^{m/R}.`$ Here $`\mathrm{\Omega }=^3\{0\}`$, $`\mathrm{\Omega }=`$ {0}, and it is often stated that $`g_C`$ has a point-like singularity (monopole) at the origin. However the geometry of $`M`$ is very different than that of a point. Namely the circles about the $`z`$-axis have length diverging to infinity as $`R`$ 0. Thus, small spheres $`R=ϵ`$ about {0} become very long in the $`\varphi `$ direction, and very short in the transverse $`\theta `$ direction, forming a very long, thin cigar. In particular, as a metric space, $`M=.`$ This is the first example where $`M`$ is non-compact but pseudo-compact. Of course $`M=\mathrm{\Sigma },`$ so that $`(M,g)`$ is complete away from $`\mathrm{\Sigma }.`$ A more detailed analysis of the Curzon singularity is given in \[SS\].
Note that one could not have solutions with both directions expanding at $`M,`$ so that area $`M=\mathrm{},`$ with $`M`$ pseudo-compact. This can be seen by use of minimal surface arguments, c.f. \[G\].
(ii). (Schwarzschild Solution). The Schwarzschild metric (0.4) is a Weyl metric, with measure $`d\mu =\frac{1}{2}dA`$ on $`[m,m]`$, where $`dA`$ is the standard Lebesgue measure on $`A`$. The resulting potential $`\nu `$ in (2.21) is the Newtonian potential of a rod on the z-axis with mass density $`\frac{1}{2},`$ given by
(2.29)
$$\nu _S=\frac{1}{2}log(\frac{R_++R_{}2m}{R_++R_{}+2m}),\mathrm{where}R_\pm ^2=r^2+(z\pm m)^2.$$
As mentioned before, the horizon $`\mathrm{\Sigma }`$ here is a smooth totally geodesic 2-sphere of radius $`2m`$ and $`M=\mathrm{\Sigma }.`$
The Weyl solution generated by the potential $`a\nu _S,`$ for $`\nu _S`$ as in (2.29) with $`a`$ $`>`$ 0 and $`a`$ $``$ 1, is not isometric or homothetic to the Schwarzschild metric. The associated Weyl metric is no longer smooth up to the horizon; in fact $`\mathrm{\Sigma }`$ is not even a 2-sphere unless $`a=1`$.
###### Remark 2.7.
More generally, consider any Weyl solution generated by a Riesz measure $`d\mu `$ satisfying (2.23). Observe that $`f=\frac{r}{u},`$ the length of the $`\varphi `$ circles in the Weyl metric, stays bounded away from 0 and $`\mathrm{}`$ on approach to supp $`d\mu ,`$ if and only if
(2.30)
$$logrC\nu logr+C,$$
for some $`C<\mathrm{},`$ since $`\nu =logu`$. From the expression (2.21), this occurs only for the Schwarzschild potential $`\nu _S.`$ Briefly, the reason for this is as follows. The estimate (2.30) implies that the potential $`\upsilon =\nu \nu _S`$ is bounded and given by convolution of $`dist^1`$ with a signed Radon measure $`d\lambda .`$ However, as in the proof of Lemma 2.4, if $`\upsilon `$ is bounded then one sees that necessarily $`d\lambda <<dA`$ and further the Radon-Nikodym derivative $`d\lambda /dA`$ is 0, a.e. $`(dA)`$. In other words, any point of non-zero density for $`d\lambda `$ w.r.t. $`dA`$ gives rise to approximating points on which $`\nu `$ is unbounded. It follows that $`d\lambda =`$ 0, and hence $`\upsilon =`$ 0.
In all other cases, where $`f`$ 0 on approach to a region in supp $`d\mu ,`$ the length of the $`\varphi `$ circles goes to 0, and hence this portion of $`\mathrm{\Sigma }`$ is singular, of dimension $``$ 1, while where $`f\mathrm{},`$ the length of the circles goes to $`\mathrm{}`$ and the corresponding portion of $`\mathrm{\Sigma }`$ is singular and non-compact, (as in the Curzon solution). Note that if supp $`d\mu `$ is compact, then in all cases, $`M`$ is pseudo-compact.
Thus among the Weyl solutions given by a Newtonian potential, only the Schwarzschild metric is smooth up to $`\mathrm{\Sigma }.`$ This gives a very simple (local) confirmation of Theorem 0.1 in this special case.
Example 2.6.(iii). (Superposition/Multiple Holes). Subharmonic functions of the form (2.21) form a convex cone. In particular, one thus has a natural linear superposition principle for Weyl solutions. This feature is another remarkable property of Weyl solutions.
For example, one may choose the measure $`d\mu =\frac{1}{2}dA`$ on two, or any number of disjoint intervals $`\{I_j\}`$ on the axis $`A`$, provided (2.27) holds. These correspond to solutions with ’multiple black holes’, each interval $`I_j`$ giving a component of $`\mathrm{\Sigma }`$ which is a 2-sphere of radius equal to the length of $`I_j.`$ Although such solutions are essentially smooth up to $`\mathrm{\Sigma },`$ they do not define smooth vacuum solutions on $`^3B_j.`$ There are cone singularities, (called struts or rods in the physics literature), along geodesics (corresponding to $`A\{I_j\})`$ joining the 2-spheres of $`\mathrm{\Sigma },`$ so that the metric $`g`$ is not locally Euclidean along such curves. Thus, the elementary flatness condition (2.14) is violated on $`AI_j.`$ Nevertheless, the curvature of such metrics is uniformly bounded everywhere. These cone singularities must be considered part of $`M,`$ so $`M`$ in this case is a union of 2-spheres joined by a collection of curves and thus connected.
Of course, the black hole uniqueness theorem, Theorem 0.1, also implies that such solutions cannot be smooth everywhere, when the number of intervals is finite. However, the proof of this result strongly uses the asymptotically flat assumption. It seems to be unknown whether there are any Schwarzschild type metrics with infinitely many black holes, i.e. metrics complete away from $`\mathrm{\Sigma }`$ and smooth up to $`\mathrm{\Sigma }`$ with $`\mathrm{\Sigma }`$ consisting of infinitely many 2-spheres, and which satisfy (2.17) or (2.21). It is natural to conjecture that such solutions do not exist, c.f. however the end of Remark 2.8 below.
###### Remark 2.8.
Whenever supp $`d\mu A`$ is not connected but compact, there will exist such cone singularities on $`A`$ supp $`d\mu `$. When supp $`d\mu ,`$ or a sufficiently small smoothing of supp $`d\mu ,`$ is interpreted to represent a matter source, then this statement corresponds exactly to the fact that there are no equilibrium (i.e. time-independent) many-body solutions in Newtonian gravity of this character, c.f. \[Bn\]. Even when supp $`d\mu `$ is non-compact and disconnected, this seems very likely to be true, c.f. the expressions for sums of Schwarzschild rods (2.29) in \[IK, p.336-337\], which generalize to infinitely many rods.
The components $`\{C_i\}`$ of $`A`$ supp $`d\mu `$ represent idealized matter sources (struts or rods) keeping the components of supp $`d\mu `$ apart in static equilibrium. The cone angle $`\alpha =\alpha _i`$ is constant on each $`C=C_i,`$ and corresponds to a concentration of scalar curvature on $`C\overline{M}`$ given by a multiple of the Lebesgue measure on $`C`$
$$s=(1\alpha )dA_C,$$
Thus the vacuum equations (0.1) are not satisfied across $`\{C_i\}.`$ If a very small tubular neighborhood of radius $`r`$ of such a rod is smoothly filled in with a perfect fluid source of constant pressure $`\rho `$ and energy density $`\mu ,`$ then one has the relation $`lim_{r0}r\rho =lim_{r0}r\mu >`$ 0. The effective mass of such rods is zero, i.e. they do not contribute to the gravitational potential $`\nu ,`$ c.f. \[I2\] for a detailed discussion.
By passing to covering spaces, it is always possible to create such cone singularities in Weyl solutions, even if none existed to begin with. For instance, for the Schwarzschild solution (0.4), with potential (2.29), take any covering, including the universal covering, of $`^3A=(S^2\{aa\})\times ^+,`$ where {$`a`$, $`a`$} are two antipodal points on $`S^2.`$ This gives a solution whose metric completion has cone singularities along two (radial) geodesics starting at the antipodal points on $`S^2=\mathrm{\Sigma }`$ and going to infinity.
Note that this discussion assumes that $`\nu `$ is given by a Newtonian potential (2.21). In fact, there are Weyl solutions $`(M,g,u)`$ everywhere smooth up to the axis $`A`$, with $`\mathrm{\Sigma }=M`$ disconnected, with no cone singularities or struts keeping the components of $`\mathrm{\Sigma }`$ apart. Namely, the B1 solution (0.9), dual to the Schwarzschild solution, has this property.
Another, more remarkable, example is given in \[KN\]. These authors construct a Weyl solution of the form (2.12), which is complete away from $`\mathrm{\Sigma }`$ and smooth up to $`\mathrm{\Sigma }`$, with $`\mathrm{\Sigma }`$ consisting of infinitely many Schwarzschild-like 2-spheres. In fact, the solution is periodic in the $`z`$-direction. This metric is not of the form (2.21), but is a limit of a sequence of solutions of the form (2.20), c.f. (D) below.
###### Remark 2.9.
If $`(M,g,u)`$ is a Weyl solution of the form (2.12) with $`\nu =\nu _u=`$ $`logu`$, then the dual solution $`(M^{},g^{},f)`$, discussed in (2.9), is also a Weyl solution of the form (2.12), with potential $`\nu _f`$ = $`logf`$ given by
(2.31)
$$\nu _f=logr\nu _u.$$
Hence if one potential is Newtonian, the dual one is not. Note that the sets $`I_u,I_f`$ where $`\nu _u`$ and $`\nu _f`$ are $`\mathrm{}`$ are disjoint, with $`\overline{I}_u\overline{I}_f=A`$. Hence, if $`\nu _u`$ is a Newtonian potential with supp $`d\mu `$ compact, so that the associated Weyl solution is asymptotically flat, then the dual Weyl solution is asymptotically cylindrical.
Another example of a potential where both terms in (2.20) are non-trivial is the situation considered (locally) in \[GH\], where $`h`$ is a smooth axi-symmetric harmonic function defined on a neighborhood of supp $`d\mu ,`$ c.f. Remark 1.5. As vacuum solutions, these metrics cannot be complete away from $`\mathrm{\Sigma },`$ as in the discussion on pure harmonic potentials.
This completes our discussion of the Positive Case.
(B).(Negative Case). Under the assumption (2.16), suppose now instead that $`\nu `$ is locally bounded below in $`^3,`$ i.e.
(2.32)
$$\underset{B_x(1)}{inf}\nu >\mathrm{},x\mathrm{\Omega }.$$
Then $`\nu `$ extends uniquely to a globally defined superharmonic function on $`^3.`$ Exactly the same discussion as in (A) above holds here, under the substitution $`\nu \nu .`$ (This corresponds to the transformation $`uu^1`$ following (1.5)). In this case, the Riesz measure is a negative measure, so that one has solutions with negative mass. Note that here the potentials $`\nu `$ or $`u`$ are unbounded above within supp $`d\mu ,`$ i.e. $`u`$ or $`\nu `$ go to $`+\mathrm{}`$.
(C).(Mixed Case). Next, one may superimpose Weyl solutions with positive and negative measures $`d\mu ,`$ i.e. consider $`\nu `$ of the form (2.21), with $`d\mu `$ a signed Radon measure. For example, one may form dipole-type solutions with potential of the form $`\nu =\nu _++\nu _{}`$ where $`\nu _+`$ and $`\nu _{}`$ are (for instance) Curzon or Schwarzschild solutions of positive and negative mass placed at different regions on the axis. This gives for instance examples of asymptotically flat solutions with mass $`m`$ assuming any value in $`.`$
More generally, since positivity is no longer assumed, the measure $`d\mu `$ may be replaced by distributions, for example weak derivatives of measures.
###### Example 2.10.
(Multipole Solutions). As a typical example, one may take potentials corresponding to derivatives of the Dirac measure based at a point $`aA,`$ i.e. the multipole potentials,
$$R^{n1}P_n(\frac{z}{R}),$$
where $`P_n`$ is the $`n^{th}`$ Legendre polynomial, or arbitrary linear combinations of such; c.f. \[MF, p.1276ff\].
Such potentials are limits of combinations of Newtonian potentials with positive and negative mass. Thus, it is reasonable to expect that there are (Newtonian) equilibrium solutions, i.e. solutions with no cone singularities on the axis. This is proved in \[Sz\], where explicit equilibrium conditions are given. Note that one may have infinitely many multipole ’particles‘ in equilibrium.
Another example is the potential of a dipole ring
$$\nu (x)=_C\frac{z(x)}{|x\xi |^3}𝑑\xi ,$$
where $`x=(r,z,\varphi )`$ and $`d\xi `$ is the Lebesgue measure on the unit circle $`C=\{r=1\}`$ in the $`z=`$ 0 plane. Here $`\nu (x)\mathrm{},`$ as $`xC`$ along the rays $`r=`$ 1, $`z>`$ 0. Hence, in this case, the set $`I=C`$ is not contained in the axis $`A`$.
(D). (Limits). Finally, one may consider potentials which are limits of potentials of the type (A)-(C) above, (besides those in Remark 2.9, Example 2.10). We consider just one important instance of this here.
###### Example 2.11.
(Kasner Metric). It is easily seen that the potential $`\nu =logr`$ generates the flat metric
$$g=dr^2+dz^2+d\varphi ^2,$$
on $`(^3)^+`$, with potential $`u=r`$. Observe that since the $`\varphi `$-lines have constant length, the function $`r`$ is now an affine (in fact linear) function on $`(^3)^+.`$
On the other hand, an equally simple computation shows that the potential $`\nu =alogr,`$ for any $`a,`$ generates the metric
$$g=r^{2a^22a}(dr^2+dz^2)+r^{22a}d\varphi ^2,u=r^a.$$
Equivalently, setting $`s=r^{a^2a+1},`$
(2.33)
$$g=ds^2+s^\alpha dz^2+s^\beta d\varphi ^2,u=s^\gamma ,$$
where $`\alpha =(2a2)/(a1+a^1),\beta =(2a^12)/(a1+a^1),\gamma =(a1+a^1)^1.`$ Here $`s^+,z`$ and $`\varphi [0,2\pi ]`$ or any other interval, including $`,`$ (by passing to covering and quotient spaces). These metrics are all non-homothetic, provided $`a[1,0)(0,1]`$; $`a=0`$ gives the flat metric with $`u=1`$ while $`a=1`$ gives the $`A3`$ metric (1.8).
The potential $`\nu =alogr`$ can be considered as the limit
(2.34)
$$\nu =lim_m\mathrm{}a[\nu _S(m)log2m],$$
where $`\nu _S(m)`$ is the Schwarzschild potential (2.29) of mass $`m`$. Thus it is a limit of potentials of the form (2.18), where both terms are non-zero, (the harmonic term $`h`$ is of course constant here).
These metrics are dual, in the sense discussed in (2.9) to the Kasner (or Bianchi I) vacuum cosmological models, with homogeneous (flat) but anisotropic space-like hypersurfaces, c.f. \[Wd, Ch.7.2\]. It is easy to see that the Kasner metrics are the only Weyl solutions $`(M,g)`$ which have an isometric $`\times `$ action, even locally. (The axisymmetric potential $`\nu `$ on $`\mathrm{\Omega }^3`$ must be invariant under an orthogonal $``$-action, hence giving a rotationally invariant harmonic function on $`^2`$. Thus the potential must be a multiple of log $`r`$).
Consider these metrics on the quotient $`M=^+\times S^1\times S^1.`$ In case $`a>0`$, we have $`\alpha <0`$, $`\beta >0`$, $`\gamma >0`$ and so $`\mathrm{\Sigma }=M.`$ As in the discussion with the Curzon metric, the $`z`$-circles have unbounded length as $`s`$ 0, so that $`M=`$ and the levels $`t=ϵ`$, $`(t(x)=dist_g(x,M)),`$ are long, thin cigars. Thus, $`M`$ is non-compact, but pseudo-compact. The end of $`\overline{M}U`$ is obviously small.
If $`a<`$ 0, then $`\alpha >`$ 0, $`\beta >`$ 0, $`\gamma <`$ 0 and so $`\mathrm{\Sigma }=\mathrm{},`$ (it occurs at infinity), with $`M=`$ {pt}. In this case, the end of $`\overline{M}U`$ is not small; the area growth of geodesic spheres is $`O(r^{1\gamma })`$.
However, none of these solutions are asymptotically flat even in a weak sense, except of course when $`a=1`$. Namely, the curvature decays only quadratically in the $`g`$-distance to $`M,`$ i.e.
$$|r|=O(t^2),$$
and not any faster. Hence, the case $`a>`$ 0 shows that the conclusion of Theorem 0.3, (asymptotically flat or small ends), cannot be strengthened to only asymptotically flat ends, while the case $`a<`$ 0 shows that the assumption (ii) on $`u`$ in Theorem 0.3 is necessary.
Another metric of this (limit) type is that constructed in \[KN\], referred to in Example 2.8. This metric has the same asymptotics as the Kasner metric.
## 3. Characterization of Asymptotically Flat Solutions.
In this section, we prove Theorem 0.3. The proof of the first statement on finiteness of the number of ends is quite easy, so we begin with this.
Throughout this section, let $`(M,g,u)`$ be a static vacuum solution with $`M`$ pseudo-compact. We recall from §0 that $`M`$ is connected and oriented. As in §0, let
$$t(x)=dist_g(x,M),$$
and suppose $`U=t^1(0,s_o),`$ so that $`UM`$ is compact. For $`r,ss_o,`$ let $`S(s)=t^1(s),A(r,s)=t^1(r,s)`$ be the geodesic spheres and annuli about $`U.`$ It is important to note that neither $`S(s)`$ nor $`A(r,s)`$ are necessarily connected, even if $`M`$ has only one end. (Of course if $`E`$ is a given end of $`M`$, then $`S_E(s)=S(s)E`$ must be connected for some sequence $`s=s_j\mathrm{}`$). Let $`S_c(s)`$ and $`A_c(r,s)`$ denote any component of $`S(s)`$ resp. $`A(r,s)`$, so that $`S(s)=S_c(s),A(r,s)=A_c(r,s).`$ Of course $`t`$ is a proper exhaustion function on $`\overline{M}U,`$ so that these sets have compact closure in $`M`$.
Let $`diam^iA_c(r,s)`$ denote the intrinsic diameter of $`A_c(r,s),`$ i.e. the diameter of the connected metric space $`(A_c(r,s),g)`$.
###### Lemma 3.1.
There exists a constant $`d_o<\mathrm{},`$ independent of $`s`$, such that the number of components of $`A(\frac{1}{2}s,2s)`$ is at most $`d_o`$ and
(3.1)
$$diam^iA_c(\frac{1}{2}s,2s)d_os.$$
In particular, the manifold $`(\overline{M}U,g)`$ has a finite number of ends $`\{E_i\}.`$
Proof: Consider first the 4-manifold $`(N,g_N),`$ as in (0.2) which is smooth and Ricci-flat outside $`\overline{U}=\pi ^1(U),`$ where $`\pi :NM`$ is projection on the first factor. Since $`\overline{U}N`$ is compact, it follows from results of \[Lu\] that Lemma 3.1 holds on $`N`$, so that in particular $`N`$ has a finite number of ends. Since $`N=M\times S^1,\overline{M}U`$ also has a finite number of ends.
The choice of the time parameter on $`N`$ defines a totally geodesic embedding $`MN`$ and we have $`t_N|_M=t`$ where $`t_N(x)=dist_N(x,M),MN.`$ A geodesic ball or annulus in $`M`$ embeds in the geodesic ball or annulus of the same size in $`N`$. Hence (3.1) also holds for $`M`$.
Lemma 3.1 of course proves the first statement of Theorem 0.3. Observe that the estimate (3.1) is invariant under rescaling of the metric $`g`$.
For the remainder of the proof, we (usually) work with a given end $`E`$ from the finite collection $`\{E_i\}.`$ The main statement of Theorem 0.3 is that if
(3.2)
$$^{\mathrm{}}\frac{1}{areaS_E(s)}𝑑s<\mathrm{},$$
then the end $`E`$ is asymptotically flat. The proof of this result is rather long, so we outline here the overall strategy. The asymptotic behavior of $`(E,g,u)`$ is studied in general by examining the structure of the possible tangent cones at infinity, defined below. Basically, tangent cones at infinity fall into two classes, according to whether the asymptotic geometry near a given divergent sequence of base points is non-collapsing or collapsing, c.f. Lemmas 1.3-1.4. The main point is to prove that under the bound (3.2), all tangent cones at infinity are flat manifolds, and further that no collapse behavior is possible. Once this is established, the proof that $`(E,g,u)`$ is asymptotically flat is relatively straightforward.
Apriori, the end $`E`$ may be very complicated topologically, for instance of infinite topological type; consider for instance that $`E`$ might be of the form $`S_{\mathrm{}}\times S^1,`$ where $`S_{\mathrm{}}`$ is any non-compact surface of infinite topological type and one end. A main idea is to use the behavior of the potential function $`u`$, in particular its value distribution theory, to control the topology and geometry of $`E`$ in the large. We have already seen in §2 that the potential $`u`$ controls quite strongly the geometry of Weyl solutions. Lemma 3.6 below is the key technical lemma which expresses this control for general static vacuum solutions (with pseudo-compact boundary). Further remarks on the strategy of proof precede the Lemmas below.
The discussion to follow, until the end of Lemma 3.6, holds in general for ends $`E`$ of static vacuum solutions with compact boundary. The estimate (3.2) will only be used after this.
We now define the tangent cones at infinity of a given end $`E`$. (While this is a commonly used terminology, such limit metric spaces are not necessarily metric cones in general).
First, we recall by Theorem 1.1 that there is a constant $`K<\mathrm{}`$ such that, $`xM,`$
(3.3)
$$|r|(x)\frac{K}{t^2(x)},|dlogu|(x)\frac{K}{t(x)}.$$
The scale-invariant estimates (3.3) give quite strong initial control on the asymptotic geometry of $`(E,g,u)`$ which allows one to get started. Observe that an immediate consequence of (3.1) and (3.3), by integration along paths in $`A_c(\frac{1}{3}s,3s),`$ is the following Harnack inequality:
(3.4)
$$\frac{supu}{infu}K_1,$$
where the sup and inf are taken over any component $`A_c(\frac{1}{2}s,2s)`$ and $`K_1`$ is independent of $`A_c(\frac{1}{2}s,2s).`$
Let $`x_i`$ be any divergent sequence of points in $`E`$, with $`t_i=t(x_i)\mathrm{}.`$ Consider the connected geodesic annuli $`A_i=A_i(\kappa )=A_c(\kappa ^1t_i,\kappa t_i),x_iA_i,`$ w.r.t. the rescaled or blow-down metric
(3.5)
$$g_i=t_i^2g;$$
here $`\kappa `$ is any fixed positive constant $`>`$ 1. By the curvature estimate (3.3), the metrics $`g_i`$ have uniformly bounded curvature on $`A_i`$ \- the curvature bound depends only on $`K`$ and $`\kappa .`$ Further, by (3.1), the diameter of $`A_i`$ w.r.t. $`g_i`$ is also uniformly bounded.
Hence, if the sequence is non-collapsing, i.e. if there is a lower volume bound $`vol_gA_i\nu _ot_i^3,`$ for some $`\nu _o>`$ 0, (equivalent to $`vol_{g_i}A_i\nu _o`$ by scaling), then Lemma 1.3 implies that a subsequence of the pointed sequence $`\{(A_i,g_i,x_i)\},`$ converges smoothly, away from the boundary, to a limiting smooth metric $`(A_{\mathrm{}}(\kappa ),g_{\mathrm{}},x_{\mathrm{}}).`$ The limit is a solution of the static vacuum equations; as noted in Lemma 1.3, the potential $`u`$ is renormalized to $`u_i=u/u(x_i),`$ so that the limit potential $`u_{\mathrm{}}`$ satisfies $`u_{\mathrm{}}(x_{\mathrm{}})=`$ 1, c.f. also (3.4). Choosing a sequence $`\kappa _j\mathrm{}`$ and a suitable diagonal subsequence, gives the maximal static vacuum solution $`(A_{\mathrm{}},g_{\mathrm{}},u_{\mathrm{}},x_{\mathrm{}}).`$ Observe here also that the estimate (3.1) implies that $`A_{\mathrm{}}=`$ {pt}.
On the other hand, if the sequence $`(A_i,g_i)`$ is collapsing, in the sense that $`vol_gA_i<<t_i^3,`$ as $`t_i\mathrm{},`$ (equivalent to $`vol_{g_i}A_i`$ 0), then Lemma 1.4 implies that $`A_i`$ is a Seifert fibered space or torus bundle over an interval. As discussed there, one may then pass to $``$ or $``$ covers $`\stackrel{~}{A}_i(\kappa )`$ of $`A_i(\kappa )`$ (or more precisely smooth interior approximations to $`A_i(\kappa ))`$ to obtain a non-collapsing sequence $`(\stackrel{~}{A}_i(\kappa ),g_i,x_i)`$, smoothly convergent to a limit $`(\stackrel{~}{A}_{\mathrm{}}(\kappa ),g_{\mathrm{}},x_{\mathrm{}})`$; ($`g_i`$ here is lifted to the cover $`\stackrel{~}{A}_i,`$ as is $`x_i`$). As above, one may then choose a sequence $`\kappa _j\mathrm{}`$ and pass to a diagonal subsequence to obtain a maximal limit $`(\stackrel{~}{A}_{\mathrm{}},g_{\mathrm{}},u_{\mathrm{}}).`$ This limit static vacuum solution has an isometric $``$ action, or $``$ action in the case of a rank 2 collapse. Hence by Proposition 2.2 it is a Weyl solution. In the latter case, the solution is then a (possibly flat) Kasner metric, c.f. Example 2.11. Topologically, the limit $`\stackrel{~}{A}_{\mathrm{}}`$ is a trivial $``$, (or $``$), bundle over a surface $`V`$, (or interval), again by Proposition 2.2.
As in §2, we will always work in a $``$-quotient $`\overline{A}_{\mathrm{}}`$ of $`\stackrel{~}{A}_{\mathrm{}},`$ (or $``$ quotient in the case of rank 2 collapse), and finite covers $`\overline{A}_i(\kappa )`$ converging to $`\overline{A}_{\mathrm{}}(\kappa ).`$ For any fixed $`\kappa >`$ 0, the manifolds $`(\overline{A}_i(\kappa ),g_i,x_i)`$ thus have uniformly bounded curvature and diameter, a uniform lower bound on their volume, and converge smoothly to the limit $`(\overline{A}_{\mathrm{}}(\kappa ),g_{\mathrm{}})(\overline{A}_{\mathrm{}},g_{\mathrm{}}).`$ The limit has a free isometric $`S^1`$ or $`S^1\times S^1`$ action, and so in particular is an $`S^1`$ or $`T^2`$ bundle. Hence $`\overline{A}_i(\kappa )`$ is also topologically an $`S^1`$ or $`T^2`$ bundle, (although not metrically).
To be definite, the finite covers are chosen so that the length of the $`S^1`$ factor or factors at the base point $`x_i\overline{A}_i(\kappa )`$ converge to $`1`$ in the limit.
Recall by Lemma 1.4 that the inclusion map of the fibers induces an injection on $`\pi _1`$. The coverings $`\overline{A}_i(\kappa )`$ are obtained by taking large finite unwrappings of the $`S^1`$ or $`T^2`$ fibers, (corresponding to taking subgroups of $`\pi _1(S^1)`$ or $`\pi _1(T^2)`$ of large but finite index). All finite covering spaces of $`S^1`$ or $`T^2`$ are still $`S^1`$ or $`T^2`$, and hence we may, and do, choose the unwrappings so that, as smooth manifolds,
$$\overline{A}_i(\kappa )=A_i(\kappa ),$$
for any $`\kappa >`$ 0. In the limit, the unwrapping of the collapse thus just corresponds to expanding the length of the collapsing $`S^1`$ factor (or factors), preserving the holonomy, if any, of the $`S^1`$ bundle; compare with the proof of Proposition 2.2.
The limit spaces $`(A_{\mathrm{}},g_{\mathrm{}},u_{\mathrm{}},x_{\mathrm{}})`$ or $`(\overline{A}_{\mathrm{}},g_{\mathrm{}},u_{\mathrm{}},x_{\mathrm{}})`$ constructed above are called tangent cones at infinity of $`(E,g,u)`$. Note that such tangent cones are only attached to some subsequence of a given divergent sequence of base points $`\{x_i\}.`$ Hence, apriori, the tangent cones at infinity could be highly non-unique as Riemannian manifolds. In general, there may be no relation between the geometry of different tangent cones based on (subsequences of) distinct divergent sequences $`\{x_i\};`$ for example, tangent cones based on sequences with $`t(x_i)=2^{i^2}`$ and $`t(x_i)=2^{i^3}.`$ The tangent cones only detect behavior of the end $`E`$ in $`g_i`$-bounded distance to the base points $`x_i.`$
On the other hand, since tangent cones at infinity attached to any divergent sequence always exist, for any $`s`$ sufficiently large, say $`ss_o`$, the geometry of $`(A_c(\frac{1}{2},2),g_s)`$ or $`(\overline{A}_c(\frac{1}{2},2),g_s)`$ is always close to that of some tangent cone at infinity $`A_{\mathrm{}}`$ or $`\overline{A}_{\mathrm{}}`$. Further, by construction, the tangent cones are always connected and, since $`M`$ is oriented, so is each tangent cone.
The following lemma is a typical application of the use of tangent cones at infinity.
###### Lemma 3.2.
Suppose the curvature $`r`$ decays faster than quadratically in the end $`(E,g,u)`$, i.e.
(3.6)
$$|r|(x)\frac{\epsilon (t)}{t^2(x)},$$
where $`\epsilon (t)`$ 0 as $`t\mathrm{}.`$ Then there is a compact set $`KE`$ such that $`EK`$ is diffeomorphic either to $`^3B`$ or to $`(^2B)\times S^1,`$ where $`B`$ is a 3-ball, (resp. a 2-ball), i.e. $`E`$ is of standard topological type. Further, the annuli $`A_E(\frac{1}{2}s,2s)`$ are connected, for all $`ss_o`$, for some $`s_o<\mathrm{}`$.
Proof: By the preceding discussion, the condition (3.6) is equivalent to the statement that all tangent cones at infinity $`(A_{\mathrm{}},g_{\mathrm{}})`$ or $`(\overline{A}_{\mathrm{}},g_{\mathrm{}})`$ of $`E`$ are flat, (as well as connected and oriented). The two possible conclusions of Lemma 3.2 correspond to the two possibilities of non-collapse and collapse in the formation of the tangent cones.
Suppose first that $`g`$ is non-collapsing on $`E`$, i.e. there exists $`\nu _o>`$ 0 such that $`vol_gA_c(\frac{1}{2}s,2s)\nu _os^3,`$ for all $`s`$ large, and all components $`A_c`$. Recall that $`A_{\mathrm{}}=\{pt\}`$ in this situation. It then follows that for $`s`$ sufficiently large, each $`A_c(\frac{1}{2}s,2s)`$ is diffeomorphic, and almost isometric to the standard flat annulus $`A=r^1(\frac{1}{2}s,2s)`$ in $`^3,`$ $`r(x)=|x|,`$ (away from the boundary). In fact, each tangent cone at infinity $`A_{\mathrm{}}`$ is isometric to $`^3\{0\}`$ in this situation. Here we are implicitly using the fact that the only complete oriented flat 3-manifold with an isolated singularity is $`^3\{0\}`$, c.f. \[AC\] for example. Similarly, a smooth approximation to $`S_c(s)`$ is diffeomorphic and almost isometric to $`S^2(s)^3.`$
By the isotopy extension theorem, these diffeomorphisms from $`A_c(\frac{1}{2}s,2s)`$ to the standard annulus may then be assembled to a global diffeomorphism, and almost isometry, of $`EK`$ into $`^3B,`$ for some compact set $`KE`$. We refer to \[AC, Thm.1.18\] for the proof of these statements, (in a slightly different but equivalent form), which are now quite standard. The main point is of course that since the family of annuli $`A_c(\frac{1}{2}s,2s)`$ as $`s`$ varies is topologically rigid, i.e. one has a unique topological type, there is no value of $`s`$ at which the topology can change or bifurcate.
Suppose on the other hand that $`g`$ is collapsing on $`A_c(\frac{1}{2}s_i,2s_i),`$ for some sequence $`s_i\mathrm{}`$, and some sequence of components $`A_c`$. As discussed above, one may then pass to suitable covers $`\overline{A}_i=\overline{A}_c(\frac{1}{2}s_i,2s_i)`$ so that, in a subsequence, $`(\overline{A}_i,g_i)`$ is diffeomorphic and almost isometric to its limit $`\overline{A}_{\mathrm{}}(\frac{1}{2},2)\overline{A}_{\mathrm{}}`$. The maximal limit $`\overline{A}_{\mathrm{}}`$ is a flat manifold with either a free isometric $`S^1`$ or $`S^1\times S^1`$ action. Hence there are two possibilities for $`\overline{A}_{\mathrm{}}`$, namely either $`V\times S^1`$ or $`^+\times S^1\times S^1`$, where $`V`$ is a flat 2-manifold and the metric is a product metric on each $`S^1`$ factor. In the former case, the diameter estimate (3.1) implies, (as in the non-collapse case above), that $`V`$ is a complete flat cone, possibly with an isolated singularity at $`\{0\}`$. Hence, although these two possibilities for the limiting metric of $`\overline{A}_{\mathrm{}}(\frac{1}{2},2)`$ are distinct, both are the same topologically, i.e. $`\overline{A}_{\mathrm{}}(\frac{1}{2},2)`$ is topologically $`I\times S^1\times S^1`$.
Now recall from the discussion on tangent cones that $`\overline{A}_i`$ is diffeomorphic to $`A_c(\frac{1}{2}s_i,2s_i)`$; metrically $`\overline{A}_i`$ approximates one of the types of flat manifolds above, with $`S^1`$ factors shrinking to very short circles. In both cases, (a smoothing of) $`\overline{S}_c(s_i)`$ is diffeomorphic and almost isometric to a flat torus $`T^2.`$
In particular, the topological type of $`\overline{A}_i`$ is distinct from that of the annuli $`A_i`$ above in the non-collapse case, which are topologically always of the form $`I\times S^2^3\{0\}`$, (for any choice of base point and component). This implies first that the family $`\{A_c(\frac{1}{2}s,2s)\}`$ must be collapsing for all $`s`$, as $`s\mathrm{}`$, and all components $`A_c`$. Second, the topological type of the annuli $`\overline{A}_s`$, and hence that of $`A_s=A_c(\frac{1}{2}s,2s)`$ is unique, and given for all $`s`$ large and all $`c`$ by $`I\times S^1\times S^1`$. Use of the isotopy extension theorem in the same way as above then proves that the end $`E`$ itself is diffeomorphic to $`(^2B)\times S^1`$.
###### Remark 3.3.
(i) It is easily seen from the vacuum equations (0.1) that the condition (3.6) is equivalent to
(3.7)
$$|logu|(x)\frac{\epsilon (t)}{t(x)},$$
as $`t\mathrm{},`$ c.f. also the proof of Theorem 1.1.
(ii). In the context of Lemma 3.2, the tangent cones at infinity $`(\overline{A}_{\mathrm{}},g_{\mathrm{}})`$ may not be unique up to isometry in the collapse case, and so may vary within the moduli space $`_1`$ of flat product metrics of the form $`V\times S^1`$ or within the moduli space $`_2`$ of flat product metrics of the form $`^+\times S^1\times S^1`$.
Note that the moduli space $`_o`$ of flat metrics on $`^3`$ or $`^3\{0\}`$ is just one point, (c.f. again \[AC\] for the latter statement for example). Similarly, by the normalization preceding Lemma 3.2 that the $`S^1`$ factors have length 1, the moduli space $`_2^{}_2`$ normalized in this way is also just one point. The moduli space $`_1^{}_1`$ where the $`S^1`$ factor has length 1 is naturally identified with $`^+`$, parametrized by the cone angle at $`\{0\}`$.
Observe however that these two moduli spaces $`_1^{}`$ and $`_2^{}`$ are disjoint; they cannot be connected (or even approximated) by a curve of flat metrics. Now the geometry of the annuli $`(A_E(\frac{1}{2}s,2s),g_s)`$ varies continously with $`s`$. By the remarks preceding Lemma 3.2, this induces a continuous variation of the possible tangent cones $`(\overline{A}_{\mathrm{}},g_{\mathrm{}})`$ in $`_1^{}`$ or $`_2^{}`$. Hence, on a given end $`E`$, one cannot obtain two different tangent cones, one of the form $`V\times S^1`$ and another of the form $`^+\times S^1\times S^1`$; c.f. also the proof of Lemma 3.7 below.
(iii). We will need a slight generalization of Lemma 3.2 for the next lemma below. Thus let $`\gamma (s)`$ be any properly embedded curve in $`E`$ with $`t(\gamma (s))\mathrm{}`$ as $`s\mathrm{}`$, and suppose (3.6), (or (3.7)), holds in the balls $`B_{\gamma (s)}(\delta t(\gamma (s)))`$, for some fixed $`\delta >0`$. Then the conclusion of Lemma 3.2 also holds.
To see this, consider the blow-downs $`g_s=t^2(\gamma (s))g`$, based at $`\gamma (s)`$, and the associated tangent cones at infinity. The scale invariant condition (3.6), when applied to $`B_{\gamma (s)}(\delta t(\gamma (s)))`$, implies that all such tangent cones are flat in $`(B_x_{\mathrm{}}(\delta ),g_{\mathrm{}})`$, and thus flat everywhere in their maximal domain $`A_{\mathrm{}}`$, by the fact that smooth solutions of the vacuum equations are real analytic. The proof then proceeds exactly as in Lemma 3.2
To prove that the estimates (3.6) and (3.7) do in fact hold on $`E`$, we need to understand in more detail the value distribution of the potential $`u`$. The main result needed for this is given in Lemma 3.6, and then (3.6)-(3.7) follow rather easily in (3.17)-(3.18) below. However, some preliminary results are required for the proof of Lemma 3.6. The main difficulty is that $`u`$ may not, in this generality, be a proper function onto its image, (c.f. the remark following the proof of Lemma 3.4). Lemma 3.4 below is a slightly weaker substitute for this property.
Let $`U=t^1(0,s_o)`$ be a neighborhood of $`M`$ as in the beginning of §3. Let $`\gamma (\tau )`$ be a maximal flow line of $`u.`$ We will say that $`\gamma `$ is divergent if $`\gamma `$ does not intersect $`U`$ at two different times, i.e. if $`\gamma (\tau )`$ exits $`U`$ at some time, then $`\gamma (\tau )`$ never reenters (a possibly distinct component of) $`U`$, and if further $`\gamma (\tau )`$ does not terminate at a critical set of $`u`$ in $`\overline{M}U`$ as $`\tau \pm \mathrm{}.`$ It follows that if $`\gamma `$ is divergent, then $`\gamma `$ is complete in at least one direction, $`(\tau +\mathrm{},\tau \mathrm{}`$ or both), and in any such direction, $`\gamma (\tau )`$ diverges to infinity in $`M`$. Since the potential $`u`$ has no local maxima or minima in $`M`$, the set of flow lines terminating on a critical set of $`u`$ in $`MU`$ is a closed set of measure 0 in $`M`$. This follows for instance from the fact that the measure $`|du|dA`$, where $`dA`$ is Lebesgue measure on the level sets of $`u`$, is preserved under the gradient flow of $`u`$, and this measure tends to 0 on approach to critical points of $`u`$. Thus among the set of flow lines not joining points of $`U`$, the divergent flow lines are generic in terms of measure on $`M`$. Of course the flow lines are curves of steepest ascent for $`u`$ as $`\tau `$ increases, and of steepest descent for $`u`$ as $`\tau `$ decreases.
###### Lemma 3.4.
There exists a compact set $`K\overline{M}`$, with $`\overline{U}K`$, such that any divergent flow line $`\gamma (\tau )`$ intersects K.
Proof: Suppose that this were not the case, so that there exist, necessarily complete, flow lines $`\gamma (\tau ),\tau ,`$ which do not intersect a given $`K\overline{U}`$. We may choose $`K`$ sufficiently large so that $`\gamma (\tau )`$ is then contained in a fixed end $`E\overline{M}U,`$ since there are only finitely many ends.
Let $`A_\gamma (\tau )=A_\gamma (\frac{1}{2}t(\gamma (\tau )),2t(\gamma (\tau )))`$ be the component of the geodesic annulus containing the base point $`\gamma (\tau )`$ and let $`E_\gamma `$ be the part of $`E`$ swept out by such annuli, $`E_\gamma =_\tau A_\gamma (\tau )E`$. As $`\tau \mathrm{}`$, the function $`u(\gamma (\tau ))`$ is monotone increasing.
If $`u(\gamma (\tau ))`$ increases to $`+\mathrm{}`$, then $`u\mathrm{}`$ uniformly as $`\tau \mathrm{}`$ in $`E_\gamma ^+=_{\tau >0}A_\gamma (\tau )E`$, by the Harnack estimate (3.4). Since $`E`$ is an end, there exists some sequence $`\tau _j\mathrm{}`$ such that, for $`t_j=t(\gamma (\tau _j))`$, the spheres $`S_\gamma (t_j)A_\gamma (\tau _j)`$ satisfy $`S_\gamma (t_j)=S_E(t_j)`$, i.e. the spheres $`S_E(t_j)`$ are connected, and hence $`A_\gamma (\tau _j)=A_E(\tau _j)`$. Thus, $`u`$ becomes uniformly unbounded on $`A_E(\tau _j)`$, as $`j\mathrm{}`$. However, as $`\tau \mathrm{}`$, the curve $`\gamma (\tau )`$ also diverges to infinity in $`E`$ and $`u(\gamma (\tau ))`$ is decreasing, (and so in particular bounded), as $`\tau \mathrm{}`$. This contradiction implies that $`u(\gamma (\tau ))u^+<+\mathrm{}`$, as $`\tau +\mathrm{}`$. Of course $`u^+>0`$.
Suppose first, (for simplicity), that $`limsup_Eu=u^+`$, i.e. $`lim_t\mathrm{}m(t)=u^+`$, where $`m(t)=sup_{S_E(t)}u`$. The maximum principle for the harmonic function $`u`$ implies that for $`t`$ sufficiently large, the function $`m(t)`$ is either monotone increasing or monotone decreasing in $`t`$ and hence approaches the value $`u^+`$ as $`t\mathrm{}`$. Consider the annuli $`A_\gamma (\tau )`$ in the scale $`g_t=t^2g`$, $`t=t(\gamma (\tau ))`$, as in (3.5). Any sequence $`\tau _i\mathrm{}`$ has a subsequence such that the corresponding annuli $`(A_\gamma (\tau _i),g_{t_i})`$ converge to a limiting domain $`A_{\mathrm{}}(\frac{1}{2},2)`$ in a tangent cone at infinity $`(A_{\mathrm{}},g_{\mathrm{}})`$, or $`(\overline{A}_{\mathrm{}},g_{\mathrm{}})`$, passing to covers as described above in the case of collapse. By construction, we then see that the potential function $`u_{\mathrm{}}`$ for this limit static vacuum solution achieves its maximal value $`u^+>0`$ at an interior point. Since $`u_{\mathrm{}}`$ is harmonic, the maximum priniple implies that $`u_{\mathrm{}}u^+`$, and hence by the vacuum equations (0.1), the limit $`(A_{\mathrm{}},g_{\mathrm{}})`$ or $`(\overline{A}_{\mathrm{}},g_{\mathrm{}})`$ is flat. This argument holds for any subsequence, and since the convergence to the limit is smooth, we see that
$$t^2r|_{A_\gamma (\frac{1}{2}t,2t)}0,t=t(\gamma (\tau )),\mathrm{as}\tau +\mathrm{},$$
by the scale-invariance of this expression.
It follows from Lemma 3.2 and Remark 3.3(iii) that the end $`E`$ is topologically standard, and the annuli $`A_E(\frac{1}{2}t,2t)`$, $`t=t(\gamma (\tau ))`$, are connected in $`E`$, for all $`\tau `$ sufficiently large. From the prior argument, this implies in particular that $`uu^+`$ uniformly at infinity in $`E`$.
However, as before, as $`\tau \mathrm{}`$, $`u(\gamma (\tau ))`$ is monotone decreasing to a value $`u_{}0`$. It follows that $`u^+=u_{}.`$ This is of course impossible, and shows that $`\gamma (\tau )`$ must have exited $`E`$ at some (negative) time.
Thus it remains to prove that $`Llimsup_Eu=u^+`$. Since the annuli $`A_E(\tau _j)`$ above are connected, the Harnack inequality (3.4), together with the fact that $`u^+<\mathrm{}`$, implies that $`L<\mathrm{}`$. Now choose points $`x_jS_\gamma (t_j)`$ such that $`u(x_j)L`$. As above, the functions $`u|_{A_E(\tau _j)}`$ have a subsequence converging to a limit harmonic function $`u_{\mathrm{}}`$ on a tangent cone at infinity based at $`x_{\mathrm{}}=limx_j`$. Then as before $`u`$ has an interior maximum at $`x_{\mathrm{}}`$ and hence $`u_{\mathrm{}}L`$; this gives $`L=u^+`$.
It follows from Lemma 3.4 and the discussion preceding it that any maximal flow line of $`u`$ intersects an apriori given large compact set $`KM`$, except those exceptional flow lines which start or end at a critical point of $`u`$ far out in $`M`$. In particular, a set of full measure in any given level set $`L`$ of $`u`$ may be connected to points in $`K`$ by flow lines of $`u.`$ In this sense, $`u`$ is ’almost proper‘, in that it behaves almost like a proper function in terms of the gradient flow.
Observe that this does not necessarily imply that the level sets of $`u`$ are compact, i.e. that $`u`$ is proper. For instance, the Weyl solution generated by the dipole potential $`\nu =\nu _++\nu _{}`$ considered in §2(IIC) satisfies (3.6), (it is even asymptotically flat), but the 0-level of $`\nu `$ is non-compact if $`\nu _+`$ and $`\nu _{}`$ are chosen so that the mass is 0. In this example, the only divergent flow lines of $`u`$ are the two ends of the $`z`$-axis.
Next, as in §2, let
$$\nu =logu.$$
The following result is quite standard.
###### Lemma 3.5.
On $`(N,g_N)`$ as in (0.2), with Riemannian metric, we have
$$\mathrm{\Delta }_N|\nu |0.$$
Proof: This standard estimate is a simple consequence of the Bochner-Lichnerowicz formula
$$\frac{1}{2}\mathrm{\Delta }|\nu |^2=|D^2\nu |^2+<\mathrm{\Delta }\nu ,\nu >+r(\nu ,\nu ),$$
on $`(N,g_N)`$, where we have dropped the subscript $`N`$ from the notation. Since $`\mathrm{\Delta }\nu =`$ 0, by (1.6), and since $`(N,g_N)`$ is Ricci-flat, this gives
$$\frac{1}{2}\mathrm{\Delta }|\nu |^2|D^2\nu |^2.$$
One computes
$$\mathrm{\Delta }|\nu |=\frac{1}{2}|\nu |^1\mathrm{\Delta }|\nu |^2\frac{1}{4}|\nu |^3||\nu |^2|^2,$$
and, by the Cauchy-Schwarz inequality $`||\nu |^2|^24|D^2\nu |^2|\nu |^2,`$ so the result follows.
Lemmas 3.4 and 3.5 lead to the following key result relating the behavior of $`|u|`$ to the area growth of geodesic spheres. This result is a straightforward consequence of the divergence theorem for proper harmonic functions $`u`$ on manifolds of non-negative Ricci curvature. Lemma 3.4 allows one to remove the assumption that $`u`$ is proper.
###### Lemma 3.6.
There is a constant $`C<\mathrm{}`$ such that for any component $`S_c(s)`$ of $`S(s)M`$, $`s1`$,
(3.8)
$$sup_{S_c(s)}|u|CareaS_c(s)^1.$$
Proof. We work on the Riemannian 4-manifold $`(N,g_N)`$ until the end of the proof. Let $`\widehat{A}_c=\widehat{A}_c(s)=\pi ^1(A_c(\frac{1}{2}s,2s))`$ and $`\widehat{S}_c=\widehat{S}_c(s)=\pi ^1(S_c(s)),`$ where $`\pi :NM`$ is projection on the first factor, with $`S^1`$ fibers. From the coarea formula, we have
(3.9)
$$_{\widehat{A}_c}|\nu |^2=_v_{L_v\widehat{A}_c}|\nu |𝑑\sigma _v𝑑v,$$
where $`L_v`$ is the $`v`$-level set of $`\nu `$ in $`N`$ and the outer integral in (3.9) is over the range of values in $``$ of $`\nu `$ in $`\widehat{A}_c.`$ Now as remarked following Lemma 3.4, up to a set $`\widehat{Z}_v`$ of measure 0 in $`L_v\widehat{A}_c,`$ all points in the set $`(L_v\widehat{A}_c)\widehat{Z}_v`$ may be joined by flow lines of $`\nu `$ to points in a fixed bounded hypersurface $`\widehat{T}`$ in $`\widehat{K}`$, independent of $`v`$, $`s`$; here $`K`$ is the compact set from Lemma 3.4, and $`\widehat{T}=\widehat{K}`$ for instance. Hence, by the divergence theorem applied to the harmonic function $`\nu `$ on $`N`$,
(3.10)
$$_{L_v\widehat{A}_c}|\nu |_{\widehat{T}}|\nu |c_1,$$
for some $`c_1<\mathrm{},`$ independent of $`s`$ and $`\widehat{A}_c`$.
Now $`|\nu |`$ is subharmonic on $`(N,g_N)`$ by Lemma 3.5, and $`diam_N^i\widehat{A}_ccs,`$ by (the proof of) Lemma 3.1. A standard sub-mean value inequality for manifolds of non-negative Ricci curvature, c.f. \[SY,Thm.II.6.2\], then gives
(3.11)
$$sup_{\widehat{S}_c(s)}|\nu |^2\frac{c_2}{vol\widehat{A}_c}_{\widehat{A}_c}|\nu |^2.$$
Hence the estimates (3.9)-(3.11) imply
(3.12)
$$sup_{\widehat{S}_c(s)}|\nu |^2c_3osc_{\widehat{A}_c}\nu (vol\widehat{A}_c)^1.$$
To estimate the right hand side of (3.12), again by (the proof of) Lemma 3.1, osc $`\nu c_4sup|\nu |s`$ on $`\widehat{A}_c.`$ Further, we claim that
$$sup_{\widehat{A}_c(s)}|\nu |csup_{\widehat{S}_c(s)}|\nu |,$$
for some constant $`c`$ independent of $`s`$ and $`\widehat{A}_c`$. To see this, by scale-invariance, it suffices to prove that $`sup_{\widehat{A}_c(1)}|\nu |csup_{\widehat{S}_c(1)}|\nu |`$ w.r.t. the rescaled metrics $`g_s=s^2g`$. By the curvature and diameter bounds (3.3) and (3.4) and Lemmas 1.3 and 1.4, the metrics $`(\widehat{A}_c(1),g_s)`$ form a compact family of metrics in the $`C^{\mathrm{}}`$ topology, unwrapping in the case of collapse. Thus, one has uniform control on the metrics $`g_s`$ on $`\widehat{A}_c(1)`$. Similarly, when normalized if necessary by additive and multiplicative constants so that $`sup_{\widehat{S}_c(1)}\nu =sup_{\widehat{S}_c(1)}|\nu |=1`$, the positive harmonic functions $`\nu `$ on $`(\widehat{A}_c(1),g_s)`$ also form a compact family of functions in the $`C^{\mathrm{}}`$ topology, i.e. a normal family. This follows by the Harnack estimate (3.4) and the Harnack principle (elliptic regularity) for harmonic functions, c.f. \[GT, Thm 2.11, Ch. 8\]. This compactness of the metrics and functions from elliptic theory proves the claim above.
Similarly, the metric compactness above also implies there is a constant $`c<\mathrm{}`$ such that
$$c^1areaS_c(1)volA_c(1)careaS_c(1),$$
w.r.t. the metrics $`g_s`$. (This estimate can also be derived directly from the Bishop-Gromov volume comparison theorem). Rescaling back to the metric $`g`$ then gives
(3.13)
$$c_5^1sareaS_c(s)volA_cc_5sareaS_c(s),$$
for some constant $`c_5<\mathrm{}.`$ Note that by definition,
(3.14)
$$area\widehat{S}_c(s)=_{S_c(s)}u𝑑A,$$
and the same for $`vol\widehat{A}_c.`$ Hence by (3.4), the estimate (3.13) holds also for $`\widehat{S}_c(s)`$ and $`\widehat{A}_c`$ in place of $`S_c(s)`$ and $`A_c`$.
Thus, by combining these estimates above, (3.12) gives
$$sup_{\widehat{S}_c(s)}|\nu |c_6area\widehat{S}_c(s)^1.$$
Using (3.14) and (3.4) again, this estimate implies (3.8).
We are now in a position to begin the proof of Theorem 0.3 itself. Observe that the previous results in §3 have not used the assumption (3.2), nor the assumption (ii) in Theorem 0.3 that $`u`$ does not approach 0 everywhere at infinity in $`E`$. Only the assumption that $`M`$ is pseudo-compact has been used. Hence, at this stage, we do not even know that $`E`$ has finite topological type. The main point initially is to prove that the estimates (3.6)-(3.7) above do hold on $`E`$ under these assumptions.
Recall that we have $`S_E(s)=S_c(s),`$ for $`S_c(s)E`$. Each geodesic ray $`\sigma (s)`$ in $`E`$, i.e. an integral curve of $`t,`$ with $`\sigma (s)S_E(s),`$ determines a component $`S_\sigma (s)=S_c(s)`$ s.t. $`\sigma (s)S_c(s);`$ the union of such components sweep out a part $`E_\sigma `$ of the end $`E`$. Of course $`E`$ is the union of $`E_\sigma `$ among all (non-homotopic) rays $`\sigma .`$
From Lemma 3.1 and from the obvious $`areaS_E(s)=areaS_c(s),`$ we have
$$^{\mathrm{}}areaS_E(s)^1𝑑s<\mathrm{}^{\mathrm{}}areaS_\sigma (s)^1𝑑s<\mathrm{},$$
for some geodesic ray $`\sigma E.`$ (Here the integrals start at some fixed value $`ss_o>`$ 0).
Hence, under the assumption (3.2), we have
(3.15)
$$^{\mathrm{}}areaS_\sigma (s)^1𝑑sK<\mathrm{},$$
for some ray $`\sigma E`$ and constant $`K`$.
By integrating along the curve $`\sigma ,`$ (3.15), Lemma 3.6 and the Harnack estimate (3.4) imply that $`u`$ is uniformly bounded in $`E_\sigma .`$ In fact, we claim that
(3.16)
$$u_{\mathrm{}}=lim_{t(x)\mathrm{}}u(x)<\mathrm{}$$
exists, where the limit is taken in $`E_\sigma .`$ To see this, let $`\gamma (s)`$ be any ”quasi-geodesic” in $`E_\sigma ,`$ i.e. $`\gamma `$ is a smooth curve with $`\gamma (s)S_\sigma (s)`$ and $`|d\gamma /ds|C_1,`$ for some $`C_1<\mathrm{}.`$ By (3.8) and (3.15), we then have $`^{\mathrm{}}|du(\gamma (s))|𝑑sCC_1K<\mathrm{},`$ and so $`u(\gamma (s_1))u(\gamma (s_2))`$ 0 whenever $`s_1,s_2\mathrm{}.`$ Hence the limit $`u_{\mathrm{}}(\gamma )`$ is well-defined. The diameter estimate (3.1) implies that all points in $`E_\sigma `$ lie on quasi-geodesics, (with a fixed $`C_1`$), in $`E_\sigma ,`$ starting on $`S_\sigma (s_o).`$ Further, the limit $`u_{\mathrm{}}(\gamma )`$ is clearly independent of $`\gamma ,`$ since for instance (3.8) and (3.15) imply that
$$osc_{A_\sigma (\frac{1}{2}s_j,2s_j)}u0,$$
on some sequence $`s_j\mathrm{}.`$ Hence (3.16) follows.
Next, since $`E`$ is an end, there exists some sequence $`t_j\mathrm{}`$ such that the geodesic spheres $`S_E(t_j)`$ are connected, and hence $`S_\sigma (t_j)=S_E(t_j).`$ By (3.16), $`u|_{S_E(t_j)}u_{\mathrm{}}`$ as $`t_j\mathrm{}.`$ The maximum principle applied to the harmonic function $`u`$ thus implies that $`u|_{A_E(t_j,t_k)}u_{\mathrm{}},`$ whenever $`t_j,t_k\mathrm{}.`$ Thus, we see that (3.16) holds where the limit is taken in the full end $`E`$, and not just in $`E_\sigma .`$
Now we use the assumption (ii) of Theorem 0.3, which says that $`u(x_j)u_o>`$ 0, for some constant $`u_o>`$ 0 and some divergent sequence $`x_jE`$. It follows from this and the existence of the limit (3.16) in $`E`$ that
$$u_{\mathrm{}}>0.$$
Hence, we may, and will, renormalize the potential function $`u`$ of the static vacuum solution $`(M,g,u)`$ so that, on $`E`$,
(3.17)
$$lim_{t(x)\mathrm{}}u(x)=1.$$
The estimate (3.17) essentially immediately implies the scale-invariant estimates
(3.18)
$$sup_{S_E(s)}|r|<<s^2,sup_{S_E(s)}|u|<<s^1,\mathrm{as}s\mathrm{},$$
strenthening the bounds (3.3). For as discussed following (3.3), (3.18) is equivalent to the statement that all tangent cones at infinity of $`E`$ are flat, (c.f. also the proof of Lemma 3.2). But, as noted above, all tangent cones at infinity are static vacuum solutions, and (3.17) implies that the limit potential $`u_{\mathrm{}}`$ satisfies $`u_{\mathrm{}}`$ 1. Hence, the static vacuum equations (0.1) of course imply the limit metrics are flat.
Lemma 3.2 now determines the topology of the end $`E`$, as one of two (standard) alternatives, according to non-collapse or collapse behavior at infinity.
Before proceeding to the analysis of these cases, note that (3.18) implies that the metrics $`g`$ and $`\stackrel{~}{g}=u^2g`$ from (1.4) are quasi-isometric on $`E`$, and almost isometric near infinity. Since $`\stackrel{~}{g}`$ has non-negative Ricci curvature on $`E`$, standard volume comparison theory implies that the area and volume ratios
(3.19)
$$\frac{area_{\stackrel{~}{g}}(\stackrel{~}{S}(s))}{s^2},\frac{vol_{\stackrel{~}{g}}(\stackrel{~}{B}(s))}{s^3}$$
are monotone non-increasing in $`s`$. Hence, their limits at $`s=\mathrm{}`$ exist, and by the equality of $`g`$ and $`\stackrel{~}{g}`$ at infinity, the limits at $`s=\mathrm{}`$ of the ratios in (3.19) w.r.t. the $`g`$ metric and $`g`$-geodesic spheres and balls also exist, and equal the $`\stackrel{~}{g}`$ limits.
To proceed further, we now separate the discussion into non-collapse and collapse cases.
Case A. (Non-Collapse).
Suppose that $`E`$ is non-collapsing at infinity, i.e. by the remarks above,
(3.20)
$$limsup_s\mathrm{}\frac{v(s)}{s^3}=lim_s\mathrm{}\frac{v(s)}{s^3}>0,$$
where $`v(s)`$ denotes $`vol_g(B_E(s)).`$ This implies, via (3.13), that $`areaS(s)^1cs^2,`$ and hence, by Lemma 3.6,
(3.21)
$$sup_{S(s)}|u|cs^2.$$
The volume condition (3.20) implies that all tangent cones at infinity $`A_{\mathrm{}}`$ of $`(E,g)`$ exist, without passing to covering spaces, and, as discussed above, are flat solutions of the static vacuum equations with potential $`u_{\mathrm{}}1`$. As discussed in the proof of Lemma 3.2, the tangent cone at infinity is unique, (up to isometry), and given by $`^3\{0\}`$, and further $`EK`$ is diffeomorphic to $`^3B`$, for some compact set $`KE`$. The blow-down metrics $`g_s=s^2g`$ on all annuli $`A_E(\frac{1}{2}s,2s)`$ converge smoothly to the flat metric on $`A(1,2)^3`$, uniformly as $`s\mathrm{}`$, and hence there are local (harmonic) coordinates on $`A_E(\frac{1}{2}s,2s)`$ in which $`g`$ has the expansion $`g_{ij}=\delta _{ij}+\gamma _{ij}`$, where $`|\gamma _{ij}(x)|`$ 0 uniformly as $`s\mathrm{}`$. Again as discussed in the proof of Lemma 3.2, these local coordinates, e.g. on $`\{A_E(2^{i1},2^{i+1})\}`$, $`i>0`$, may be assembled into a global chart, mapping $`EK`$ onto $`^3B`$, c.f. \[AC\] for further details if desired. With respect to such a chart, the metric $`g_{ij}`$ has the form
(3.22)
$$g_{ij}=\delta _{ij}+\gamma _{ij},$$
on all of $`EK`$, with $`|\gamma _{ij}|0`$ uniformly at infinity in $`E`$. In other words, $`g`$ is $`C^o`$ asymptotic to the flat metric at infinity.
To prove that the metric on $`E`$ is (strongly) asymptotically flat, as defined preceding Theorem 0.1, consider again the metric $`\stackrel{~}{g}=u^2g.`$ From (1.4) and (3.21), the curvature of $`\stackrel{~}{g}`$ decays as
$$|\stackrel{~}{r}|Ct^4,$$
as $`t\mathrm{}.`$ (Note also that $`t`$ and $`\stackrel{~}{t}`$ are approximately equal for $`t`$ large).
It follows that the expansion (3.22) may be improved, for $`\stackrel{~}{g},`$ to
(3.23)
$$\stackrel{~}{g}_{ij}=\delta _{ij}+O(t^2),$$
in a suitable (harmonic) coordinate chart. We refer to \[BKN\] or \[BM\] for instance for further details here. Briefly, elliptic regularity theory applied to the equations (1.4)-(1.5), together with the curvature decay above, implies that the $`2^{nd}`$ derivatives of the metric $`\stackrel{~}{g}`$ in the coordinate chart decay as $`t^4,`$ so that the metric $`\stackrel{~}{g}`$ decays to the flat metric at a rate of $`t^2.`$ Hence
(3.24)
$$g_{ij}=u^2\stackrel{~}{g}_{ij}=(1+2\upsilon )\delta _{ij}+O(t^2),$$
where $`\upsilon =1u.`$ Here we are using that fact that since $`|u|=O(t^2),`$ $`u`$ has an expansion of the form $`u=`$ 1 $`+O(t^1).`$ Further, since log $`u`$ is harmonic w.r.t. $`\stackrel{~}{g},`$ the decay (3.21) and (3.23) implies that $`\mathrm{\Delta }_flogu=O(t^4)`$ for $`t`$ large, where $`\mathrm{\Delta }_f`$ is the flat Laplacian on $`^3.`$ This means that $`u`$ has an expansion $`u=`$ 1 $`\frac{m}{t}+O(t^2),`$ where $`m`$ is the mass of $`E`$ defined in (0.7).
In particular, these estimates show that the end $`E`$ is asymptotically flat in the sense preceding Theorem 0.1.
Note that, to first order in $`t^1,`$ the function $`\upsilon =1u`$ corresponds to the Green’s function in $`^3,`$ i.e. the fundamental solution of the Laplacian, weighted by the mass $`m`$. It is of course possible to have $`m=`$ 0, as for instance for the dipole-type Weyl solutions in §2(IIC), or also $`m<0`$. Further, since $`u`$ has been normalized so that $`u1`$ at infinity in $`E`$, the expression (0.7) for the mass is equivalent to the usual definition
(3.25)
$$m_E=\frac{1}{4\pi }_{S_E(s)}<u,t>dA,$$
where $`s`$ is sufficiently large so that $`S_E(s)E=\mathrm{}`$. This is because the expression (3.25) is independent of $`s`$, since $`u`$ is harmonic, and the fact that it is asymptotic to the expression (0.7) as $`s\mathrm{}`$.
This completes the analysis of Case A.
Case B. (Collapse).
Under the standing assumption (3.15), suppose that the end $`E`$ is collapsing at infinity, i.e.
(3.26)
$$limsup_s\mathrm{}\frac{v(s)}{s^3}=lim_s\mathrm{}\frac{v(s)}{s^3}=0.$$
We will prove that this situation is impossible. The results preceding Case A remain valid, so that (3.17)-(3.18) hold, all tangent cones at infinity $`\overline{A}_{\mathrm{}}`$ of $`E`$ are flat products of the form $`V\times S^1`$ or $`^+\times S^1\times S^1`$, where $`V`$ is a flat 2-dimensional cone. Further $`EK`$ is diffeomorphic to $`(^2B)\times S^1`$, for some compact set $`KE`$.
By (3.15) and Lemma 3.6, we have
(3.27)
$$^{\mathrm{}}sup_{S(s)}|u|(s)𝑑sK_1<\mathrm{}.$$
The main point is now to show that $`(E,g)`$ itself, (and not just its tangent cones), is asymptotic to a flat quotient of $`^3,`$ and hence has at most quadratic volume growth of geodesic balls or linear area growth of geodesic spheres. This is done in the following result, which is a strengthening of Lemma 3.2.
###### Lemma 3.7.
Under the assumptions (3.26) and (3.27) above, there is a compact set $`KE`$ such that $`(EK,g)`$ is quasi-isometric to a flat product $`(^2B)\times S^1`$ or $`^+\times S^1\times S^1.`$
Proof: As in Case A, it is useful to work with the metric $`\stackrel{~}{g}=u^2g;`$ again, this makes no significant difference, since (3.17) holds. All the metric quantities below are thus in the $`\stackrel{~}{g}`$ metric. For notational simplicity however, we drop the tilde from the notation.
Let $`t_{\mathrm{}}(x)=lim_s\mathrm{}(dist(x,S_E(s))s)`$. As in the construction of Busemann functions, the limit here exists, c.f. \[Wu\] for a discussion of such functions. By construction, $`t_{\mathrm{}}`$ is a Lipschitz distance function, i.e. $`t_{\mathrm{}}`$ realizes everywhere the distance between its level sets. Observe that on $`^3,t_{\mathrm{}}`$ is just the distance function to $`\{0\}^3,`$ on $`V\times S^1,t_{\mathrm{}}`$ is the distance function to $`\{0\}V`$ pulled back to $`V\times S^1,`$ for any cone $`V`$ with vertex {0}, while on $`^+\times S^1\times S^1,t_{\mathrm{}}`$ is the distance function on $`^+`$ pulled back to the total space.
By renormalization, (as with the potential $`u`$), $`t_{\mathrm{}}`$ induces a distance function $`\overline{t}_{\mathrm{}}`$ on each tangent cone $`\overline{A}_{\mathrm{}}`$ by defining $`\overline{t}_{\mathrm{}}(x)=lim(t_{\mathrm{}}(x)/t_{\mathrm{}}(x_i)),`$ where $`t(x_i)\mathrm{}`$ and $`x_i`$ are the base points converging to the base point $`x_{\mathrm{}}\overline{A}_{\mathrm{}}.`$ Thus $`\overline{t}_{\mathrm{}}`$ is the function above on $`V\times S^1`$ or $`^+\times S^1\times S^1.`$
The map $`t_{\mathrm{}}:(E,\stackrel{~}{g})^+`$ is distance non-increasing, and preserves distance along the integral curves of $`t_{\mathrm{}};`$ thus where smooth, it is a Riemannian submersion. We will show that $`t_{\mathrm{}}`$ gives rise to a Lipschitz quasi-isometry by examining the asympotics of the second fundamental form of its level sets.
Thus, let $`\sigma (s)`$ be any geodesic ray in $`E`$ which is an integral curve of $`t_{\mathrm{}},`$ and let $`B=B_\sigma (s)`$ denote the second fundamental form of the level surface $`t_{\mathrm{}}^1(s)`$ at $`\sigma (s).`$ The form $`B`$ is well-defined and smooth along any such ray $`\sigma `$. Recall that $`B`$ satisfies the Riccati equation
(3.28)
$$B^{}+B^2+R_T=0,$$
where $`T`$ is the unit tangent vector along $`\sigma .`$ Consider the behavior of $`sB(s)`$ as $`s\mathrm{}.`$ This quantity is scale-invariant, and thus converges smoothly, (in subsequences), to the limiting expression $`\overline{s}B_{\mathrm{}}(\overline{s})`$ on any tangent cone at infinity. Since the parameters $`s`$ and $`t_{\mathrm{}}|_\sigma `$ are the same up to additive constants, $`\overline{s}=\overline{t}_{\mathrm{}}.`$ Similarly, by the definition of $`\overline{t}_{\mathrm{}},B_{\mathrm{}}`$ is the second fundamental form of the levels $`\overline{t}_{\mathrm{}}.`$ Hence, either
(i): $`\overline{s}B_{\mathrm{}}(\overline{s})=(d\theta /|d\theta |)^2,`$ when the tangent cone is of the form $`V\times S^1,`$ and $`\theta `$ is the angle variable about $`\{0\}V,`$ or
(ii): $`\overline{s}B_{\mathrm{}}(\overline{s})=`$ 0, when the tangent cone is of the form $`^+\times S^1\times S^1.`$
Thus $`\overline{s}B_{\mathrm{}}(\overline{s})`$ is either of rank 1, with eigenvalue 1, or identically 0. Note that the expression in case (i) is independent of the cone $`V`$, i.e. the cone angle at $`\{0\}`$.
As noted preceding Lemma 3.2 and in Remark 3.3(ii), the geometry of $`\overline{A}(\frac{1}{2}s,2s)`$ smoothly approximates that of a limit tangent cone $`\overline{A}_{\mathrm{}},`$ for $`s`$ large, and varies continuously in $`s`$. Since the two alternatives (i) and (ii) above for the structure of $`B_{\mathrm{}}`$ are rigid, it follows that all tangent cones $`\overline{A}_{\mathrm{}}`$ are of the same type, i.e. they are all of the form $`V\times S^1,`$ or all of the form $`^+\times S^1\times S^1.`$
The main task now is to show that the deviation of $`sB(s)`$ from its limit $`\overline{s}B_{\mathrm{}}(\overline{s})`$ has bounded integral. To do this, we use the Riccati equation, and estimate the decay of the curvature term $`R_T,`$ using basically standard methods in comparison geometry, c.f. \[P, Ch.6.2\] for instance.
Thus, from (1.4), the sectional curvature $`K`$ of $`(M,\stackrel{~}{g})`$ satisfies
$$K_{XZ}=|\nu |^20,K_{XY}=|\nu |^20,$$
where $`Z=u/|u|,`$ $`X,Y`$ are vectors orthogonal to $`Z`$, and $`\nu =`$ log $`u`$. Hence $`|R_T||\nu |^2.`$ Substituting this in (3.28) gives
(3.29)
$$|B^{}+B^2||\nu |^2.$$
Let $`\lambda `$ be any eigenvalue of $`B`$, with unit eigenvector $`e`$; (note that $`B`$ is symmetric). Observe then that $`s\lambda (s)`$ converges either to 1 or to 0, as $`s\mathrm{}.`$ The estimate (3.29) when applied to $`(e,e)`$ gives
$$\pm (s\lambda ^{}+s\lambda ^2)s|\nu |^2.$$
Integrate this by parts along any finite interval $`I`$ to obtain
$$|s\lambda |_I+_I(s\lambda ^2\lambda )ds|_Is|\nu |^2ds.$$
If $`s\lambda 0`$ as $`s\mathrm{}`$, choose $`I`$ to be any interval on which $`s\lambda =0`$ at $`I`$ and $`s\lambda 0`$ on $`I`$, so that $`s\lambda `$ has a definite sign on $`I`$. If there are no such boundary points, choose $`I`$ to be an infinite half-line. Similarly, if $`s\lambda 1`$, choose $`I`$ to be intervals such that $`s\lambda =1`$ at $`I`$ with $`s\lambda 10`$ on $`I`$. Then summing up the estimate above over all such intervals gives
$$_\sigma |s\lambda ^2\lambda |𝑑s_\sigma s|\nu |^2𝑑s+C_o,$$
for some constant $`C_o<\mathrm{}`$. Now the estimate (3.18), together with (3.27) gives
$$s|\nu |^2𝑑sC_1,$$
for some constant $`C_1<\mathrm{}.`$
Suppose first $`s\lambda (s)`$ 1 as $`s\mathrm{}.`$ We then obtain
$$|s\lambda (\lambda \frac{1}{s})|𝑑sC_2<\mathrm{},$$
and hence
(3.30)
$$|\lambda \frac{1}{s}|𝑑sC_3<\mathrm{}.$$
Similarly, if $`s\lambda (s)`$ 0 as $`s\mathrm{},`$ one obtains
(3.31)
$$|\lambda |𝑑sC_3<\mathrm{}.$$
Now the second fundamental form $`B`$ gives the logarithmic derivative of the norm of Jacobi fields formed by the family of $`t_{\mathrm{}}`$-rays in $`E`$ starting at some level $`t_{\mathrm{}}^1(s_o).`$ Thus, if $`J`$ is any Jacobi field formed from the $`t_{\mathrm{}}`$-conguence, and $`v=J/|J|,`$ we have along any $`t_{\mathrm{}}`$-ray,
$$B(v,v)=\frac{d}{ds}(log|J|(s)).$$
Suppose first that the end $`E`$ has tangent cones at infinity of the form $`^+\times S^1\times S^1.`$ Then (3.31) implies the uniform bound
(3.32)
$$C_4^1|J|(s)C_4,$$
with $`C_4=e^{C_3}.`$ This means that the geometry of the level surfaces of $`t_{\mathrm{}}`$ is uniformly bounded as $`s\mathrm{},`$ i.e. the diameter and area of the level surfaces is uniformly bounded away from 0 and $`\mathrm{}.`$ It is then clear that there is a Lipschitz quasi-isometry of $`(EK,\stackrel{~}{g})`$ to $`^+\times S^1\times S^1`$ induced by $`t_{\mathrm{}}.`$ Since $`\stackrel{~}{g}`$ and $`g`$ are also quasi-isometric, by (3.17), the lemma follows in this case.
If $`E`$ has tangent cones at infinity of the form $`V\times S^1,`$ then there is a basis of Jacobi fields whose elements satisfy either (3.32), or, from (3.30),
(3.33)
$$C_4^1s|J|(s)C_4s.$$
As before, this implies that $`t_{\mathrm{}}`$ gives rise to a quasi-isometry of $`(EK,g)`$ to $`(^2B)\times S^1.`$
Of course Lemma 3.7, in both cases, immediately implies that
(3.34)
$$areaS_E(s)cs,$$
for some constant $`c<\mathrm{}.`$ However, (3.34) violates the standing assumption (3.15), (or (3.2)). It follows that no end $`(E,g)`$ can satisfy the assumptions of Case B.
Together with Case A, this completes the proof of the second statement in Theorem 0.3.
We now turn to the proof of the last statement in Theorem 0.3. We will assume that the end $`E`$ is small, i.e.
(3.35)
$$_E\frac{1}{areaS_E(s)}𝑑s=\mathrm{},$$
and derive a contradiction from the assumptions $`supu<\mathrm{}`$ and $`m_E0`$.
The proof is based on a result of Varopoulos \[V\] which states that ends of Riemannian manifolds satisfying (3.35) are parabolic, i.e. admit no non-constant positive superharmonic functions $`v`$ which tend uniformly toward their infimum at infinity. (Actually, the result in \[V\] is a condition on the volume growth of geodesic balls, but this is equivalent to the bound (3.35) under the estimate (3.13)). We will prove that the potential $`u`$ is such a non-constant function, giving the required contradiction.
Thus, suppose first that
(3.36)
$$sup_Eu<\mathrm{},$$
Arguing as in (3.9), but now on $`EM`$ in place of $`N`$, we have
(3.37)
$$_E|u|^2=_v_{L_vE}|u|𝑑A_v𝑑v,$$
and as in (3.10),
$$_{L_vE}|u|_T|u|C.$$
Thus, these estimates imply that
(3.38)
$$_E|u|^2<\mathrm{},$$
so that
(3.39)
$$_{A_E(s,\mathrm{})}|u|^20\mathrm{as}s\mathrm{}.$$
On the other hand, again referring to the proof of Lemma 3.6, since almost all (in terms of measure) points in $`L_vE`$ may be joined by flow lines of $`u`$ to a fixed bounded surface $`T`$ in $`K`$, by the divergence theorem there is a subsurface $`T^{}T`$ such that
(3.40)
$$_{L_vE}|u|_T^{}|u|c,$$
for some constant $`c>`$ 0. It follows from (3.37)-(3.40) that
(3.41)
$$osc_{A_E(s,\mathrm{})}u0\mathrm{as}s\mathrm{}.$$
By assumption (ii) in Theorem 0.3, we may thus assume w.l.o.g. that
(3.42)
$$lim_t\mathrm{}u=1,$$
in $`E`$. Thus, as noted in (3.18), $`|r|\epsilon (t)/t^2,`$ where $`\epsilon (t)`$ 0 as $`t\mathrm{}`$, and so Lemma 3.2 holds on $`E`$.
Now choose a smooth approximation $`S`$ to a large geodesic sphere $`S_E(s)`$ with $`S_E(s)E=\mathrm{}`$, so that $`S`$ separates $`E`$ into two components, one being the outside containing the end $`E`$.
We claim that if, in addition to (3.36),
(3.43)
$$m_E0,$$
then there is a set of flow lines $`\gamma `$ of $`u`$ or $`u`$, starting on $`S`$, of positive measure on $`S`$, and pointing out of $`S`$, which never intersects $`S`$ again at later times. Hence such $`\gamma `$ diverge to infinity in $`E`$, since up to a set of measure 0, $`\gamma `$ does not terminate in a critical point of $`u`$. (Compare with the earlier discussion regarding divergent flow lines and Weyl dipole solutions concerning Lemma 3.4).
To see this, suppose instead that all flow lines say of $`u`$ which initially point out of $`S`$ eventually intersect $`S`$ again, with the exception of those terminating in critical points. Consider the measure
(3.44)
$$d\mu =<u,\nu >dA,$$
on $`S`$, where $`dA`$ is the Lebesgue measure and $`\nu `$ is the unit outward normal on $`S`$; $`d\mu `$ is absolutely continuous w.r.t. $`dA`$. Since $`u`$ is harmonic, the divergence theorem implies that the gradient flow of $`u`$ preserves the measure $`d\mu `$ in the following sense. Let $`D`$ be a domain in $`S`$ and let $`\mathrm{\Omega }`$ be the domain in $`E`$ formed by a collection of flow lines outside $`S`$, whose endpoints form a smooth surface $`D^{}`$. If $`\nu ^{}`$ denotes the unit outward normal to $`\mathrm{\Omega }`$ at $`D^{}`$, then the flow from $`D`$ to $`D^{}`$ carries the measure $`d\mu `$ to the measure $`d\mu ^{}=<u,\nu ^{}>dA^{}`$, where $`dA^{}`$ is Lebesgue measure on $`D^{}`$. In particular, the flow preserves the masses of the measures.
Thus, under the assumptions above, the gradient flow (with varying flow-times), induces a homeomorphism of $`SZ`$ into itself, where $`Z`$ is a set of Lebesgue measure 0, corresponding to flow lines terminating in critical points. However, this homeomorphism inverts the direction of $`u`$ w.r.t. the fixed normal $`\nu `$, and hence maps domains $`D_+`$ on which the measure $`d\mu `$ is positive onto domains $`D_{}`$ on which $`d\mu `$ is negative, in such a way that that
(3.45)
$$m_\mu (D_+)=|m_\mu (D_{})|.$$
This of course implies that the total mass of $`d\mu `$ on $`S`$ is 0. However using (3.44) and the remarks concerning (3.25), the mass $`m_E`$ of $`E`$ equals the total mass of $`d\mu `$. This contradiction proves the claim.
We may now complete the proof as follows. Assuming $`E`$ is an end satisfying (3.36) and (3.43), there exists an open set $`𝒪`$ of flow lines $`\gamma =\gamma (\tau )`$ of either $`u`$ or $`u`$ which start on a set of positive measure on $`S`$ and diverge to infinity in $`E`$. Consider the former case, which corresponds to $`m_E>0`$. A generic flow line of $`u`$ in $`E`$ tends to the maximal value of $`u`$ in $`E`$ and hence a generic flow line in $`𝒪`$ also tends this maximal value. By (3.42), it then follows that sup$`{}_{EK}{}^{}u=1`$, for some compact set $`KE`$. The function $`v=u`$ is thus a bounded (non-constant) harmonic function on $`E`$, which tends uniformly to its infimum at infinity. Hence $`E`$ cannot be parabolic. This contradiction shows that (3.35) cannot hold for $`E`$, and thus, by the proofs in Cases A and B above, the end $`E`$ is asymptotically flat. The proof in case $`m_E<`$ 0 is the same.
This completes the proof of Theorem 0.3. ∎
Remark 3.8.(i). There are (non-flat) static vacuum solutions with a small end, namely the Kasner metrics (2.33), with $`a>`$ 0. These solutions have $`volB(s)s^{2\delta },`$ $`volS(s)s^{1\delta },`$ where $`\delta =(a+a^11)^1(0,1)`$, and hence the end is small. Note that $`us^\delta `$ is unbounded. There are other Weyl solutions which are complete away from $`\mathrm{\Sigma }`$ with $`M`$ pseudo-compact and with one small end, (take for instance the potential given by the Green’s function on $`^2\times S^1`$, see (ii) below), but all known examples with small ends are either asymptotic to the Kasner metric at infinity or have faster than quadratic curvature decay, i.e. satisfy (3.18).
It is an open problem to understand in more detail the structure of small ends of static vacuum solutions. It follows from the results above in §3 that all tangent cones at infinity of $`E`$ are collapsing, and hence they are all Weyl solutions. But the metric uniqueness of tangent cones at infinity is unknown, as is the question of whether small ends have finite topological type.
(ii). We construct an example which illustrates the sharpness of the last statement of Theorem 0.3. Let $`G_1`$ and $`G_2`$ be the Green’s functions for the Laplacian on the flat product $`^2\times S^1`$, with poles at $`(0,p_i)`$, for $`p_1`$,$`p_2`$ distinct points in $`S^1`$. Here we consider $`S^1`$ as the $`z`$-axis in $`^3`$ quotiented out by an isometric $``$-action. As in (2.34), $`G_i(x)=G(x,p_i)`$, viewed as a function on the universal cover $`^3`$, may be written as
(3.46)
$$G_i=lim_n\mathrm{}\left[\underset{j=n}{\overset{n}{}}\frac{1}{r_j}c_n\right],$$
where $`r_j`$ is the Euclidean distance to the collection of lifts of $`p_i`$ in $`^3`$ and $`c_n`$ is a suitable normalizing constant with $`c_n\mathrm{}`$ as $`n\mathrm{}`$, chosen so that $`G(x,p_i)`$ is finite. Thus $`G_i`$ is an axisymmetric and $`z`$-periodic harmonic function on $`(^2)^+\times S^1`$, where the $`S^1`$ now means rotations about the $`z`$-axis in $`^3`$. Hence the potential
$$\nu =G_1G_2$$
generates a Weyl solution as in (2.12), with $`u=e^\nu `$ and which has an isometric $``$-action along the $`z`$-axis. Let $`M`$ denote the $``$ quotient of this solution. Then the metric boundary of $`M`$ is pseudo-compact. Since $`G_ilogr`$ as $`r\mathrm{}`$,
$$u1,$$
at infinity in $`M`$. The end $`E`$ = $`^+\times T^2`$ is small and has mass 0 in the sense of (0.7). These solutions of course resemble the dipole-type solutions discussed in §2(IIC), but with a collapsing end.
Remark 3.9.(i). Although we will not detail it here, an examination of the proof shows that Theorem 0.3 holds for non-vacuum static solutions of the Einstein field equations (1.2), provided suitable decay conditions are imposed on the stress-energy tensor $`T`$.
This is the case for example, if $`T`$ is any tensor with compact support, or more generally if $`T`$ satisfies an estimate of the form $`|T|(x)ct^3(x),`$ for some $`c>`$ 0 and all $`x`$ with $`t(x)s_o`$, together with $`\frac{1}{u}\mathrm{\Delta }uL^1(MU).`$ This latter condition is needed to obtain the bound (3.10). By (1.2), note that $`\frac{1}{u}\mathrm{\Delta }u=\frac{1}{2}trT.`$ The starting estimate (3.3) may be obtained by applying (a suitable version of) \[An1,Thm.3.3\].
(ii). Also, Theorem 0.3 can be given a finite or quantitative formulation, i.e. one can relax the assumption of completeness, in basically the same way as the local estimates (3.3) follow from the non-existence of global static vacuum solutions with $`M=\mathrm{}`$, c.f. \[An1,App.\].
Thus, if (3.15) holds and $`(M,g)`$ is ’sufficiently large‘, depending on $`K`$, then sufficiently far out in $`(M,g)`$, the metric is close to a flat metric. We leave a precise formulation and proof, (based on Theorem 0.3), to the reader.
Remark 3.10. As noted in §0, it would be of interest to prove that $`(M,g,u)`$ has a unique end. Under the hypotheses of Theorem 0.3, we conjecture this is the case at least when $`M`$ is complete away from $`\mathrm{\Sigma }.`$ If $`M`$ is in addition smooth up to $`\mathrm{\Sigma },`$ this has been proved by Galloway \[G\]. Following essentially the same arguments as in \[G\], it is not difficult to show that if $`M`$ is complete away from $`\mathrm{\Sigma }`$ and if the Riemannian 4-manifold $`N`$ admits a compact smoothing of a neighborhood of $`\mathrm{\Sigma }`$ having non-negative Ricci curvature, then $`M`$ has a unique end.
Remark 3.11. We point out that Theorem 0.3 and Theorem 0.1 are false in higher dimensions, due to the existence of Einstein metrics on compact manifolds, which are not of constant curvature, in dimensions $``$ 3. (The equations (0.1) on any $`n`$-dimensional manifold $`M^n`$ generate Ricci-flat manifolds on $`N^{n+1}).`$
Thus, let $`(\mathrm{\Sigma },g)`$ be any compact $`(n2)`$ dimensional Einstein manifold, with $`Ric_g=(n3)g`$ and define the warped product metric $`\overline{g}`$ on $`^2\times \mathrm{\Sigma }`$ by
$$\overline{g}=dt^2+\frac{4(f^{}(t))^2}{(n2)^2}d\varphi ^2+f^2(t)g,$$
where $`f`$ is the unique function on $`[0,\mathrm{})`$ such that $`f(0)=`$ 1, $`f^{}>`$ 0 and $`(f^{})^2=`$ 1 $`f^{1n}.`$ A simple computation shows that $`(^2\times \mathrm{\Sigma },\overline{g})`$ is complete and Ricci-flat, c.f. \[Bes, p.271\] and the space-like hypersurface $`^+\times \mathrm{\Sigma },`$ with metric $`dt^2+f^2(t)g,`$ is a solution to the static vacuum equations, with potential $`u=f^{},`$ (up to a multiplicative constant). Thus the horizon is $`\mathrm{\Sigma }`$ and the solution is smooth up to $`\mathrm{\Sigma }`$ and complete away from $`\mathrm{\Sigma }.`$
This metric is asymptotically conical, i.e. asymptotic to the complete (Euclidean) cone on $`(\mathrm{\Sigma },g)`$, but is asymptotically flat only in the case that $`(\mathrm{\Sigma },g)=S^{n2}(1),`$ corresponding to the $`n`$-dimensional Schwarzschild metric.
July, 1998: revision July, 2000 |
warning/0001/quant-ph0001068.html | ar5iv | text | # Quantum Decoherence from Adiabatic Entanglement
## I Introduction
In quantum measurement process, wave packet collapse (WPC, also called von Neumann’s projection or wave function reduction ) physically resembles the disappearance of interference pattern for Young’s two-slit experiment in the presence of a “which-way” detector. Associated with the wave-particle duality, this phenomenon of losing quantum coherence is referred to as the so-called quantum decoherence. In fact, before a measurement to observe “which-way” the particle actually takes, the quantum particle seems to move from a point to another along several different ways simultaneously. This just reflects the wave feature of a quantum particle. The detection of “which-way” means a probe for the particle feature, which leads to the disappearance of wave feature or quantum decoherence.
Further explanation for decoherence phenomena was made in the view point of complementarity by Niels Bohr based on Heisenberg’s position-momentum uncertainty : A particular interaction between a classical instrument (detector) and the measured quantum system can be regarded as a quantum measurement; but once enough data about the states of the quantum system is “read out” from the motion configuration of the detector, the interaction unavoidably destroys the interference pattern. According to Heisenberg’s uncertainty principle, to locate the position of a particle to the uncertainty of order $`\mathrm{\Delta }x`$ along the direction orthogonal to its moving direction, the “which-way” measurement must kick the momentum of the particle to an uncertainty of order $`1/\mathrm{\Delta }x`$ and thereby washes out the spatial interference pattern. (In this paper the Plank constant is taken to be unity). Bohr’s argument sounds correct, but a recent experiment on Brrag’s reflection of cold atoms shows that Schrodinger’s concept of entangled state, rather than the unavoidable measurement distribution, is crucial for the wave-particle duality in this “which-way” experiment. Another “which-way” experiment , which uses the electron Aharonov-Bohm interference with a quantum point contact, also manifests the importance of quantum entanglement. Actually, similar gedenken experiments using photon and neutron have been considered before \[6-7\]
A quantum entangled state \[8-10\] such as
$$|\mathrm{\Psi }=\underset{n}{}C_n|S_n|D_n(|S|D$$
$`(1.1)`$
for any $`|S`$ and $`|D`$ is a coherent superposition of states of a quantum system of many particles or of a single particle with many degrees of freedom. It involves a correlation between the states $`|S_n`$ of the quantum system and the states $`|D_n`$ of the detector. Once the detector is found in a state $`|D_n,`$ the total system must collapse into a certain component $`|S_n|D_n`$. Then one can infer the state $`|S_n`$ of the quantum system. The interference pattern can be obtained from the total wave function $`|\mathrm{\Psi }`$ by “summing over” all possible states of the detector. Assuming the states of the detector are normalized, we have
$$\underset{m}{}|D_m|x|\mathrm{\Psi }|^2=\underset{n}{}|C_n|^2|S_n(x)|^2+\underset{nm}{}C_m^{}C_nS_m^{}(x)S_n(x)D_m|D_n$$
$`(1.2)`$
where $`S_n(x)=x|S_n`$ is the state of the quantum system in position representation. The second term on the r.h.s of the above equation is responsible for the interference pattern. It is easy to see the interference fringes completely vanish when the states of the detector are orthogonal to one another , i.e., when $`D_m|D_n=\delta _{m.n}`$. In this situation, an ideal quantum measurement results from the ideal entanglement with the orthogonal correlated components $`|D_n`$ , in which one can distinguish the state of detector very well. Mathematically, by using the reduced density matrix
$$\rho =\text{Tr}_D(|\mathrm{\Psi }\mathrm{\Psi }|)=\underset{n}{}|C_n|^2|S_nS_n|+\underset{mn}{}C_m^{}C_n|S_nS_m|D_m|D_n$$
$`(1.3)`$
which is obtained by tracing out the variables of detector, the above-mentioned decoherence phenomenon can be equivalently expressed as a projection or reduction of the reduced density matrix from a pure state $`\rho =_{m,n}|S_nS_m|`$ to a mixed state $`\widehat{\rho }=_n|S_nS_n|.`$
It is noticed that, so long as the “ which-way” information already stored in the detector could be read out, the interference pattern has been destroyed without any data read out in practice . In this sense the environment surrounding the quantum system behaves as a detector to realize a “measurement-like ” process. This is because the environment never needs to read out the data. Thus, the above argument is also applicable to the analysis of decoherence problem of an interfering quantum system coupling to the environment \[8-10\].In this kind of problems, the environment is imagined as an objective detector detecting the states of the quantum system and thereby the detector states $`|D_n`$ are thought to be the macroscopic quantum states of the environment. Provide the environment couples with the quantum system and produce an ideal entanglement, the quantum system must lose its coherence. It is worthy to point out that this simple entanglement conserves the energy of the quantum system while destroying the quantum coherence. The loss of energy of the quantum system can be separately discussed in the quantum dissipation theory well developed in recent years \[11-16\]
In our previous works on quantum measurement theory\[17-23\] , we investigate how an ideal entanglement appears in the macroscopic limit that the number $`N`$ of particles making up the detectors approaches infinity. It was found that the factorization structure
$$F_{m,n}=D_m|D_n\underset{j=1}{\overset{N}{}}D_m^{[j]}|D_n^{[j]}$$
$`(1.4)`$
concerning the overlapping of detector-states plays a crucial part in quantum decoherence. Here, $`|D_n^{[j]}`$ are the single states of those blocks constituting the detector,and $`F_{m,n}`$ is called decoherence factor. Since each factor $`D_m^{[j]}|D_n^{[j]}`$ in $`F_{m,n}`$ has a norm less than unity, the product of infinite such factors may approach zero. This investigation was developed based on the Hepp-Coleman mode and its generalizations\[24-27\]. In 1998, this theory was applied to the analysis of the universality of the environment influences on quantum computing process \[29-31\].Parallelly, the classical limit that certain quantum numbers ( such as angular momentum) are huge is also investigated in our previous works .
However we have not got a totally-satisfactory answer to the question why the large system entangling with the small system behaves so classically in such limit situations. In fact,concerning the transition of the detector from quantum status to classical status , there were only some vague presentations\[17, 19, 21,\] . In a general situation the classical feature of the large system can not simply be characterized by large quantum numbers, and thus what is responsible for the classical feature remains unclear yet. Besides, all of our previous discussions about quantum decoherence are based on interaction of particular forms, namely the non-demolition interaction . In this paper, using Born-Oppenheimer (B-O) approximation , we universally consider the decoherence problem for a quantum system coupling to a large system through a general interaction. This basic approach can be applied to analyzing influences exerted by environment and detector as well. Our discussion is also involved with a fundamental problem that the physicist can not avoid completely : how does the time reversal symmetry implied by the Schrodinger equation on the microscopic scale turn into the time reversal asymmetry manifested by quantum decoherence or quantum dissipation on the macroscopic scale?
This paper is organized as follows. We describe in Sec.2 the adiabatic factorization of slow and fast dynamic variables in terms of the B-O approximation and show how the interaction of the large object with a quantum system causes a quantum entanglement dynamically. In Sec.3,incorporating the semi-classical approach to the quasi-classical motion of slow variable in a smooth potential, we manifest that, driven by the adiabatically-effective Hamiltonians, the final states of the large object initially in an appropriate state are orthogonal to one another,and their entanglement with the quantum system leads to decoherence. In Sec.4, Sec.5 and Sec.6, the universal formalism in Secs.2 and 3 is illustrated by three explicit examples : a. A particle with spin $`\frac{1}{2}`$ moves slowly in an inhomogeneous magnetic field of varying direction; b. A two level quantum system interacting with a very large spin; c. A quantized cavity field is coupled with a simple harmonic oscillator. The first illustration is similar to the Stern-Gerlach experiment . It stresses that the classical properties of the large system can be understood in terms of the macroscopic distinguishability of its quantum states. The second illustration reflects a simple presentation of quantum- classical transition when the quantum number is huge . The third illustration has certain practical significance as it is relevant to the problem of detecting gravitational wave by intracavity dynamics . It demonstrates the necessity of choosing a quasi-classical initial state of the large system to realize the quantum coherence of the quantum system. In Sec.7, based on the adiabatic approach for quantum decoherence, we discuss the spatial localization of the macroscopic object resulting from the adiabatic entanglement between its collective coordinate and the dynamic variables of particles constituting it. This study provides us with a possible solution to the Schroedinger cat paradox.Concluding remarks are given in the end.It includes a brief discussion about the development of quantum dynamic theory of decoherence. In connection with the coupled channel theory, one of whose concrete realization is the B-O approach, our discussion reveals the possibility of generalizing our present work to the case with a complicated and hence more practical interaction than the over-simplified interaction as shown in the two examples presented in Sec. 4 , 5 and 6.
## II Quantum Entanglement via Born-Oppenheimer Approach
In a very wide sense, any interaction between two quantum systems can cause an entanglement between them. In general, it then realizes a quantum measurement in a certain meaning. This is because one quantum system in different states can act on another with different effects correspondingly. However, this entanglement and its relevant quantum measurement is generally not very ideal because the usual interaction can not produce a one-one correspondence between the states of the two systems. Indeed, only a very particular interaction or its effective reduction can lead to an ideal entanglement and thereby an ideal quantum measurement. Nevertheless, fortunately , so long as one of the two systems can be separated adiabatically and behaves classically,as we will prove in the following, any interaction can result in an ideal entanglement in the evolution of the total system through its adiabatic reduction based on Born-Oppenheimer (B-O) approximation.
From the view point of BO approximation, we consider a total quantum system (“molecular”) with two sets of variables, a fast (“electric”) one $`q`$ and a slow (nuclear) one $`x`$. Resolving the dynamics of fast variables for a given motion of the slow subsystem, we obtain certain quantum states labeled by $`n`$ for the fast part. To the first order approximation, the left effective Hamiltonian governing the slow variables involves an external scalar potential $`V_n(x)`$ and an magnetic-like vector potential $`A_n(x)`$ induced by the fast variables . The latter is called the induced gauge potential or Berry’s connection. If we assume the motions of the slow subsystem are “classical”, we naturally observes that, due to the back-actions of the fast part , there are different induced forces
$$F_n=_xV_n(x)+\frac{d}{dt}x\times (_x\times A_n)$$
$`(2.1)`$
exerting on the slow part. Their direct physical effects are that the information of the “fast” states labeled by $`n`$ is recorded in the different motion configurations of the slow part. An entanglement just stems from this correlation between the quantum states of the fast subsystem and the classical motion configurations of the slow subsystem. In spirit of this physically-intuitive observation, we study the production of such quantum entanglement from the adiabatic separation of slow and fast variables based on the B-O approach.
Let us consider the interaction between a quantum system $`S`$ with fast dynamic variable $`q`$ and the large system $`E`$ with slow variable $`x.`$ The former with the Hamiltonian $`H_s=H_s(q)`$ can be regarded as a subsystem soaked in an environment or a measured system monitored by a detector, and the latter with the Hamiltonian $`H_E=H_E(x)`$ as the environment or the detector accordingly. In general the interaction Hamiltonian is written as $`H_I=H_I(x,q)`$ . For a fixed value of slow variable $`x`$ of $`E`$ , the dynamics of the quantum system is determined by the eigen-equation
$$[H_s(q)+H_I(x,q)]|n[x]=V_n(x)|n[x]$$
$`(2.2)`$
Both the eigen-values $`V_n[x]`$ and the eigen-state $`|n[x]`$ depend on the slow variable $`x`$ as a given parameter.
Usually, the variation of the Hamiltonian $`H_s(q)+H_I(x,q)`$ with $`x`$ can cuase transition from an energy level $`V_n(x)`$ of the quantum system to another level $`V_m(x)`$. But within the spatial domain $`R`$ to which the slow variable $`x`$ belongs, if the variable $`x`$ changes so slowly that the adiabatic conditions \[39-42\]
$$\left|\frac{n[x]|_x|m[x]\frac{d}{dt}x}{V_m(x)V_n(x)}\right|=\left|\frac{n[x]|\{_xH_I(x,q)\}|m[x]\frac{d}{dt}x}{\{V_m(x)V_n(x)\}^2}\right|1$$
$`(2.3)`$
hold for any two of the different energy levels {$`V_n(x)`$} , this transition can be physically neglected and then the BO approximation works as an effective approach. Let $`|\mathrm{\Phi }_{n,\alpha }`$ be the full eigen- function of the full Hamiltonian $`H=H_E(x)+H_s(q)+H_I(x,q)`$ for the total system formed by the large system plus the quantum system. The B-O approximation treats it as a partially factorized function
$$x|\mathrm{\Phi }_{n,\alpha }=\varphi _{n,\alpha }(x)|n[x]$$
$`(2.4)`$
of the slow and fast variables $`x`$ and $`q.`$ Here, the set of slow components $`\{\varphi _{n,\alpha }(x)=x|\varphi _{n,\alpha }\}`$ and the corresponding eigen-values $`\omega _{n,\alpha }`$ are obtained by solving the effective eigen-equation
$$H_n(x)\varphi _{n,\alpha }(x)=\omega _{n,\alpha }\varphi _{n,\alpha }(x)$$
$`(2.5)`$
The effective Hamiltonian $`H_n(x)`$ is defined by
$$H_n(x)=H_{nE}(x)+V_n(x)$$
$`(2.6)`$
where $`H_{nE}(x)`$ is an gauge-covariant modification of $`H_E(x)`$. It was obtained by replacing the momentum operator $`p=i\mathrm{}_x`$ with its gauge-covariant form $`p=i\mathrm{}_xA_n(x)`$ . Here, $`A_n(x)=in[x]|_xn[x]`$ is a $`U(1)`$ gauge potential induced by the motion of the quantum system. In the classical limit that the slow part behaves classically, an effective dynamics of interaction between quantum and classical objects naturally results from the effective Hamiltonians or its relevant Lagrangian.
The completeness relations $`_{n,\alpha }|\mathrm{\Phi }_{n,\alpha }\mathrm{\Phi }_{n,\alpha }|=1`$for the full eigen-functions $`|\mathrm{\Phi }_{n,\alpha }`$ can be expressed in $`xrepresentation`$ as
$$\underset{n,\alpha }{}𝑑x𝑑x^{}\varphi _{n,\alpha }(x^{})\varphi _{n,\alpha }(x)|xx^{}||n[x]n[x]|=1$$
$`(2.7)`$
which is equivalent to
$$\underset{n}{}|xx||n[x]n[x]|=|xx|,\underset{\alpha }{}|\varphi _{n,\alpha }\varphi _{n,\alpha }|=1$$
$`(2.8)`$
After obtaining the complete set {$`\varphi _{n,\alpha }(x)|n[x]`$} of eigenstates of the total system, we can now consider how the entanglement appears in the adiabatic dynamic evolution. Let the total system be initially in the state $`|\mathrm{\Psi }(t=0):`$
$$x|\mathrm{\Psi }(t=0)=\{\underset{n}{}c_n|n[x]\}\varphi (x)$$
$`(2.9)`$
The first component of the initial state $`|\mathrm{\Psi }(t=0)`$ is a superposition of the eigenstates of the quantum system while the second one a single pure state. Expanding $`|\mathrm{\Psi }(t=0)`$ in terms of the complete set {$`\varphi _{n,\alpha }(x)|n[x]`$}, we have the evolution wave function at time $`t`$
$$x|\mathrm{\Psi }(t)=\underset{n,\alpha }{}c_n\varphi _{n,\alpha }|\varphi \mathrm{exp}[i\omega _{n,\alpha }t]|n[x]\varphi _{n,\alpha }(x)$$
$`(2.10)`$
where we have used the completeness relation eq.(2.7). In terms of the effective Hamiltonian $`H_n(x)`$ related to each single adiabatic state $`|n[x],`$ the above wave function is rewritten in a concise form
$$x|\mathrm{\Psi }(t)=\underset{n}{}c_n|n[x]x|D_n(t)$$
$`(2.11)`$
with
$$|D_n(t)=\underset{\alpha }{}\varphi _{n,\alpha }|\varphi e^{i\omega _{n,\alpha }t}|\varphi _{n,\alpha }=\mathrm{exp}[iH_nt]|\varphi (x)$$
$`(2.12)`$
The full wave function$`|\mathrm{\Psi }(t)`$ is obviously an entangled state. Starting from the same initial state $`|\varphi `$ at $`t=0`$, the large system will be subject to different back-actions defined by $`(V_n,A_n)`$ from the different adiabatic states $`|n[x]`$ of the quantum system. Then it evolves to a superposition of different final states $`|D_n(t)`$. This intuitive argument shows us that, there indeed exists an entanglement between two quantum systems with an quite general interaction,if one of them moves so slowly that their dynamic variables can be adiabatically factorised according to the B-O approximation. Roughly speaking, in the B-O approach, the slow subsystem is usually referred to as heavy particles (such as nucleons ) while the fast one as light particles (such as the electrons). So it is reasonable to expect the slow subsystem to behave as a classical object.
## III Decoherence: Transition from Quantum to Classical
In this section we will discuss under what conditions the large system, the environment or the detector, can behave classically so that the quantum system entangled with it could completely lose its coherence and approach the classical limit.
Consider the reduced density matrix of the quantum system
$`\rho _s(t)`$ $`=`$ $`\text{Tr}_D(|\mathrm{\Psi }(t)\mathrm{\Psi }(t)|)={\displaystyle \underset{n}{}}|C_n|^2|n[x]n[x]|+`$ (3.2)
$`+{\displaystyle \underset{nm}{}}C_mC_n^{}|m[x]n[x]|D_n(t)|D_m(t)`$
obtained by ”summing over” the variables of the large system. The off-diagonal term responsible for interference is proportional to the overlapping $`F_{n,m}=D_n(t)|D_m(t)`$ of the two large system states.Were there no large system interacting with it, the quantum system would be completely coherent for $`\rho _s(t)=|\phi (t)\phi (t)|`$ is a pure state. Here $`|\phi (t)=\mathrm{exp}(iH_st)|\phi `$ is a free evolution state of the large system. Mathematically, the effect of the adiabatic effective interaction is to multiply the off-diagonal term of the reduced density matrix by the decoherence factor $`F_{n,m}.`$ A complete decoherence is defined by $`F_{n,m}=0`$ while a complete coherence by $`F_{n,m}=1(mn).`$
Before considering how the decoherence factor $`F_{n,m}`$ becomes zero for the large system, we need to review some known arguments about the meaning of the classical limit of the motion of the large system. According to a widely accepted viewpoint , in the classical limit, the expectation value of an observable for certain particular states should recover its classical value forms. These particular states can give definite classical trajectories of particle in this limit. Usually we call them quasi-classical states. A coherent state or its squeezed version is a typical example of such states. According to Landau and Lifshitz , in general, a quasi-classical state is a particular superposition $`_nc_n\varphi _n`$ with the non-zero coefficients $`c_n`$ only distributing around a large quantum number $`\stackrel{~}{n}`$ . Then the correspondence principle requires that $`\stackrel{~}{n}\mathrm{},\mathrm{}0`$ and the product $`\stackrel{~}{n}\mathrm{}`$ approaches a finite classical action. In such a limit, the expectation of an observable will take the Fourier series of its corresponding classical quantity ; or strictly speaking, it takes the Fej$`\stackrel{^{}}{e}`$r’s arithmetic mean of the partial sums of the Fourier series . In this sense the mean-square deviation of the observable is zero; and accordingly the mean of the position operator defines a classical path. Physically, the zero mean-square deviation of the position operator implies the zero width of each wave packet $`x|D_m(t),`$and the overlapping $`F_{n,m}=D_n(t)|D_m(t)`$ of zero width wave packets must vanish. From such a semi-classical picture, we will clearly see in the following how the decoherence factor $`F_{n,m}`$ approaches zero dynamically as the large system becomes classical.
In the semi-classical approach, for a heavy particle, the initial state $`|\phi `$ can be regarded as a very narrow wave packet of width $`a`$. Since the heavy particle has a large mass $`M`$ it hardly spreads in the evolution because without the environment induced quantum dissipation the width of the wave packet at time $`t`$ is
$$w(t)=a\sqrt{1+\frac{t^2}{4M^2a^4}}$$
$`(3.2)`$
Then we describe the large system as an moving wave packet with the center along a classical path $`x(t)`$ on a manifold with local coordinates $`x.`$ For a proper initial state $`|\phi ,`$we will see that the wave packet will split into several narrow peaks with the centers along different paths determined by different motion equations governed by the effective forces $`F_n=_xV_n(x)+\frac{d}{dt}x\times (_x\times A_n)`$ with effective potentials $`(V_n(x),A_n).`$ Usually,the widths of these peaks are almost of the same order as that of the original wave packet and each peak is correlated to an adiabatic quantum state $`|n[x]`$ for a large mass. Except for some moments at which the centers of two or more peaks coincide, these narrow peaks hardly overlaps with one another. In this sense, the large system starting from a narrow initial state can reach a superposition of those states orthogonal to one another. Thus we approximately have $`F_{n,m}=0`$ in the classical limit for $`mn.`$
With reference to the useful analysis in ref., we present an explicit but sketchy calculation to justify the above physically-intuitive observation about $`F_{n,m}=0`$ in the classical limit. Assume the large system to be a heavy particle with very large mass $`M`$. In the duration $`\tau `$ of the adiabatic interaction with the quantum system, if the condition $`v\tau (\mathrm{\Delta }p/M)`$ $`\tau \mathrm{\Delta }x`$ holds, the momentum $`p_0`$ of the free heavy particle can not be changed notably. Thus the contributions of the kinetic term and the induced gauge potential can be ignored in the wave function evolution of the free heavy particle under this condition. From this consideration we can approximately write down
$$|D_n(t)=e^{iH_nt}|\varphi (x)e^{iV_n(x)t}|\varphi (x)$$
$`(3.3)`$
The approximation requires that the effective potential $`V_n(x)`$ is satisfactorily smooth or the interaction $`H_I(x,q)`$ is a smooth function of $`x.`$ So we can use
$$V_n(x)V_n(0)+F_nx;F_nV_n(0)$$
$`(3.4)`$
to re-express the decoherence factor
$$F_{n,m}=\phi |\mathrm{exp}\left(it\delta F(m,n)x\right)|\phi $$
$`(3.5)`$
Here, $`\delta F(m,n)=F_mF_n`$ is the difference of two external forces exerted by two adiabatic potentials $`V_m(x)`$ and $`V_n(x).`$ Then the role of the back-action of the quantum system on the large system is summing up the momentum shift by a quantity $`\delta F(m,n)t`$ with respect to the initial state $`|\phi `$. Obviously, when the width $`\sigma =a^1`$ of the initial wave packet $`p|\phi `$ in the momentum space is much less than the momentum shift $`\delta F(m,n)t`$, the large system will adiabatically evolves into states orthogonal to one another. In fact, if the initial state is chosen to be a Guassian wave packet $`x|\phi =\frac{\sigma }{\sqrt{\pi }}\mathrm{exp}[\frac{1}{2}\sigma ^2x^2]`$ of width $`x=\frac{1}{\sigma },`$ the decoherence factor is a Guassian decaying function of time $`t`$
$$F_{n,m}=\mathrm{exp}\left(\frac{\delta F(m,n)^2}{4\sigma ^2}t^2\right)$$
$`(3.6)`$
As the evolution time $`t`$ approaches infinity or if we have a very narrow width $`\sigma `$ , $`F_{n,m}0`$ and a quantum decoherence results from the dynamical evolution automatically.
Generally, we consider a system described by $`H_n=p^2/2M+V_n(x)`$ without the induced gauge field. Define $`x_c=\phi |x|\phi `$ and $`p_c=\phi |p|\phi `$ for an initial state $`|\phi `$. In the classical regime one may expect that the variations $`\xi `$ $`=xx_c`$ and $`p_\xi =pp_c`$ are small compared with $`x_c`$ and $`p_c.`$Accordingly the potential can be expanded as
$$V_n(x)V_n(x_c)+V_n^{}(x_c)\xi +\frac{1}{2}V_n^{\prime \prime }(x_c)\xi ^2.$$
$`(3.7)`$
So, approximately the Heisenberg equations of motion become
$$\frac{d}{dt}x=\frac{p}{M},\frac{d}{dt}p=V_n^{}(x_c)V_n^{\prime \prime }(x_c)\xi $$
$`(3.8)`$
Sandwiched by the initial state $`|\phi `$, the above equations turn into the classical equations of motion
$$\frac{d}{dt}x_c=\frac{p_c}{M},\frac{d}{dt}p_c=V^{}(x_c).$$
$`(3.9)`$
Now we turn to Schrödinger’s picture. The evolution of the initial state is governed by $`i\mathrm{}_t|\phi (t)=H_n|\phi (t)`$. Introduce the following time-dependent translation
$$|\varphi (t)=\mathrm{exp}\left\{\frac{i}{\mathrm{}}\left(\theta (t)+x_cp_\xi p_c\xi \right)\right\}|\phi (t)S(t)|\phi (t)$$
$`(3.10)`$
where $`\theta (t)`$ is determined by $`\dot{\theta }_t=p_c^2/2M+V(x_c)`$. Then straightforward calculation gives
$$i\mathrm{}_t|\varphi (t)=\left(\frac{p_\xi ^2}{2M}+\frac{1}{2}M\omega _t^2\xi ^2\right)|\varphi (t)$$
$`(3.11)`$
where $`M\omega _t^2=V^{\prime \prime }(x_c)`$ . This exactly describes an oscillator with time-dependent frequency. The above direct derivation shows that in the non-inertial frame moving along the classical orbit, every quasi-classical system looks like a time-dependent oscillator whose frequency depends on the orbit. This fact is an established conclusion and illustrated in Fig 1. Actually, it is present in many textbooks about path integral. But our argument here is based on a clear physical picture and is applicable to the three dimensional case after a slight generalization.
Denote by $`|0`$ the vacuum state of the harmonic oscillator with frequency $`\omega _0`$ which is equivalent to a Gaussian wave packet of width $`\sigma _0^1=\sqrt{2m\omega _0/\mathrm{}}`$.Suppose that initially the system is in the state $`S^{}(0)|0`$, a coherent state whose center lies at $`(x_c(0),p_c(0))`$. At time $`t`$,the center of the wave packet is obviously at $`(x_c(t),p_c(t))`$, and it is reasonable to expect that the width of the wave packet becomes $`\sigma _t^1=\sqrt{2m\omega _t/\mathrm{}}`$,since the frequency of the time-dependent oscillator changes very slowly. For two different potentials $`V_1(x)`$ and $`V_2(x)`$ the macroscopic distinguishability is ensured when the width sum of the two evolved packets is less than their orbital difference, that is, when
$$\sigma _{1t}+\sigma _{2t}|x_{c1}(t)x_{c2}(t)|$$
One cannot expect that this condition can always be fulfilled for all time $`t`$. The orbital difference is determined by something like $`|V_1^{}V_2^{}|`$ and the width is determined by the second derivative of the potential. But their relation is not very clear to us at present. What is clear is,to have $`D_1(t)|D_2(t)=0`$ one should require the points that fail the inequality form a zero measure set. On the other hand, the adiabatic approximation also imposes some restrictions on the potential. To clarify the situation further more sophisticated considerations are needed.
## IV From Macroscopic Distinguishability to Decoherence
In the context of quantum measurement,a variant of the Stern-Gerlach (SG) experiment provides an illustration of the above formalism. Quantum measurement is mutationally an observing process that “reads out” the system states from the “macroscopically distinguishable” states of the detector. As is shown in the above, if the large particle moves slowly enough, an adiabatic eigen-state of the quantum system will be correlated to one of the detector states in the B-O approximation. So the adiabatic correlation
$$|1[x]|D_1(t),|2[x]|D_2(t),\mathrm{}\mathrm{}|n[x]|D_n(t)$$
$`(4.1)`$
between the system states $`|n[x]`$ and the detector states $`|D_n`$ defines a quantum measurement. In the classical limit, this measurement is thought to be ideal for $`|D_n(n=1,2,\mathrm{})`$ are orthogonal to one another, i.e., $`|D_n`$ are shown to be “classically- or macroscopically distinguishable”. Once the detector is found in the state $`|D_n`$, we can infer that the system is just in the state $`|n`$. In the following we will quantitatively analyze the dynamical realization of such an adiabatic measurement in a variant of the SG experiment.
The original Stern-Gerlach (SG) experiment can be considered as a quantum measurement process detecting the spin states of particles from their spatial distribution. The WPC or quantum decoherence can be described in an dynamical evolution governed by the interaction between the space- and spin- degrees of freedom. In its variant, a spin-$`\frac{1}{2}`$ particle initially in a certain superposition state enters an inhomogeneous magnetic field of amplitude $`B(x)`$ with varying direction $`𝐧(x)=(\mathrm{sin}\theta \mathrm{cos}kx,\mathrm{sin}\theta \mathrm{sin}kx,\mathrm{cos}\theta )`$ where $`\theta `$ is fixed. Its configuration is shown in Fig.2. A simple experiment though it is,it is among the candidates of experiments proposed to test the Berry’s phase or its corresponding induced gauge field for a neutron in a static heliacal magnetic field \[46, 42 \]. In the usual S-G experiment, the direction of the magnetic field is along the fixed $`x`$axis, but in our present model the polarization direction varies as the position $`x`$ changes.
The spatial variable is considered to be the slow system while the spin -variable to be fast as a quantum system. Corresponding to the eigenvalues $`V_\pm (x)=\pm \mu B(x)`$, the adiabatic eigenstates of the spin-Hamiltonian $`H_{spin}=\mu B\left(x\right)𝐧(x)\sigma `$ are
$$|\chi _+[x]=\left[\begin{array}{c}\mathrm{cos}\frac{\theta }{2}e^{ikx}\\ \mathrm{sin}\frac{\theta }{2}\end{array}\right],|\chi _{}[x]=\left[\begin{array}{c}\mathrm{sin}\frac{\theta }{2}e^{ikx}\hfill \\ \mathrm{cos}\frac{\theta }{2}\hfill \end{array}\right]$$
Here $`\sigma =(\sigma _x,\sigma _y,\sigma _z)`$ is the Pauli spin operator and $`\mu `$ the gyromagnetic ratio. Let the incoming beam be initially in a superposition of the adiabatic eigen-states $`|\psi =c_+|\chi _+[x]+c_{}|\chi _{}[x]`$ along a certain polarization direction depending on $`x`$. When the particle moves so slowly that the adiabatic condition
$$|\frac{d}{dt}xk\mathrm{sin}\theta /\mu B(x)|1$$
$`(4.2)`$
holds, to the lowest order of the B-O approximation , the total initial state $`|\mathrm{\Psi }(0)=\{c_+|\chi _+[x]+c_{}|\chi _{}[x]\}|\varphi (x)`$ will evolves into an entangled state
$$|\mathrm{\Psi }(t)=c_+|\chi _+[x]|D_+(t)+c_{}|\chi _{}[x]|D_{}(t)$$
$`(4.3)`$
Here, $`|D_\pm (t)=\mathrm{exp}[iH_\pm t]|\varphi (x)`$ are the spatial states governed by the effective Hamiltonians
$$H\pm =\frac{1}{2M}(i_xA_\pm )^2+V_\pm (x)$$
$`(4.4)`$
The effective scalar potentials $`V_\pm (x)`$ and the induced vector potentials $`A_\pm =\frac{1}{2}k(1\pm \mathrm{cos}\theta )`$ are determined from the adiabatic spin eigenstates $`|\chi _+[x]`$and $`|\chi _{}[x]`$. In the semi-classical picture, because the particles in the adiabatic spin states $`|\chi _+[x]`$and $`|\chi _{}[x]`$ separately suffer two forces $`F_\pm =\frac{}{x}V_\pm (x)`$ of opposite directions along $`𝐱`$, they will finally form two macroscopically -distinguishable spots on the detecting screen, each of which is correlated to one of the spin states. This spin-space correlation process enables people to pick out different spin states according to the spatial distribution.
To analyze this measurement process in details we assume the spatial part $`\varphi (x)`$ in the initial state is a Gaussian wave packet
$$|\varphi (x)=\left(\frac{1}{2\pi a^2}\right)^{\frac{1}{4}}𝑑xe^{\frac{x^2}{4a^2}}|x$$
$`(4.5)`$
distributing along direction $`x`$ with the center at the original point. Here $`a`$ is the initial width of the atom beam. Adopting the semi-classical method, we have the linear approximation $`B(x)[_xB\left(x=0\right)]x`$ and $`f=\mu _xB\left(x=0\right)`$. Factorizing the evolution operator $`U_\pm (t)=\mathrm{exp}[iH_\pm t]`$ by Wei-Norman method (see Appendix 1), we exactly obtains,in position representation, the following effective wave functions $`|D_\pm (t)`$at time $`t`$
$$x|D_\pm (t)=\left(\frac{a^2}{2\pi ^3}\right)^{\frac{1}{4}}\left(\frac{\pi }{a^2+\frac{it}{2M}}\right)^{\frac{1}{2}}e^{i\mathrm{\Omega }_\pm \left(t\right)iftx}\mathrm{exp}\left[\frac{\left(xx_{\pm c}(t)\right)^2}{4(a^2+\frac{it}{2M})}\right]$$
$`(4.6)`$
where
$$\mathrm{\Omega }_\pm (t)=\frac{f^2t^3}{6M}+\frac{1}{2}ft^2A_\pm $$
It is seen from Eq.(4.6) that the Gaussian wave packets $`x|D_\pm (t)`$ center on the classical trajectories
$$x_{\pm c}(t)=\frac{1}{2}\frac{f}{M}t^2\frac{A_\pm }{M}t$$
$`(4.7)`$
They have the different group speeds $`v_\pm =\frac{f}{M}t`$ $`\frac{A_\pm }{M}`$ along the opposite directions, but have the same width $`a(t)=a\sqrt{1+t^2/(4M^2a^2)}`$ spreading with time. It is obvious that the motions of the wave packet centers obey the classical motion law that a particle of mass $`M`$ forced by $`f`$ will move with the acceleration $`f/M`$. The quantum character of this motion is mainly reflected in the spreading of the wave-packets. The induced gauge fields $`A_\pm `$ are constant, but they change the initial value of $`\frac{d}{dt}x`$ according to the corresponding classical Hamilton equation.
$$M\frac{d^2}{dt^2}x=f;\frac{d}{dt}x=\frac{p}{M}\frac{A_\pm }{M}$$
$`(4.8)`$
This means that the zero initial value of the canonical momentum $`p=M\frac{d}{dt}x+A_\pm `$ determines the initial velocity $`\frac{d}{dt}x(0)=A_\pm /M.`$ The quantum effects of $`A_\pm `$ are to contribute the additional phases $`\frac{1}{2}ft^2A_\pm `$in the wave functions.
The macroscopic distinguishibility of wave-packets in quantum measurement requires that the distance between the two wave-packets should be larger than the width of each wave packet, i.e.
$$ft^2k\mathrm{cos}\theta ta\sqrt{M^2+\frac{t^2}{4a^2}}$$
$`(4.9)`$
This condition is easily satisfied for a long time evolution.
To analyze the decoherence quantitatively, we compute the norm of the decoherence factor $`F(t)=|D_+(t)|D_{}(t)|`$. The extent of quantum coherence depends totally on this overlapping integral. We can explicitly integrate it
$$F(t)=\mathrm{exp}\left[a^2f^2t^2\frac{1}{8a^2}\left(\frac{f}{M}t^2\frac{k\mathrm{cos}\theta }{M}t\right)^2\right]$$
$`(4.10)`$
It is obvious that the decoherence process indeed happens as $`t\mathrm{},`$ but it does not obey the simple exponential law $`e^{\gamma t}.`$ In a long time scale, the temporal behavior of decoherence is described by $`F(t)\mathrm{exp}\left[\frac{f^2t^4}{8a^2M^2}\right]`$ and the characteristic time of the decoherence process can be defined by $`F(\tau _d)=e^1`$ , that is
$$\tau _d=\sqrt{\frac{2\sqrt{2}Ma}{f}}$$
$`(4.11)`$
This shows that the long time behavior of decoherence is indepedendent of the spatial details of interaction , which is caused by the configuration of the external field.
## V Decoherence Resulting from Large Spin
There is a second illustration to show the happening of decoherence owning to the adiabatic separation of two systems. Based on our previous investigation about quantum decoherence in the classical limit , we assign an arbitrary spin $`j`$ to interact with a two-level system (such as a spin-$`\frac{1}{2}`$ system) through a coupling of particular form. Let $`𝐉=(\widehat{J}_x,\widehat{J}_y,\widehat{J}_z)`$ be the angular momentum operator of the large system and $`\sigma =(\sigma _x,\sigma _{y,}\sigma _z)`$ be the Pauli matrix describing the quasi-spin of the two-level quantum system with energy-level difference $`\omega _s`$. The full Hamiltonian of this model is
$$H_I=\omega _s\sigma _z+\omega J{}_{𝐳}{}^{}+𝐟(𝐉)\sigma _x,$$
$`(5.1)`$
The general interaction $`𝐟(𝐉)\sigma _x`$ is linear with respect to the variable of the quantum system while it depends on the variable $`𝐉`$ through a function $`𝐟(𝐉)`$. Two free Hamiltonians $`\omega _s\sigma _z`$ and $`\omega J_𝐳`$ were introduced to consider the energy-exchange between the quantum system and the large system.
The interaction $`𝐟(𝐉)\sigma _x`$ can not well distinguish the states $`|\pm \frac{1}{2}`$ of the quantum system for $`|\pm \frac{1}{2}`$ are not the eigen- states of the interaction Hamiltonian. So, in general, this model can not well describe a quantum measurement process and thus can not give a good description of quantum decoherence. However, if we think $`𝐉`$ as the slowly-changing variable relative to the fast one $`\sigma ,`$determined by the B-O approximation under the adiabatic condition, the effective potential $`V_\pm =\pm \sqrt{\omega _s^2+𝐟(𝐉)^2}`$ of the large system will clearly distinguish the adiabatic eigen-states $`|u_+[𝐉]=(\mathrm{cos}\frac{\vartheta }{2},\mathrm{sin}\frac{\vartheta }{2})^T`$ and $`|u_{}[𝐉]=(\mathrm{sin}\frac{\vartheta }{2},\mathrm{cos}\frac{\vartheta }{2})^T`$. Here, the angle parameter $`\vartheta =\mathrm{arg}\mathrm{tan}(\frac{𝐟(𝐉)}{\omega _s})`$ depends on the slow variable $`𝐉.`$ Then, the adiabatic separation of the spin-$`\frac{1}{2}`$ and the large spin system will result in a quantum decoherence.
In fact, because of the introduction of the arbitrary spin $`j`$ , which labels the $`2j+1`$-dimensional irreducible representation of the rotation group SO(3), we are able to consider the behaviors of the quantum dynamics governed by this model Hamiltonian in the classical limit with infinite spin $`j`$. The reason why the limit with infinite $`j`$ is called classical is that the mean square deviations of the components $`\widehat{J}_x,`$ and $`\widehat{J}_y`$ enjoy the following limit feature $`\frac{\mathrm{\Delta }\widehat{J}_x}{j}=\frac{\mathrm{\Delta }\widehat{J}_y}{j}=\frac{1}{\sqrt{2j}}0`$ as$`j\mathrm{}`$.
To solve the dynamical evolution of the total system explicitly, we choose a particular form of interaction : $`𝐟(𝐉)=\sqrt{g^2J_x^2\omega _s^2}.`$ Taking this particular form is equivalent to making a linear approximation for the effective potential $`V_\pm [𝐉].`$ With this particular form the effective Hamiltonians $`H_\pm =\omega J{}_{𝐳}{}^{}+`$ $`V_\pm [𝐉]`$ can be expressed as an rotation of the simple spin-Hamiltonian $`H_o=\sqrt{g^2+\omega ^2}J{}_{𝐳}{}^{},𝐢.𝐞.,`$
$$H_\pm =\mathrm{exp}[i\widehat{J}_y\varphi _\pm ]H_o\mathrm{exp}[i\widehat{J}_y\varphi _\pm ]$$
$`(5.2)`$
where the polar angle $`\varphi _\pm `$ is defined by $`\mathrm{tan}\varphi _\pm =\pm \frac{g}{\omega }`$.
According to the quantum angular momentum theory, the eigen-states of $`H_\pm `$ can be constructed as
$$|j,m(\varphi _\pm )=\mathrm{exp}[i\widehat{J}_y\varphi _\pm ]|j,m=\underset{m=j}{\overset{j}{}}d_{m^{},m,}^j(\varphi _\pm )|j,m^{}$$
$`(5.3)`$
where $`|j,m`$ is a standard angular momentum state and $`d_{m^{},m,}^j(\varphi )=j,m^{}|\mathrm{exp}[i\widehat{J}_y\varphi ]|j,m`$ is the corresponding $`d`$ function ; the corresponding eigen-values are $`E_m=m\sqrt{g^2+\omega ^2}`$.
Here, we should remark that the exact solvability of the above model largely depends on the particular form of the function $`𝐟(𝐉)`$. If this is not the case, the above method can not work well and then certain semi-classical approximation methods should be used to deal with the effective Hamiltonian in its classical limit with very large $`j`$. If the coupling function $`𝐟(𝐉)`$depends on $`𝐉`$ quite slightly, we can generally linearize the above effective potential $`V_\pm (𝐉)`$ to realize the particular form.
We are concerned with classical characters of the large-spin system. Let us suppose it is initially assigned the adiabatic ground state $`|j,m=j(\varphi )`$ with the lowest magnetic quantum number $`m=j`$. In quantum measurement theory, the choice of ground state is required by a stable measurement. Starting with its initial state
$$\psi (0)=(C_+|u_+[𝐉]+C_{}|u_{}[𝐉])|j,j(\varphi )$$
$`(5.4)`$
the effective Hamiltonians (5.2) evolves the large spin system into an entanglement state
$$\psi (t)=C_+|u_+[𝐉]D_+(t)+C_{}|u_{}[𝐉]D_{}(t),$$
$`(5.5)`$
with
$$D_\pm (t)=\mathrm{exp}[\pm i\widehat{J}_y\varphi ]\mathrm{exp}[it\stackrel{}{𝐉}{}_{𝐳}{}^{}\sqrt{g^2+\omega ^2}]\mathrm{exp}[i\widehat{J}_y\varphi ]|j,j(\varphi )$$
$`(5.6)`$
Using the explicit expressions of the $`d`$-function $`d_{m^{},m,}^j(\varphi _\pm ),`$ we can calculate the overlapping $`D_{}(t)D_+(t)`$ , obtaining
$$F(j;t)=|D_{}(t)D_+(t)|=\left|1\mathrm{sin}^22\varphi \mathrm{sin}^2\frac{\sqrt{g^2+\omega ^2}}{2}t\right|^j.$$
$`(5.7)`$
The above formula directly manifests the happening of quantum decoherence in the classical limit $`j\mathrm{}`$ . In fact, in a nontrivial case with $`\varphi 0`$ , $`|1\mathrm{sin}^2\frac{t}{2}\sqrt{g^2+\omega ^2}\mathrm{sin}^22\varphi |`$ is usually a positive number less than 1. In the classical limit with $`j\mathrm{}`$ , its $`jth`$ power $`|D_{}(t)D_+(t)|`$ must approach for $`tt_n2nt/\sqrt{g^2+\omega ^2},n=0,1,2\mathrm{}`$. At those instances $`t_n,`$quantum coherence revivals as so-called quantum jumps (see Fig.3). Then, as far as the present model is concerned, we reach the conclusion that, if the large spin system behaves classically, the decoherence can be dynamically realized for the entangled quantum system. In traditional quantum measurement, the detector was pre-required as a purely classical object to reduce the coherent superposition instantaneously. But now it is proved that the WPC occurs as the quantum detector moves slowly to approach the classical limit. This means in our treatment the detector is essentially still a quantum object. Thus it has the advantage of dealing with the problem of quantum measurement consistently within the framework of quantum theory.
## VI Intracavity Dynamics with Classical Source
Our third example about decoherence in quantum adiabatic process is the intracavity dynamics with a classical source, which is associated with the interferometric detection of the gravitational wave by a squeezed light .
We consider a cavity with two end mirrors (as in Fig.4), one of which is fixed while the other is treated as a simple harmonic oscillator of frequency $`\mathrm{\Omega }`$ and mass $`M`$ with the position and momentum $`x`$ and $`p.`$ The radiation pressure force of the cavity field on the moving mirrors is proportional to the intracavity photon density. Let $`a^{}`$ and $`a`$ be the creation and annihilation operators of the cavity with a single mode of frequency $`\omega `$. The cavity-mirror coupling is described by an interaction Hamiltonian $`H_I=gx`$ $`a^{}`$ $`a`$ where $`g`$ is the coupling constant depending on the electric dipole. In the radio frequency range the cavity field can be prescribed as a macroscopic current. From this consideration we describe the cavity field dynamics with the Hamiltonian $`H_c=\omega a^{}`$ $`a+f(a^{}`$ $`+a)`$. This cavity field -mirror coupling system can also be used to detect the photon number in the cavity by the motion of mirror. Obviously, the motion of the mirror is slow with respect to the oscillation of the cavity field. Thus we can use the B-O approximation to approach the quantum decoherence problem in the measurement of the cavity field. Most recently, the special case of this model without classical source has been used as a scheme probing the decoherence of a macroscopic object .
Coupled with the mirror and the classical source, the adiabatic eigen-states
$$|n[x]=\frac{1}{\sqrt{n!}}[a^{}+\lambda (x)]^n|0$$
$`(6.1)`$
of the cavity field for displacement $`\lambda (x)=\frac{f}{\omega +gx}`$ are determined by
$$\{[\omega +gx]a^{}a+f(a^{}+a)\}|n[x]=v_n(x)|n[x]$$
$`(6.2)`$
with the corresponding eigen-values $`v_n(x)=n(\omega +gx),n=0,1,2,\mathrm{}`$ Under the B-O approximation, the effective Hamiltonians are also referred to as the forced harmonic oscillators in the same renormalization external potential (RNEP)$`V_{rne}=\frac{f^2}{\omega +gx}`$ . Under the adiabatic condition
$$\left|\frac{(n1)[x]_x|n[x]\frac{d}{dt}x}{\omega +gx}\right|\frac{|ngf\frac{d}{dt}x|}{\omega ^3}1$$
$`(6.3)`$
$`\mu ,`$the RNEP $`V_{rne}`$ can be linearized as $`\frac{f^2}{\omega }[1\frac{gx}{\omega }|`$ . Then the effective Hamiltonians can be rewritten as $`H_n=\mathrm{\Omega }b^{}b+g_n(b^{}+b)`$ in terms of
$$b=\frac{M\mathrm{\Omega }x+ip}{\sqrt{2M\mathrm{\Omega }}},g_n=\frac{g(nf^2/\omega ^2)}{\sqrt{2M\mathrm{\Omega }}}=\mu \left(n\frac{f^2}{\omega ^2}\right)$$
$`(6.4)`$
For each effective Hamiltonian $`H_n,`$the corresponding evolution is a displacement operator
$$D[\alpha _n(t)]=\mathrm{exp}\left(\alpha _n(t)b^{}\alpha _n(t)^{}b\right)$$
$`(6.5)`$
with $`\alpha _n(t)=g_n(\mathrm{exp}[i\mathrm{\Omega }t]1)/\mathrm{\Omega }.`$
Let the initial state of the mirror be a well-defined quasi-classical state, a coherent state $`|\alpha `$ and the initial state of the cavity be a superposition $`|c(0)=_nc_n|n[x]`$of the adiabatic states. The evolution governed by the effective Hamiltonian $`H_n`$ leads to an entangled state
$$|\psi _I(t)=\underset{n}{}c_n|n[x]D[\alpha _n(t)]|\alpha \underset{n}{}c_n|n[x]|D_n(t)$$
$`(6.6)`$
for the total system. The overlapping of the mirror states in this entanglement can be computed and its norm is
$$|D_m(t)|D_n(t)|=\mathrm{exp}\left((nm)^2\frac{2\mu ^2}{\mathrm{\Omega }^2}\mathrm{sin}^2\frac{\mathrm{\Omega }t}{2}\right)$$
$`(6.7)`$
The changing rate $`\frac{d}{dt}x(`$ the velocity ) of the slow variable $`x`$ is proportional to $`\mathrm{\Omega }.`$ In the adiabatic limit, $`\mathrm{\Omega }`$ is very small. So we can rationally consider the limit $`\mathrm{\Omega }0`$ for a fixed $`\mu `$. Then an ideal entanglement appears in this limit case for the overlapping becomes an non-linear exponential decaying factor
$$|D_m(t)|D_n(t)|=\mathrm{exp}\left(\frac{1}{2}(nm)^2\mu ^2t^2\right)$$
$`(6.8)`$
This result is quite similar to that of the Cini model in van Hove limit . This decay phenomenon was first illustrated in ref.. Mathematically,it results from the fact that in the strong coupling limit, the period of the oscillation is very large in comparison with the small frequency $`\mathrm{\Omega }`$.
Another interesting situation arise when the mirror is initially prepared in a Fock state $`|n=\frac{1}{\sqrt{n}!}\left(a^{}\right)^n|0`$. To show a macroscopic, but non-classical dynamic behavior, the Fock state should possess a very large occupation number $`n`$. The overlapping for the initial Fock state can be expressed as
$$F(t,n)=n|D[\alpha _k(t)]D[\alpha _l(t)]|n$$
$$=\mathrm{exp}[\frac{1}{2}(lk)^2\frac{\mu ^2}{\mathrm{\Omega }^2}\mathrm{sin}^2\frac{\mathrm{\Omega }t}{2}]L_n\left((lk)^2\frac{\mu ^2}{\mathrm{\Omega }^2}\mathrm{sin}^2\frac{\mathrm{\Omega }t}{2}\right)$$
$`(6.9)`$
in terms of the Laguerre polynomial $`L_n(z)`$. Fig.5 shows $`F(t,n)`$ as a function of time $`t`$ for different $`j`$. In fact, according to the theory of special function, $`L_n(z)`$ approaches the zero-order Bessel function $`J_0(\sqrt{n}z)`$ when $`n\mathrm{}`$, hence ,
$$F(t,n)e^{\frac{1}{2}(lk)^2\mu ^2t^2}L_n((lk)^2\mu ^2t^2/4)e^{\frac{1}{2}(lk)^2\mu ^2t^2}J_0(\sqrt{n}(lk)^2\mu ^2t^2)$$
$`(6.10)`$
The zero-order Bessel function of real variable $`\zeta \sqrt{n}`$ is a decaying-oscillating function and approaches zero as $`n`$ tends to infinity. Therefore, when the cavity is occupied by a large number of photons, the macroscopic feature of the detector (the end mirror ) dynamically decoheres the initial pure state of the cavity.
## VII Localization of Macroscopic Object Through Adiabatic Entanglement
As another interesting application of the above adiabatic approach for decohernce, we will discuss how the adiabatic entanglement result in the spatial localization of a macroscopic object. This discussion is devoted to consider the quantum decoherence of the slow part rather than that of the fast part, which has been studied in previous sections.
The localization problem originated from the correspondence between Einstein and Born and is closely related to the Schrodinger cat. They observed that, usually in a spatially-localized state, a macroscopic object can only be described by a time-dependent localized wave packet, which is a coherent superposition of the eigen-states of the center-of-mass Hamiltonian $`H_0=p^2/2M`$ . Though the spreading of an initially well localized wave packet can be reasonably ignored for the macroscopic object with very large mass, Einstein argued that the superposition of two narrow wave packets is no longer narrow with respect to the macro-coordinate, but it is still a possible state of the macroscopic object. So a contradiction to the superposition principle arises because of the requirement that the wave function of a macroscopic object must be narrow. To solve this problem, Wigner , Joos and Zeh present the so called scattering -induced -decoherence mechanism (or WJZ mechanism): scattering of photons or atoms off a macroscopic object records the information of its position to form a quantum measurement about the position. In this mechanism the interference terms between different positions are destroyed by the generalized ”which-way”detection . In spirit of Omnes ’s observation , we argue that, mentioning macroscopicness implies the requirement that the macroscopic object must contain a large number of internal blocks. Then the macroscopic object,coupled to the internal variables, should be described by collective variables subject to the interaction similar to that concerning the external scattering in WJZ mechanism. In this section we will show that the spatial localization of a macroscopic object can be caused by an ideal entanglement between its collective position (or center-of-mass ) and internal variables . This entanglement results from their adiabatic separation..
Let $`x`$ and $`q`$ be, respectively, the collective position and internal variables of a macroscopic object with the collective Hamiltonian $`H_s=p^2/2M`$ ($`[x,p]=i`$). To consider how different positions affect the quantum coherence of the internal motion of the macroscopic object, we suppose that the total system is initially in a product state
$$|\mathrm{\Psi }_x(t=0)=|x|\varphi $$
$`(7.1)`$
where the first component $`|x`$ is the eigen-state of the collective position operator $`x`$ while $`|\varphi `$ is an arbitrary initial pure state of the internal degrees of freedom. Usually, the collective motion acts on the internal motion in certain ways and the back-action of the internal motion can not be neglected physically. So this generic interaction can not produce an ideal entanglement between the collective position and the internal states of the macroscopic object. By an argument similar to that by Joos and Zeh , who deal with quantum decoherence and its relevant localization by considering the scattering of external particles by the macroscopic object, we see only when the back-action is negligiblly small, can the interaction between the collective and internal states realize a ”measurement-like process”:
$$|x|\varphi U(t)|x|\varphi =|x(t)S(x;t)|\varphi $$
$`(7.2)`$
Here, $`U(t)`$ is the total evolution matrix and $`|x(t)`$ $`=U_0(t)|x`$ represents the free evolution in the absence of the coupling to the internal variables; $`S(x,t)`$, acting on the internal states, denotes the effective $`Smatrix`$ parametrized by the collective position $`x`$. If the collective motion is initially described by a wave packet $`|\phi =\phi (x)|x𝑑x,`$then the reduced density matrix of the collective motion is
$$\rho (x,x^{},t)=\phi (x,t)\phi ^{}(x^{},t)\varphi |S^{}(x^{};t)S(x;t)|\varphi $$
$`(7.3)`$
Considering the translational invariance of the scattering process, Joos and Zeh showed that, in $`xrepresentation`$, the off-diagonal terms take the following form
$$\varphi |S^{}(x^{};t)S(x;t)|\varphi \mathrm{exp}(\mathrm{\Lambda }t|xx^{}|^2)$$
$`(7.4)`$
This means the decoherence factor is a damping function with the localization rate $`\mathrm{\Lambda }`$ , which depends on the total cross section.
Now, the question arises whether the negligibility of the back-action is the unique cause for the appearance of the above mentioned ”measurement-like process”. If not, what are the other causes beyond it? To resolve this problem, we assume the Hamiltonian $`h(q,x)=H_i(q)+W(x,q)`$ describes the motion of the internal variables $`q`$ coupling to the collective variable $`x`$. For a fixed value of the slow variable $`x`$, the eigen-state $`|n[x]`$ and the corresponding eigen-values $`V_n[x]`$ are determined by the eigen- equation $`h(q,x)|n[x]=V_n(x)|n[x]`$. Regarding $`x`$ and $`q`$ as the slow and fast variables respectively in the BO adiabatic approach , we approximately obtain the complete set {$`\varphi _{n,\alpha }(x)|n[x]`$} of eigenstates of the total system, where $`\varphi _{n,\alpha }(x)`$ come from the eigen-equation $`H_n\varphi _{n,\alpha }(x)=E_{n,\alpha }\varphi _{n,\alpha }(x)`$ and $`H_n=p^2/M+V_n[x]`$ is the effective Hamiltonian correlated to the internal state $`|n[x]`$. Here, we do not consider the induce gauge potential. Then,we can see how the “measurement-like process” naturally appears as a result of the adiabatic dynamic evolution.
In fact, under the BO approximation, we can expand the factorized initial state $`|\mathrm{\Psi }(0)`$ $`=`$ $`|x|\varphi `$ in terms of the adiabatic basis {$`x|n,\alpha \varphi _{n,\alpha }(x)|n[x]`$} and then we obtain the total wave function
$$|\mathrm{\Psi }(t)=\underset{n,\alpha }{}\varphi _{n,\alpha }|xn[x]|\varphi e^{iE_{n,\alpha }t}|n,\alpha $$
$$=\underset{n}{}n[x]|\varphi 𝑑x^{}x^{}|e^{iH_nt}|x|x^{}|n[x^{}]$$
$`(7.5)`$
where we have used the single-component completeness relation $`_\alpha |\varphi _{n,\alpha }\varphi _{n,\alpha }|=1.`$Generally, the propagator $`K(x^{},x,t)=x^{}|e^{iH_nt}|x|x^{}`$ is not diagonal for $`|x`$ is not an eigen-state of $`H_n.`$ However, in the large mass limit , we can prove that, to the first order approximation , $`K(x^{},x,t)`$ takes a diagonal form proportional to a $`\delta function.`$ In fact,in this case, the kinetic term $`p^2/2M`$ can be regarded as a perturbation in comparison with the effective potential $`V_n(x).`$ Using Dyson expansion to the first order of $`\frac{1}{M}`$, we have
$$e^{iH_nt}=e^{iV_nt}\left(1i_0^te^{iV_nt^{}}\frac{p^2}{2M}e^{iV_nt^{}}𝑑t^{}+\mathrm{}\right)$$
$$=e^{iV_nt}\left(1i\frac{p^2t^2}{2M}+i\frac{t^2}{4M}(p_xV_n+[_xV_n]p)\frac{it^3_xV_n^2}{6M}+cdots\right)$$
$`(7.6)`$
Since $`x^{}|P^n|xf(x^{})𝑑x=0`$ for n=1,2,…, we have
$$K(x^{},x,t)=e^{iV_n[x]t}[\delta (xx^{})+\frac{i}{2M}_0^t𝑑\tau e^{iV_n[x^{}]\tau }\frac{^2}{x^2}\delta (xx^{})e^{iV_n(x)\tau }]$$
$`(7.7)`$
We notice this simple result has the following physical explanation: the evolution state of a heavy particle for very large $`M`$, which is almost steady, is approximately an eigenstate of the position operator if it is initially in a state with a fixed position. Then,it follows that, in the large-mass limit, the wave function $`|\mathrm{\Psi }(t)`$ can be factorized approximately: $`|\mathrm{\Psi }(t)=|xS(x,t)|\varphi `$ where the entangling $`Smatrices`$
$$S(x,t)=\underset{n,}{}e^{ihV_nt}|n[x]n[x]|$$
$`(7.10)`$
are defined in terms of the adiabatic projection $`|n[x]n[x]|`$.
According to our previous argument about the factorized structure of $`Smatrix`$ in the dynamic theory of quantum measurement \[ \], if the internal degree of freedom has many components, e.g.,if $`q=(q_1,q_2,\mathrm{}q_N)`$ ,then in their normal non-interaction modes , $`S(x;t)`$ can be factorized as:
$$S(x;t)=\underset{j=1}{\overset{N}{}}S__j(x;t)$$
$`(7.11)`$
with
$$S__j(x;t)=e^{ih_j(q_j,x)t}$$
$`(7.12)`$
with $`h(q,x)=_jh_j(q_j,x)`$. Of course in the derivation of the above factorized structure for the $`Smatrix`$ , we have made some simplifications. Roughly speaking, we have assumed that the potential takes the form of direct sum and the eigenstate the form of direct product
$$V_n=\underset{j}{}V_{nj}(q_j),|n[x]=\underset{j=1}{\overset{N}{}}|n_j[x]$$
$`(7.13)`$
neglecting the higher order terms $``$ $`O(\frac{1}{M}).`$
For the initial state $`|\varphi =_{j=1}^N`$ $`|\varphi _j`$ factorized with respect to internal components, the reduced density matrix
$$\rho (x,x^{},t)=\phi (x)\phi ^{}(x^{})F_N(x^{},x,t):$$
$`(7.14)`$
can be re-written in terms of the so called decoherence factor
$$F_N(x^{},x,t)=\underset{j=1}{\overset{N}{}}F^{[j]}(x^{},x,t)\underset{j=1}{\overset{N}{}}\varphi _j|S_{q_j}^{}(x^{};t)S_{q_j}(x;t)|\varphi _j.$$
$`(7.15)`$
This factor is expressed as an $`N`$-multiple product of the single decohering factors $`F^j(x,x^{})=`$ $`\varphi _j|S_{q_j}^{}(x^{};t)S_{q_j}(x;t)|\varphi _j`$with norms less than unity. Thus in the macroscopic limit $`N\mathrm{}`$ , it is possible that $`F_N(x^{},x,t)`$ $`0,`$ for $`x^{}x`$. In fact, this factor reflects almost all the dynamic features of the influence of the fast part on the slow part. Physically, an infinite $`N`$ means that the object is macroscopic since it is made of infinite number of particles in that case. On the other hand, the happening of decoherence at infinite $`N`$ manifests a transition of the object from the quantum realm to the classical realm.Here,as expected,the physical picture is consistent.
As to the localization problem raised by Einstein and Born , we , based on the above argument, comment that one can formally write down the wave function of a macroscopic object as an narrow pure state wave packet, but it is not the whole of a real story. Actually, the statement that an object is macroscopic should physically imply that it contains many particles. So a physically correct description of its state must concern its internal motions coupling to the collective coordinates (e.g., its center-of-mass) . Usually, one observe this collective coordinate to determine whether two spatially-localized wave packets can interfere with each other. If there does not exist such interference, one may say that, the superposition of two narrow wave packets for the macro-coordinate is no longer a possible pure state of the macroscopic object. Indeed, because the “which-way” information of the macro-coordinate is recorded by the internal motions of particles making up the macroscopic object, the induced decohernce must destruct the coherence in the original superposition so that the state of the macroscopic object is no longer pure.
The present argument also provides a possible solution for the Schroedinger cat paradox. If we consider the Schroedinger cat as a macroscopic object consisting of many internal particles, then we can never observe anything corresponding to the interference between the dead and the living cats because the macroscopically- dead and the macroscopically- living states of the cat are correlated to the corresponding internal states. In this sense, we conclude that the Schroedinger cat paradox is not a paradox at all in practice. Rather, it essentially arises from overlooking the internal motions of a macroscopic cat or the multi-particle scattering off it. Turning to the problem of quantum coherence of a subsytem within the total system, from the above argument we also conclude that the classical characters of both the quantum system and the large system entangled with it are “correlated” physically: when the large system transits from quantum to classical realms, the quantum system has to act in the same way. In other words, you can never see a coherent superposition of microscopic states entangled to a live or a dead cat’s states if the Schrodinger cat is classical. In the presence of a classical cat, the quantum system entangling with it should lose its own coherence. Actually, if , in a classical manner,one asks experimentally what a quantum system really does , then the quantum system would behave physically like a classical object . This is just the quantum mystery physicists have to face.
To make a deeper elucidation of the above general arguments about the localization of a macroscopic object of mass $`M`$, we model the macroscopic object as consisting of $`N`$ two level particles, which are fixed at certain positions to form a whole without internal spatial motion. The collective position $`x`$ is taken to be its mass-center or any reference position on it while the internal variables are the quasi-spins associated with two level particles. Generally, if we assume the back-action of the internal variables on the collective position is relatively small, the model Hamiltonian can be written as
$$H=\frac{P^2}{2M}+\underset{j=1}{\overset{N}{}}[f_j(x)|e_jg_j|+f_j^{}(x)|g_je_j|]+\underset{j=1}{\overset{N}{}}\omega _j[|e_je_j||g_jg_j|]$$
$`(7.16)`$
where $`|g_j`$ and $`|e_j`$ are the ground and the excited states of the $`j`$ ’th particle and $`f_j(x)`$ denote the position-dependent couplings of the collective variable to the internal variables. Let $`l_j`$ be the relative distance between the $`j`$ ’th particle and the reference position $`x.`$We can further assume $`f_j(x)=f(x+l_j).`$ Physically,we may think that these couplings are induced by an inhomogeneous external field , e.g., they may be the electric dipole couplings of two-level atoms in an inhomogeneous electric field.
We remark that the above model enjoys some universality under certain conditions, compared with various environment models inducing both dissipation and decoherence of quantum processes. In fact, Caldeira and Leggett have pointed out that any environment weakly coupling to a system may be approximated by a bath of oscillators under the condition that “each environmental degree of freedom is only weakly perturbed by its interaction with the system”. We observe that any linear coupling only involves transitions between the lowest two levels (ground state and the first excitation state) of each harmonic oscillator in the perturbation approach though it has many energy levels. Therefore in such a case we can also describe the environment as a combination of many two level subsystems without losing generality .To some extent, these arguments justify our choosing the two level subsystems to model the internal motion of the macroscopic object.We will soon see its advantage:the localization characters can be manifested naturally and clearly.
Now let us calculate the $`S__j(x;t)`$ for this concrete model. The single-particle Hamiltonian $`h_j(x)=\omega _j(|e_je_j||g_jg_j|)+(f_j(x)|e_jg_j|+h.c)`$ has the $`x`$-dependent eigenvalues
$$V_{jc}=n\mathrm{\Omega }_j(x)\pm \sqrt{|f_j(x)|^2+\omega _j^2}(n=\pm )$$
$`(7.17)`$
and the corresponding eigen-vectors$`|n_j[x]`$ are
$$|+_j[x]=\mathrm{cos}\frac{\theta _j}{2}|e_j+\mathrm{sin}\frac{\theta _j}{2}|g_j,$$
$`(7.18)`$
$$|_j[x]=\mathrm{sin}\frac{\theta _j}{2}|e_j\mathrm{cos}\frac{\theta _j}{2}|g_j,$$
$`(7.19)`$
where $`\mathrm{tan}\theta _j=\frac{f_j(x)}{\omega _j}.`$ Then we explicitly the corresponding single-particle $`Smatrix`$
$$S__j(x;t)=\left(\begin{array}{cc}\mathrm{cos}(\mathrm{\Omega }_jt)i\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{cos}\theta _j,& i\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{sin}\theta _j\\ i\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{sin}\theta _j,& \mathrm{cos}(\mathrm{\Omega }_jt)+i\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{cos}\theta _j\end{array}\right)$$
$`(7.20)`$
Here in the derivation we have used the formula $`\mathrm{exp}[i\stackrel{}{\sigma }\stackrel{}{A}]=\mathrm{cos}A+i\stackrel{}{\sigma }\stackrel{}{n_A}\mathrm{sin}A`$ for a given vector $`\stackrel{}{A}`$ of norm $`A`$ along the direction $`\stackrel{}{n_A}.`$Having obtained the above analytic results about $`Smatrix,`$ we can further calculate the single-particle decoherence factors $`F^{[j]}(x^{},x,t)g_j|S__j^{}(x^{};t)S__j(x;t)|g_j`$for a given initial state $`|\varphi =_{j=1}^N`$ $`|g_j`$. For simplicity we use the notation $`f(x^{})=f^{}`$ .We have
$$F^{[j]}(x^{},x,t)=\{\mathrm{sin}(\mathrm{\Omega }_j^{}t)\mathrm{sin}\theta _j^{}\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{sin}\theta _j+$$
$$\mathrm{cos}(\mathrm{\Omega }_j^{}t)\mathrm{cos}(\mathrm{\Omega }_jt)+\mathrm{sin}(\mathrm{\Omega }_j^{}t)\mathrm{cos}\theta _j^{}\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{cos}\theta _j^{}\mathrm{cos}\theta _j$$
$`(7.22)`$
$$+i\{\mathrm{cos}(\mathrm{\Omega }_j^{}t)\mathrm{sin}(\mathrm{\Omega }_jt)\mathrm{cos}\theta _j\mathrm{sin}(\mathrm{\Omega }_j^{}t)\mathrm{cos}\theta _j^{}\mathrm{cos}(\mathrm{\Omega }_jt)\}\}$$
In the weakly coupling limit with $`g_j\omega _j`$ and the coupling $`f_jg_jx`$,we have $`\mathrm{sin}\theta _j\theta _j\frac{f_j}{\omega _j},\mathrm{cos}\theta _j1\frac{1}{2}\theta _j^2`$ and $`\mathrm{\Omega }_j\omega _j.`$Thus, the decohering factors can be simplified as
$$F^{[j]}(x^{},x,t)1(xx^{})^2\frac{|g_j|^2}{2\omega _j^2}\mathrm{sin}^2(\omega _jt)+\frac{i|g_j|^2}{4\omega _j^2}\{x^2x^2)\mathrm{sin}(2\omega _jt)$$
$`(7.23)`$
Consequently, the temporal behavior of the decoherence is determined by
$$F(x^{},x,t)=\mathrm{exp}\left\{(xx^{})^2\underset{j=1}{\overset{N}{}}\frac{|g_j|^2}{2\omega _j^2}\mathrm{sin}^2(\omega _jt)+(x^2x^2)\underset{j=1}{\overset{N}{}}\frac{i|g_j|^2}{4\omega _j^2}\mathrm{sin}(2\omega _jt)\right\}$$
$`(7.24)`$
In the case of continuous spectrum, the sum $`R(t)=`$ $`\underset{j=1}{\overset{N}{}}\frac{g_j^2}{2\omega _j^2}\mathrm{sin}^2\left(\omega _jt\right)`$can be re- expressed in terms of a spectrum distribution $`\rho (\omega _k)`$ as $`R(t)=_0^{\mathrm{}}\frac{\rho (\omega _k)g_k^2}{2\omega _k^2}\mathrm{sin}^2\omega _kd\omega _k.`$From some concrete spectrum distributions, interesting circumstances may arise. For instance, when $`\rho (\omega _k)=\frac{4}{\pi }\gamma /g_k^2`$ the integral converges to a negative number proportional to time t , precisely, $`S(t)=\gamma t`$ .
Therefore, our analysis recovers the result
$$\rho (x,x^{},t)=\phi (x)\phi ^{}(x^{})e^{\gamma t(xx^{})^2}\mathrm{exp}[i\pi (x^2x^2)s(t)]$$
$`(7.25)`$
for the reduced density matrix of the macroscopic object , which was obtained by Joos and Zeh through the multi particle external scattering mechanism and by Zurek separately through Markov master equation. Here, $`s(t)=_{j=1}^N\frac{\mathrm{sin}(2\omega _jt)}{\pi \omega _j^2}`$ is a time-dependent multi-period function. This shows that the norm of the decoherence factor is exponentially decaying and as $`t\mathrm{},`$ the off-diagonal elements of the density matrix vanish simultaneously! Due to the presence of the oscillating factor $`s(t)`$ of multi-period, $`\rho (x,x^{},t)`$ seems very complicated.But on the other hand, the simple decaying norm of $`\rho (x,x^{},t)`$ can well serve to describe decoherence of the macroscopic object. Consider now a similar example by Joos and Zeh .We take a coherent superposition of two Gaussian wave packets of width $`d`$
$$\phi (x)=\frac{1}{\sqrt[4]{8\pi d^2}}\left\{\mathrm{exp}\left(\frac{(xa)^2}{4d^2}\right)+\mathrm{exp}\left(\frac{(x+a)^2}{4d^2}\right)\right\}$$
$`(7.26)`$
The norm of the corresponding reduced density matrix
$$|\rho (x,x^{},t)|=\underset{k,l=0}{\overset{1}{}}P_{kl}(x,x^{},t)$$
$`(7.27)`$
contains 4 peaks
$$P_{11}(x,x^{},t)=\frac{1}{\sqrt{8\pi d^2}}e^{\gamma t(xx^{})^2}\mathrm{exp}[\frac{(xa)^2}{4d^2}\frac{(x^{}a)^2}{4d^2}]$$
$$P_{10}(x,x^{},t)=\frac{1}{\sqrt{8\pi d^2}}e^{\gamma t(xx^{})^2}\mathrm{exp}[\frac{(xa)^2}{4d^2}\frac{(x^{}+a)^2}{4d^2}]$$
$$P_{01}(x,x^{},t)=\frac{1}{\sqrt{8\pi d^2}}e^{\gamma t(xx^{})^2}\mathrm{exp}[\frac{(x+a)^2}{4d^2}\frac{(x^{}a)^2}{4d^2}]$$
$`(7.28)`$
$$P_{00}(x,x^{},t)=\frac{1}{\sqrt{8\pi d^2}}e^{\gamma t(xx^{})^2}\mathrm{exp}[\frac{(x+a)^2}{4d^2}\frac{(x^{}+a)^2}{4d^2}]$$
centering respectively around the points $`(a,a),(a,a),(a,a)`$ and $`(a,a)`$ on $`xx^{}`$-plane. The heights are respectively $`1/\sqrt{8\pi d^2},e^{4\gamma ta^2}/\sqrt{8\pi d^2},e^{4\gamma ta^2}/\sqrt{8\pi d^2}`$and $`1/\sqrt{8\pi d^2})`$. Obviously, two peaks with centers at $`(a,a)`$ and $`(a,a)`$ decays with time while the other two keep their heights constant . Fig.6. shows this time-dependent configuration at t=0,and a finite t. As $`t\mathrm{}`$, two off-diagonal terms $`P_{10}`$ and $`P_{01}`$decay to zero so that the interference of the two Gaussian wave packets are destroyed . In this sense, we say that the pure state $`\rho (x,x^{},t=0)=𝑑x\phi (x)\phi (x^{})|xx^{}|`$ becomes a mixture
$$\rho (t)=𝑑x\phi (x)\phi ^{}(x)|xx|$$
$`(7.29)`$
in $`x`$representation.
Interference of two plane waves of wave vector $`k_1,k_2`$ provides us another simplest example. Without decoherence induced by its internal motions or the external scattering , their coherent superposition $`\phi (x)=\sqrt{\frac{1}{4\pi }}[e^{ik_1x}+e^{ik_2x}]`$ yields a spatial interference described by the reduced density matrix
$$\rho _0(x,x^{},t)=\frac{1}{4\pi }\{e^{ik_1(xx^{})}+e^{ik_2(xx^{})}+$$
$$\mathrm{exp}[i(\frac{k_1^2tk_2^2t}{2m}+k_2xk_1x^{})]+\mathrm{exp}[i(\frac{k_2^2tk_1^2t}{2m}+k_1xk_2x^{}]\}$$
$`(7.30)`$
Under the influence of internal motions , it becomes
$$\rho (x,x^{},t)\rho _0(x,x^{},t)e^{\gamma t(xx^{})^2}$$
for large mass. We see that the difference created by decoherence is only reflected in the off-diagonal elements,and the pure decoherence (without dissipation) does not destroy the interference pattern described by the diagonal term
$$\rho (x,x,t)=\rho _0(x,x,t)=\frac{1}{2\pi }\{1+\mathrm{cos}[\frac{k_1^2tk_2^2t}{2m}+(k_2k_1)x]\}$$
$`(7.32)`$
This simple illustration tells us that the present quantum decoherence mechanism may not have to do with the interference pattern of the first order coherence, but it does destroy the higher order quantum coherence: $`\rho (x,x^{},t)0`$ as $`t\mathrm{}.`$ In fact, Savage, Walls and Yu have shown that, due to the induced loss of energy, quantum dissipation is responsible for the disappearance of the interference pattern of the first order coherence . The influences of internal motions or external scattering on the decoherence of a macroscopic object may be very complicated. Intuitively, these dynamic effects should depend on the details of interaction between the collective variables and the internal and external degrees of freedom. Pratically,we can classify these influences into two species,namely, quantum dissipation and quantum decoherence, and then study them separately by different models.
## VIII Concluding Remarks
We remark that the quantum decoherence of a small system, resulting from a transition of the entangling “large system” from quantum to classical, is certainly an irreversible process. This is because the density matrices $`\rho _s(0)`$ and $`\rho _s(t)`$ have different ranks for $`t0`$. Thus they can not be transformed into each other through an unitary time-evolution matrix. If we only consider a closed system, as a postulate with certain classical elements in it, its WPC or quantum decoherence can not be derived from Schröedinger equation based on the basic laws of quantum mechanics. Since quantum mechanics was founded, physicists have wished to add this WPC postulate to the axiom system of quantum mechanics. von Neumann and Wigner made the first attempt and considered the measurement detector plus the measured system as a total system called a“universe” satisfying Schröedinger equation. They hoped that, projected on the system, the evolution of the “universe” leads to wave packet collapse naturally. However,because it did not take the macroscopic or classical character of detector into account , this approach brings philosophical difficulty: If the observation of the detector force the measured system to decohere, the detector must decohere in advance. So the second detector is needed to monitor the first one, and the third one is needed to monitor the second one and so on. By this argument in logic, a chain of detectors should be introduced in sequence (we usually call it von Neumann ’s chain),and at the end of this von Neumann ’s chain, there should exist a pair of eyes as a special detector ,which is required to be classical.
In 1972, to avoid the introduction of this chain of detectors, Hepp and Coleman raised a dynamical description for the WPC via a simple exactly-solvable model. They emphasized that, if the macroscopic limit of the first detector is considered appropriately, detectors other than the first one are not necessary. Using the macroscopic character of the first detector is crucial to the solution of this problem. Later Namiki, Nakazato and Pascazio et al generalized this work to put forward various new models for quantum measurement. In 1993, after carefully analyzing these models and taking the classical limit of detector into account, one of the authors (CPS) found that the essence implied by these models is a factorization structure . By exact- solvable models and approximately-solvable models as well,it is shown that when the effective evolution of the detector can be factorized,in the macroscopic limit that the number of particles composing the detector approachs infinity, quantum decoherence or WPC will appear naturally.
Previously we also associated quantum decoherence problem with the requirement that the result of a measurement should be macroscopically observable so that an ideal entanglement happens dynamically . But all discussions about the interaction induced quantum decoherence strongly rely on the particular forms of interaction, namely, the interaction $`H_I(q,x)`$ of non-demolition type, which depends on the variable $`q`$ of the measured system, but commutes with its free Hamiltonian. Therefore, a fatal defect inherent in the previous works,including our own works,concerning the study of environment induced decoherence based on quantum measurement theory is that the question why should nature choose such a particular form of interaction remains unanswered. In some sense, the present work in this paper has well tackled this problem.Indeed, through the B-O adiabatic separation of the quantum and quasi-classical variables,we have demonstrated that in the adiabatic limit, the effective interaction reduced from a quite general coupling just takes such a particular form.
In a wide sense, the adiabatic entanglement can be well understood in the picture of coupled channels , which is an extensive generalization of B-O approximation. Consider a total system whose wave function depends on two set of variables, $`q`$ and $`x`$. Let $`Q`$ be an operator only acting on the function of $`q`$ and has a complete set of eigen-vectors $`\{|n\}`$ with the corresponding eigenvalues $`v_n.`$ Since $`\{|n\}`$ forms a complete basis of the Hilbert space of all functions of $`q,`$ the total eigenfunction $`\mathrm{\Psi }_E(x,q)`$ of the full Hamiltonian $`H=H_E(x)+H_s(q)+H_I(x,q)`$ with eigen-value $`E`$ can be regarded as a function of $`q`$ for a given $`x`$ and then can be expressed as $`\mathrm{\Psi }_E(x,q)=\varphi _n(x)|n`$. The channel wave function $`\varphi _n(x)`$ is defined by the coupled channel equations
$$H_{nn}(x)\varphi _n(x)+\underset{mn}{}H_{nm}(x)\varphi _m(x)=E\varphi _n(x)$$
$`(8.1)`$
The matrix elements $`H_{mn}(x)=m|H|n_q`$ are defined in terms of the $`q`$-function space “integral”.Under a certain condition, if the off-diagonal elements can be neglected physically, an effective non-demolition Hamiltonian $`H_{eff}=H_{Eeff}(x)+H_{seff}+H_{in}(x)`$:
$$H_{Eeff}=diag.[H_{11}^E(x),H_{22}^E(x),\mathrm{}..,H_{dd}^E(x)]$$
$$H_{seff}=diag.[\lambda _1,\lambda _2,\mathrm{}.,\lambda _d],$$
$`(8.2)`$
$$H_{in}(x)=diag.[H_{11}^s(x),H_{22}^s(x),\mathrm{}..,H_d^s(x)]$$
can be partially diagonalized in the ‘channel space’. Here $`H_{mm}^A(x)=m|H_A|m_q`$ for $`A=E,S`$ and $`\lambda _m=m|H_s(q)|m_q`$ are constants. Obviously, the non-demolition condition \[$`H_{seff},`$ $`H_{in}(x)]=0`$ holds as $`H_{seff}`$ is a constant matrix. In the B-O approximation the channel operator $`Q`$ is taken to be $`Q[x]=H_s(q)+H_I(x,q),`$which is parametrized by $`x.`$ The adiabatic condition maintains that, only the diagonal elements play a dominant role and the off-diagonal elements can be neglected for very small channel-channel coupling . Therefore, it can be concluded that there may exist a more universal mechanism beyond B-O approximation to realize the quantum decoherence dynamically originated from the basic interaction , which is related to the theory of coupled channels.
Finally we point out that the presence of non-demolition interaction is only a necessary condition for quantum decoherence to appear.Sufficient conditions should include the requirement that the large system be classical so that its final states could be orthogonal to one another. In this paper,we have regarded the spin-system with a very large spin and the harmonic oscillator initially in a coherent state as classical objects. Then within the semi-classical framework,even in the case of a general potential motion, we are able to relate the macroscopic distinguishibility of the quantum states of the large system to its classical limit behaviors. However, there are still vague points in the definition of the quantum-classical division for the large system. This problem is deeply rooted in the following more fundamental and more challenging issue: why or in what sense does a general large system behave classically. If we imagine that, beside the considered quantum system, there is another system coupling with the large system to decohere it, then the present problem will be trapped into an evil logic chain. One notices the difficulty here is very similar to that faced by von Neumann and Wigner about sixty years ago . Though new experiments have been revitalizing the study of decoherence problem and progress is being made,it seems that there is still a long way to go to finally understand quantum irreversible process completely.To reach this goal,one should first find a satisfactory definition for the so called quantum-classical boundary.At present it is very unclear to us how to do this without recourse to particular physical systems.
## Acknowlegement
This work is supported by direct grant (Project ID:2060150) from The Chinese University of Hong Kong. It is also partially supported by the NFS of China. One of the authors (CPS) wishs to express his sincere thanks to P.T.Leung, C.K.Law and K.Young for many useful discussions.
Appendix. Wei-Norman Algebraic Solution For $`H=\frac{1}{2M}(𝐩A)^2+fx`$
Let $`U(t)`$ be the evolution operator of a quantum system with the effective Hamiltonian
$$H=\frac{1}{2m}(𝐩A)^2+f𝐱$$
$`(a1)`$
where $`A`$ is a constant induced gauge potential . Neglecting the constant term, we can rewrite the effective Hamiltonian
$$H=\frac{1}{2m}𝐩^2\frac{A}{M}𝐩+fx$$
$`(a2)`$
as an element of the Lie algebra £ generated by $`\{p^2,p,x,1\}.`$
According to Wei-Noman’s algebraic theorem , a solution of the Schroedinger equation for the evolution operator $`U(t)`$ must be an element belonging to the Lie group related to the Lie algebra £. Since the commutation relations are closed among these four elements, the solution $`U(t)`$ is assumed to have a factorized form
$$U(t)e^{\alpha \left(t\right)𝐩^2}e^{\beta \left(t\right)𝐩}e^{\gamma \left(t\right)𝐱}e^{\mu (t)}$$
$`(a3)`$
Its Schroedinger equation defines an solvable system of equations about the time-dependent parameters $`\alpha \left(t\right),\beta \left(t\right),\gamma \left(t\right)`$ and $`\mu (t):`$
$$\frac{d}{dt}\alpha \left(t\right)=\frac{i}{2M}$$
$$\frac{d}{dt}\beta \left(t\right)2i\alpha \left(t\right)\frac{d}{dt}\gamma \left(t\right)=\frac{iA}{M}$$
$`(a4)`$
$$\frac{d}{dt}\gamma \left(t\right)=if$$
$$\frac{d}{dt}\mu \left(t\right)i\beta \left(t\right)\frac{d}{dt}\gamma \left(t\right)=0$$
The solution is
$$\alpha \left(t\right)=\frac{it}{2M}=i\stackrel{\mathrm{~}}{\alpha }(t)$$
$$\beta \left(t\right)=\frac{ift^2}{2M}+\frac{iAt}{M}=i\stackrel{\mathrm{~}}{\beta }\left(t\right)$$
$`(a6)`$
$$\gamma \left(t\right)=itf$$
$$\mu \left(t\right)=\frac{if^2t^3}{6M}+\frac{iAft^2}{2M}=i\stackrel{\mathrm{~}}{\mu }\left(t\right)$$
The action of the evolution operator
$$U_k(t)e^{\stackrel{\mathrm{~}}{i\alpha }\left(t\right)𝐩^2}e^{\stackrel{\mathrm{~}}{i\beta }\left(t\right)𝐩}e^{iftx}e^{i\stackrel{\mathrm{~}}{\mu }(t)}$$
$`(a3)`$
transforms the initial state in momentum representation
$$\phi (p,0)=p|\phi (0)=(\frac{2a^2}{\pi })^{\frac{1}{4}}e^{a^2p^2}$$
$`(a3)`$
into
$$\phi (p,t)=p|U(t)|\phi (0)=(\frac{2a^2}{\pi })^{\frac{1}{4}}e^{i\stackrel{\mathrm{~}}{\mu }(t)}e^{\stackrel{\mathrm{~}}{i\alpha }\left(t\right)p^2}e^{\stackrel{\mathrm{~}}{i\beta }\left(t\right)p}e^{a^2(p+ft)^2}$$
$`(a3)`$
Then in momentum representation we can easily calculate the overlap of two entangled states:
$`|\phi ^{}(t)|\phi (t)|`$ $`=`$ $`\mathrm{exp}[a^2t^2(f^2+f^2)]\times `$
$`|\mathrm{exp}\{{\displaystyle \frac{[\stackrel{\mathrm{~}}{\beta }\left(t\right)\stackrel{}{\stackrel{\mathrm{~}}{\beta }}\left(t\right)+i2a^2t(f+f^{})]^2}{8a^2}}\}|`$ |
warning/0001/hep-th0001065.html | ar5iv | text | # On Domain-wall/QFT dualities in various dimensions
## Introduction
Anti-de Sitter (AdS) gravity has attracted much attention due to the conjectured correspondence to a conformal field theory (CFT) on the boundary of the AdS spacetime leading to the so-called AdS/CFT correspondence. (for a review, see ). The AdS/CFT correspondence has been extended to a DW/QFT correspondence for D$`p`$–branes in ten dimensions. In this talk we extend the discussion of to general two-block $`p`$–branes in various dimensions.
In Section 1 we first discuss some general facts on Domain-Wall and anti-de Sitter spacetimes. In Section 2 we calculate the near-horizon geometries of a generic $`p`$–brane. In Section 3 we discuss the field theory limit for the general case. The formulae we give in this Section are applied to the special case of 0–branes we discuss in the Section 4. In this final Section we discuss the quantum mechanics of 0–branes, or extreme black holes, in various dimensions.
## 1 Domain-Walls and anti-de Sitter spacetimes
Domain-wall (DW) spaces, i.e. spaces of co-dimension 1, occur as solutions to the equations of motion of a (super-)gravity action with a dilaton scalar $`\varphi `$ and a (D-1)–form gauge potential (for a review, see ). They are $`p`$–branes with worldvolume dimension $`p+1`$ which is one less than the dimension $`D`$ of their target spacetime, i.e. $`D=p+2`$. The part of the supergravity action needed to describe the DW solution is given by (we use the Einstein frame and mostly plus signature):
$$S(D,b)=\mathrm{d}^Dx\frac{\sqrt{g}}{2\kappa _D^2}\left[R\frac{4}{D2}(\varphi )^2\frac{1}{2D!}g_s^{2k}\left(\frac{e^\varphi }{g_s}\right)^bF_D^2\right],$$
(1)
where $`\kappa _D`$ is the gravitational coupling constant with (ignoring constants)
$$\kappa _D^2\mathrm{}_s^{D2}g_s^2,$$
(2)
and $`b`$ is the dilaton coupling parameter. We also introduced a parameter $`k`$ defined by
$$k(D,b)=\frac{b}{2}+2\frac{D1}{D2}.$$
(3)
Solving the equations of motions following from (1) one finds:
$`\mathrm{d}s^2`$ $`=`$ $`H^{\frac{4\epsilon }{(D2)\mathrm{\Delta }_{\mathrm{DW}}}}\mathrm{d}x_{D1}^2+H^{\frac{4\epsilon (D1)}{(D2)\mathrm{\Delta }_{\mathrm{DW}}}2(\epsilon +1)}\mathrm{d}y^2,`$
$`e^\varphi `$ $`=`$ $`g_sH^{\frac{(D2)b\epsilon }{4\mathrm{\Delta }_{\mathrm{DW}}}},`$ (4)
$`g_s^kF_{01\mathrm{}D2y}`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{4}{\mathrm{\Delta }_{\mathrm{DW}}}}}_yH^\epsilon ,`$
with $`\epsilon `$ a parameter and where $`g_s`$ and $`\mathrm{\Delta }_{\mathrm{DW}}`$ (which is invariant under dimensional reduction) are defined by
$`g_s`$ $`=`$ $`e^{\varphi (H=1)},`$
$`\mathrm{\Delta }_{\mathrm{DW}}(D,b)`$ $`=`$ $`\frac{1}{8}(D2)b^22{\displaystyle \frac{D1}{D2}}.`$ (5)
The function $`H`$ is harmonic on the 1-dimensional transverse space with coordinate $`y`$:
$`H`$ $`=`$ $`c+Q_1y,y>0,`$
$`H`$ $`=`$ $`c+Q_2y,y<0,`$ (6)
with $`c,Q_1,Q_2`$ constant. The domain-wall is positioned at the discontinuity $`y=0`$. Different choices of $`\epsilon `$ correspond to different choices of coordinates and lead to different expressions for the metric . For instance, one can choose $`\epsilon `$ such that the powers of the harmonic function become equal or opposite.
Shifting the position of the domain wall to infinity, e.g. $`y+\mathrm{}`$ allows us to discard the constant $`c`$ in the harmonic function and, furthermore there is only one side of the domain wall. We can eliminate $`\epsilon `$ if we define a mass parameter $`m`$ by :
$$Q\epsilon =m,$$
(7)
with $`Q=Q_1`$. Making a co-ordinate transformation and by going to the so-called “dual frame” metric $`g_{}`$, indicated by a subscript star,
$$Qy=e^{Q\lambda },g_{}=e^{b\varphi }g_E,$$
(8)
we find for the solution (1)
$$\mathrm{d}s_{}^2=e^{2m\lambda (\frac{2+\mathrm{\Delta }_{\mathrm{DW}}}{\mathrm{\Delta }_{\mathrm{DW}}})}\mathrm{d}x_{D1}^2+\mathrm{d}\lambda ^2,\varphi =\lambda \frac{(D2)bm}{4\mathrm{\Delta }_{\mathrm{DW}}}.$$
(9)
In the dual frame (8) the domain-wall solution (9) describes an $`\mathrm{AdS}_D`$ spacetime with a linear dilaton, the latter breaking the full conformal structure of the AdS spacetime. This observation will be the key to generalizing the AdS/CFT duality to a DW/QFT duality.
## 2 Near-horizon Geometries of $`p`$–branes
Our starting point is the $`D`$-dimensional action
$$S(D,a,p)=\mathrm{d}^Dx\frac{\sqrt{g}}{2\kappa _D^2}\left[R\frac{4}{D2}(\varphi )^2\frac{g_s^{2k}}{2(d+1)!}\left(\frac{e^\varphi }{g_s}\right)^aF_{d+1}^2\right],$$
(10)
which contains three independent parameters: the target spacetime dimension $`D`$, the dilaton coupling parameter $`a`$ and a parameter $`p`$ specifying the rank $`Dp2`$ of the field strength $`F`$. The parameter $`k`$ is a generalization of (3) and is given by
$$k(D,a,p)=\frac{a}{2}+2\frac{p+1}{D2}.$$
(11)
We have furthermore introduced two useful dependent parameters $`d`$ and $`\stackrel{~}{d}`$ which are defined by
$$\{\begin{array}{ccc}\hfill d& =& p+1\mathrm{dimension}\mathrm{of}\mathrm{the}\mathrm{worldvolume},\hfill \\ \hfill \stackrel{~}{d}& =& Dp3\mathrm{dimension}\mathrm{of}\mathrm{the}\mathrm{dual}\mathrm{brane}\mathrm{worldvolume}.\hfill \end{array}$$
(12)
We next consider the following class of diagonal “two-block” $`p`$–brane solutions (using the Einstein frame)<sup>1</sup><sup>1</sup>1For later convenience, we give the solution in terms of the magnetic potential of rank $`Dp3`$. The p-brane solution is electrically charged with respect to the p+1–form potential.:
$`\mathrm{d}s^2`$ $`=`$ $`H^{\frac{4\stackrel{~}{d}}{(D2)\mathrm{\Delta }}}\mathrm{d}x_d^2+H^{\frac{4d}{(D2)\mathrm{\Delta }}}\mathrm{d}x_{\stackrel{~}{d}+2}^2,`$
$`e^\varphi `$ $`=`$ $`g_sH^{\frac{(D2)a}{4\mathrm{\Delta }}},`$ (13)
$`g_s^{(2k)}\stackrel{~}{F}_{m_1\mathrm{}m_{\stackrel{~}{d}+1}}`$ $`=`$ $`\pm \sqrt{{\displaystyle \frac{4}{\mathrm{\Delta }}}}ϵ_{m_1\mathrm{}m_{\stackrel{~}{d}+1}m}_mH,`$
where
$$g_s^{(2k)}\stackrel{~}{F}=\left(\frac{e^\varphi }{g_s}\right)^ag_s^k{}_{}{}^{}F.$$
(14)
We use a constant, i.e. metric independent, Levi-Civita tensor. Furthermore, $`g_s=e^{\varphi (H=1)}`$ and $`\mathrm{\Delta }`$ is a generalization of the $`\mathrm{\Delta }_{\mathrm{DW}}`$ in the previous section defined by
$$\mathrm{\Delta }=\frac{1}{8}(D2)a^2+\frac{2d\stackrel{~}{d}}{D2},$$
(15)
which is invariant under reductions and oxidations (in the Einstein frame). The function $`H`$ is harmonic over the $`\stackrel{~}{d}+2`$ transverse coordinates and, assuming that
$$\stackrel{~}{d}2,0,$$
(16)
(i.e. no constant or logarithmic harmonic) this harmonic function is given by
$$H=1+\left(\frac{r_0}{r}\right)^{\stackrel{~}{d}}.$$
(17)
Here $`r_0`$ is an integration constant with the dimensions of length. It is related to the mass and charge of the $`p`$–brane a follows. The mass $`\tau _p`$ per unit p–volume is given by the ADM–formula:
$`\tau _p`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _D^2}}{\displaystyle _{M^{Dp1}}}d^{Dp2}\mathrm{\Sigma }^m(^nh_{mn}_mh^b{}_{b}{}^{})`$ (18)
$`=`$ $`{\displaystyle \frac{2(Dp3)}{\mathrm{\Delta }\kappa _D^2}}r_0^{Dp3}\mathrm{\Omega }_{Dp2}.`$
On the other hand, the charge $`\mu _p`$ per unit p-volume is given, in terms of the same integration constant $`r_0`$, by the Gauss-law formula
$`\mu _p`$ $`=`$ $`{\displaystyle \frac{1}{2\kappa _D^2}}{\displaystyle (d^{Dp2}\mathrm{\Sigma })^{m_1\mathrm{}m_{Dp2}}g_s^{(2k)}\stackrel{~}{F}_{m_1\mathrm{}m_{Dp2}}}`$ (19)
$`=`$ $`\pm \sqrt{{\displaystyle \frac{\mathrm{\Delta }}{4}}}\tau _p.`$
Hence, the p–brane solution satifies the BPS bound
$$\tau _p=\sqrt{\frac{4}{\mathrm{\Delta }}}|\mu _p|.$$
(20)
To derive an expression for $`r_0`$ in terms of the string parameters $`g_s`$ and $`\mathrm{}_s`$, which fixes the scaling of $`H`$, one must add a source term to the supergravity bulk action. Using the no-force condition<sup>2</sup><sup>2</sup>2Alternatively, one can use a scaling argument, see Appendix B of . and the fact that in the string frame the electric (p+1)–form potential $`C_{p+1}`$ is proportional to $`g_s^k`$, as follows from the action (10), we find that
$$\tau _p\frac{1}{\mathrm{}_s^{p+1}g_s^k}.$$
(21)
Comparing with (18) and using (2) we deduce that, for a single brane,
$$\left(\frac{r_0}{\mathrm{}_s}\right)^{\stackrel{~}{d}}g_s^{2k}.$$
(22)
The two-block solutions (2) include the (supersymmetric) domain-wall spaces of the previous section. They correspond to the case $`\stackrel{~}{d}=1`$, $`\epsilon =1`$ and $`r_0=1/m`$. The solutions also include the known branes in ten and eleven dimensions (M2, M5, D$`p`$, F1, NS5 etc.) as well as branes in lower dimensions. If the branes under consideration preserve any supersymmetries we can set
$$\mathrm{\Delta }=\frac{4}{n},$$
(23)
where generically $`32/2^n`$ is the number of unbroken supersymmetries.
We now consider the limit for which the constant part in the harmonic function is negligible. As in the previous section we make a co-ordinate transformation and go to the dual frame
$$\left(\frac{r_0}{r}\right)=e^{\lambda /r_0}g_{}=e^{(\frac{a}{\stackrel{~}{d}})\varphi }g_E.$$
(24)
After these manipulations we can write the near-horizon metric as
$$\mathrm{d}s_{}^2=e^{2(1\frac{2\stackrel{~}{d}}{\mathrm{\Delta }})\lambda /r_0}\mathrm{d}x_d^2+\mathrm{d}\lambda ^2+r_0^2\mathrm{d}\mathrm{\Omega }_{\stackrel{~}{d}+1}^2,\varphi =\lambda \frac{(D2)a\stackrel{~}{d}}{4\mathrm{\Delta }r_0},$$
(25)
which has an $`AdS_{d+1}S^{\stackrel{~}{d}+1}`$ geometry and a linear dilaton.
Reducing over the $`\stackrel{~}{d}+1`$ angular variables of the sphere we end up with a gauged supergravity in $`d+1`$ dimensions of the form (1) supporting a domain-wall solution of the form (1). The precise relation between the parameters of the action (1) and its solution (1) in terms of those of the action (10) and its solution (25) can be found in section 2.3 of .
Summarizing, in this Section we showed that in the dual frame, defined by (24), all $`p`$–branes solutions (2) have a near horizon geometry $`\mathrm{DW}_{d+1}S^{\stackrel{~}{d}+1}`$. The domain-wall metric has all the isometries of an AdS space. These isometries are broken in the full background because of the presence of a linear dilaton.
## 3 The field theory limit
In this section we will set up the framework for the DW/QFT duality similar to the analysis of but for arbitrary dimensions. The analysis is similar to the AdS/CFT duality. There one finds (in the low-energy limit) that supergravity in the near-horizon geometry of a large number of non-dilatonic branes is dual to the conformal field theory living on these branes (which are located at the boundary of the near-horizon geometry). This so-called holographic principle lies at the heart of the AdS/CFT conjecture.
As we showed in the previous sections, the near-horizon geometry of a general $`p`$-brane in the dual-frame is equivalent to that of a non-dilatonic $`p`$-brane. It is therefore natural to assume that the duality might be extended. The presence of the dilaton makes the AdS background into a domain-wall background and the conformal field theory into a general quantum field theory, hence the name DW/QFT duality. In the following we will only consider D-branes and their intersections in lower dimensions.
One starts with a D$`p`$-brane configuration in string theory and takes the low-energy limit in which only the massless modes survive and in which the theories in the bulk and on the brane decouple. Next, one introduces a new energy-scale and a dimensionless coupling constant which are the free parameters of the decoupled theories. The energy-scale can be given in several ways:
If one probes the stacked branes by another D$`p`$-brane, a natural energy-scale is given by the mass of the endpoint of a stretched string, which acts like a W–boson (see Figure 1)
$$E_W=U=\frac{r}{l_s^2}.$$
(26)
This was the approach taken by , all branes have the same energy-scale but the disadvantage is that the near-horizon geometry when written in the $`U`$ co-ordinates does not take the same form for all branes. One can also probe the branes by a supergravity field $`\psi `$ (suppressing possible quantum numbers). By solving the wave equation for this field one finds the following energy scale
$$E_\psi =u=\frac{r^\beta }{r_0^{\beta +1}},\beta =\frac{2\stackrel{~}{d}}{\mathrm{\Delta }}1,r_0=l_s(g_s^{2k}N)^{\frac{1}{\stackrel{~}{d}}}.$$
(27)
In it was shown that this energy-scale corresponds to the so-called holographic energy-scale used in entropy calculations. Although the energy depends on the parameters of the brane solution, one finds that the near-horizon geometry (9) in these holographic $`u`$ co-ordinates is particularly simple:
$$\mathrm{d}s_{}^2=r_0^2\left[^2\left\{u^2\mathrm{d}x_d^2+\left(\frac{\mathrm{d}u}{u}\right)^2\right\}+\mathrm{d}\mathrm{\Omega }_{\stackrel{~}{d}+1}^2\right].$$
(28)
For all branes this metric has the form $`AdS(r_0)_{d+1}S^{\stackrel{~}{d}+1}(r_0)`$ differing only in the (relative) radii of the sphere and anti-de-Sitter space.
The quantum field theory on the brane is in general not explicitly known, but dimensional analysis enables us to extract some information. If we assume that the theory is given by some $`q`$-form potential on the $`d=p+1`$ dimensional worldvolume
$$S_{wv}=\tau _p\mathrm{d}^{p+1}\xi \mathrm{Tr}F_{q+1}^2,$$
(29)
then we have for the coupling constant scaling behavior
$$g_{ft}g_s^kl_s^\alpha \mathrm{with}\alpha =p2q1.$$
(30)
For 10-dimensional D$`p`$-branes we have a vector-multiplet (i.e. $`q=1`$) on the worldvolume reproducing the familiar $`\alpha =p3`$ scaling behavior for the Yang-Mills theory on the D$`p`$-brane worldvolume. The effective dimensionless coupling constant $`\lambda `$ is given by
$$\lambda _W=g_{ft}NE_W^\alpha ,\lambda _\psi =g_{ft}NE_\psi ^\alpha .$$
(31)
Since we have two different probes, we can in principle construct two different dimensionless couplings, constructed by either $`E_W`$ or $`E_\psi `$. Since both probes have a sensible interpretation, it must be possible to use them both, independently of the low-energy limit. This leads to the following constraint on the scaling behavior in terms of the brane-solution parameters :
$$\alpha =\mathrm{\Delta }\stackrel{~}{d}=a\frac{D2}{4}.$$
(32)
Combining this with the definitions of the probe-scales $`E_W`$ and $`E_\psi `$ we find the following relations between them and their corresponding dimensionless couplings $`\lambda _W`$ and $`\lambda _\psi `$
$$\frac{E_W}{E_\psi }=\lambda _W^{\frac{2}{\mathrm{\Delta }}},\lambda _W^\beta =\lambda _\psi .$$
(33)
The decoupling limit is now given by
$$\begin{array}{cc}\text{take}\hfill & E/E_s0,E_s=\mathrm{}_s^1,\hfill \\ \text{at fixed}\hfill & \{\begin{array}{cc}\lambda _W=g_{\mathrm{ft}}NE_W^\alpha \mathrm{and}E/E_\mathrm{W}\mathrm{or},\mathrm{equivalently},\mathrm{see}(\text{33}),\hfill & \\ \lambda _\psi =g_{\mathrm{ft}}NE_\psi ^\alpha \mathrm{and}E/E_\psi \hfill & \end{array}.\hfill \end{array}$$
The dimensionless coupling made from Newton’s constant in this limit becomes
$$G_N=g_s^2\left(\frac{E}{E_s}\right)^{D2}=\left(\frac{\lambda }{N}\right)^2\left(\frac{E}{E_s}\right)^{D22\alpha }$$
(34)
so that if $`\alpha `$ becomes too large one has to take $`N`$ to infinity to ensure the decoupling of gravity.
In the AdS/CFT correspondence one has in the field theory side the parameters $`(\lambda ,\frac{1}{N})`$ which correspond to the curvature and string coupling on the supergravity side. They generalize to the effective tension of a string in the dual frame and the dilaton, respectively:
$$\tau _{}l_s^2=\lambda _\psi ^{\frac{2}{\mathrm{\Delta }\beta }},e^\varphi =\frac{\tau _{}^{\frac{\stackrel{~}{d}}{2}}}{N}.$$
(35)
We now have the following ranges of validity on the two sides of the duality
* perturbative SYM:
$$\lambda _W1\{\begin{array}{ccc}\alpha <0:E_W\mathrm{}:\hfill & \text{UV-free,}\hfill & \\ \alpha >0:E_W0:\hfill & \text{IR-free.}\hfill & \end{array}$$
* classical SUGRA:
$$\{\begin{array}{ccc}\tau _{}l_s^2\hfill & 1:\hfill & \text{no stringy corrections,}\hfill \\ e^\varphi \hfill & 1:\hfill & \text{no quantum corrections.}\hfill \end{array}$$
Using the above formulae one can easily see that classical supergravity describes strongly coupled, large N field theory. However, the conformal invariance which in the AdS/CFT duality facilitates computations in the strongly coupled field theory is now broken so that any direct check of a DW/QFT duality is ruled out. For specific examples, we refer to .
As long as $`\beta `$ is positive, the supergravity and field theory are valid in different regimes (i.e. $`\lambda _\psi 1`$ and $`\lambda _W1`$) and one can find consistent phase diagrams as in . However, for negative $`\beta `$, e.g. in case of the D6-brane, this seems no longer true at first sight: $`\lambda _\psi 1`$ and $`\lambda _W1`$.
As one can see by using the relation (33) between the two couplings, this corresponds to $`\lambda _W1`$ and $`\lambda _W1`$. This matches nicely with the explanation of , namely that the low-energy Hilbert space of the field theory has two separate sectors: one describing nearby (in U co-ordinates) brane probes and one describing supergravity probes far away (also in U co-ordinates) from the brane.
Finally we note that when the dilaton becomes larger one can perform an S-duality transformation. If the curvature is small with respect to the S-dual string scale, then one can still trust the supergravity approximation.
## 4 Dynamics of 0–branes
In this Section we consider the special case of 0–branes in various dimensions. The (bosonic) Lagrangian for a particle with mass $`m`$ and charge $`q`$, moving in the string frame near-horizon background of $`N`$ stacked 0–branes reads
$$=me^\varphi \sqrt{|\dot{x}^\mu \dot{x}^\nu g_{\mu \nu }^S|}+qA_\mu \dot{x}^\mu ,$$
(36)
where the dot represents derivatives with respect to the worldline time. The D-dimensional 0–brane solution in the dual frame $`g_{\mu \nu }^{}`$ is given by the expression (28), taken for $`p=0`$. Introducing the canonical momentum $`P_\mu =\frac{}{\dot{x}^\mu }`$ one can write down the mass-shell constraint in the string frame for the probe particle as
$$(P_\mu qA_\mu )(P_\nu qA_\nu )g_S^{\mu \nu }=m^2e^{2\varphi }.$$
(37)
We would like to solve this equation for $`P_t=`$. For this purpose, we transform the mass-shell equation (37) to the dual frame and substitute the metric and gauge field of the solution (28). For the gauge field (being in the temporal gauge) we can write:
$$\frac{A_t}{u}=r_0e^{\frac{D4}{D3}\varphi }M,$$
(38)
and we find for the Hamiltonian
$``$ $`=`$ $`\left({\displaystyle \frac{u}{}}\right)^3{\displaystyle \frac{P_u^2}{2f}}+{\displaystyle \frac{u}{}}{\displaystyle \frac{g}{2f}},`$
$`f`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(qM+\sqrt{(mM)^2+({\displaystyle \frac{uP_u}{}})^2+\stackrel{}{L}^2}\right),`$
$`g`$ $`=`$ $`\stackrel{}{L}^2+(m^2q^2)M^2.`$
We can write this Hamiltonian in a rather suggestive form by making the following transformation
$$\frac{u}{}=\frac{(2)^2M}{x^2}.$$
(39)
after which the Hamiltonian takes on the following form
$``$ $`=`$ $`{\displaystyle \frac{P_x^2}{2f}}+{\displaystyle \frac{g}{2fx^2}},`$
$`f`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(q+\sqrt{m^2+\left({\displaystyle \frac{xP_x}{2M}}\right)^2+\left({\displaystyle \frac{\stackrel{}{L}}{M}}\right)^2}\right),`$
$`g`$ $`=`$ $`4\left(\stackrel{}{L}^2+(m^2q^2)M^2\right).`$
The dynamics of the system is now generated by the Poisson brackets with respect to the canonical co-ordinate and momentum $`(x,P_x)`$. The transition to quantum mechanics is the usual one. Defining two other operators
$$𝒟=\frac{1}{2}xP_x,𝒦=\frac{1}{2}fx^2,$$
(40)
we find that the system possesses a classical $`SL(2,)`$-symmetry
$$\{𝒟,\}=\{𝒟,𝒦\}=𝒦\{,𝒦\}=2𝒟.$$
(41)
This is only true when the effective coupling constant $`\lambda _\psi `$ of the model under consideration is fixed under scale transformations. Whenever we have a non-trivial dilaton the model will not be conformally invariant by itself. Only when introducing a transformation to keep $`\lambda _\psi `$ fixed under conformal transformations, will the model be invariant under what are called generalized conformal transformations .
Taking the limit $`Mm\mathrm{}`$ and at the same time $`M(mq)0`$ we find a “non-relativistic” Hamiltonian of the form:
$$=\frac{P_x^2}{2m}+\frac{\stackrel{}{L}^2}{2mx^2}.$$
(42)
This model, as first considered in , is in fact an integrable model so that it might offer a concrete test of the $`DW_2/QFT_1`$ duality in this particular limit.
## Acknowledgements
The work reported here is based upon hep-th/9907006. We thank our collaborators Klaus Behrndt and Jan Pieter van der Schaar for numerous discussions. This work is supported by the European Commission TMR programme ERBFMRX-CT96-0045, in which E.B. and R.H. are associated to the University of Utrecht. |
warning/0001/cond-mat0001046.html | ar5iv | text | # Extended 𝑑_{𝑥²-𝑦²}-wave superconductivity
## 1 Introduction
Power laws in the low-temperature asymptotic behaviour of several linear response electronic properties provide complementary evidence for $`d`$-wave symmetry of the order parameter (OP) $`\mathrm{\Delta }_\stackrel{}{k}`$ of high-$`T_c`$ superconductors Annett:90 ; Annett:96 as well as preliminary evidence of ‘exotic’ shapes in the OP of heavy fermion superconductors, such as UPt<sub>3</sub> Sigrist:91 . This has to be contrasted to an “activated” behaviour $`\mathrm{exp}(\beta \mathrm{\Delta }_{\mathrm{min}})`$, appropriate of $`s`$-wave superconductors, or, in the case of mixed symmetry, of superconductors with a non-vanishing $`s`$-wave contribution to their OP, where $`\mathrm{\Delta }_{\mathrm{min}}=\mathrm{min}_\stackrel{}{k}|\mathrm{\Delta }_\stackrel{}{k}|>0`$. In the case of a non-empty nodal manifold for the superconducting excitation spectrum $`E_\stackrel{}{k}`$, defined as the locus of points in $`\stackrel{}{k}`$-space such that $`E_\stackrel{}{k}=0`$, a large number of quasiparticles can be created near such nodes, thus dominating all the low-temperature electronic properties Lee:97 . An exact analysis allows one to relate the exponent of the leading power of the low-$`T`$ expansion of a given linear response function to the dimension of the Fermi manifold (defined as the locus of states in $`\stackrel{}{k}`$-space with vanishing dispersion relative to the Fermi level in the normal state, $`\xi _\stackrel{}{k}=0`$) and the topological nature of the nodal manifold, *viz.* a collection of points, of line segments, or of surface patches Volovik:96 ; Annett:90 .
On the basis of group theoretical arguments, the simplest choice for a $`d`$-wave gap function on a square lattice is $`\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }g(\stackrel{}{k})`$, where $`\mathrm{\Delta }`$ is a $`T`$-dependent parameter, and
$$g(\stackrel{}{k})=\frac{1}{2}(\mathrm{cos}k_x\mathrm{cos}k_y)$$
(1)
is the first basis function associated with the $`d`$-wave irreducible representation of the appropriate crystal point group, $`C_{4v}`$ Annett:90 . We remark that $`g(\stackrel{}{k})`$ is generated, together with an extended $`s`$-wave term proportional to
$$h(\stackrel{}{k})=\frac{1}{2}(\mathrm{cos}k_x+\mathrm{cos}k_y),$$
(2)
by a nearest-neighbour interaction term in real space. Here and in the following we shall measure the wavevectors in units of the appropriate inverse lattice spacings. Proportionality to Eq. (1) allows $`\mathrm{\Delta }_\stackrel{}{k}`$ to vanish linearly at a given point along the Fermi line, which for most cuprate superconductors can be modelled by the tight-binding expansion:
$$\xi _\stackrel{}{k}=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)+4t^{}\mathrm{cos}k_x\mathrm{cos}k_y\mu =0,$$
(3)
where $`t=0.25`$ eV, $`t^{}=0.45t`$ measure nearest and next-nearest neighbour hopping, respectively, and $`\mu `$ is the chemical potential.
On the other hand, increasing experimental evidence above all from angle-resolved photoemission spectroscopy (ARPES) suggests a richer structure in $`\stackrel{}{k}`$-space for the OP of the underdoped high-$`T_c`$ superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> Mesot:99 . In particular, the superconducting gap near the nodal points turns out to be flatter than predicted by the simple assumption $`\mathrm{\Delta }_\stackrel{}{k}g(\stackrel{}{k})`$ Mesot:99 . Such a feature is consistent with the observation of whole ungapped segments of the Fermi line above $`T_c`$ in the pseudogap regime of underdoped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> Norman:98 , and will of course serve as a constraint for a microscopic understanding of the pairing mechanism.
Quite remarkably, qualitatively similar deviations from a $`g(\stackrel{}{k})`$-like dispersion have been evidenced in the $`\stackrel{}{k}`$-dependence of the antiferromagnetic gap in the related insulating compounds X<sub>2</sub>CuO<sub>2</sub>Cl<sub>2</sub> (X = Ca, Sr) Ronning:98 . Such a finding has been interpreted in terms of an interrelation between the antiferromagnetic phase of the parent insulator and the underdoped regime of the intervening superconductor Zacher:99 .
In this paper, we argue that such extended structures in the superconducting OP, interpolating between point and line nodes, can be included in the definition of $`\mathrm{\Delta }_\stackrel{}{k}`$ as higher order terms in $`g(\stackrel{}{k})`$. We shall then look for their signatures in the low-temperature asymptotic electronic properties of the superconducting cuprates, as corrections to the predicted power-law behaviour. In deriving our results analytically, we will specifically consider the interlayer pair-tunneling (ILT) mechanism of high-$`T_c`$ superconductivity Chakravarty:93 , which has been shown to accurately reproduce most of the observed gap features Angilella:99 .
## 2 Extended $`d`$-wave gap within the ILT model
A distinguishing feature of the ILT mechanism, compared to other proposed models of HTSC, is that superconductivity is driven by a gain in kinetic, rather than potential, energy as temperature is lowered below the critical temperature $`T_c`$. It is assumed that coherent single particle hopping between adjacent CuO<sub>2</sub> layers in the cuprates is suppressed by the non-Fermi liquid character of the normal state (*e.g.* due to spin-charge separation), while interlayer coherent tunneling of Cooper pairs is allowed as soon as a superconducting condensate is established. Confined coherence Clarke:97 within CuO<sub>2</sub> layers in the normal state is indeed largely motivated by the absence of coherent transport along the $`c`$-axis, whereas a comprehensive theoretical understanding of it is still lacking. However, there is now abundant *experimental* evidence that $`c`$-axis transport in the normal state indeed is incoherent, while that in the superconducting state may not be vanderMarel . This seems to warrant attention being paid to unconventional models of high-$`T_c`$ superconductivity based on relieving $`c`$-axis frustrated kinetic energy. Recent findings Moler:98 ; Tsvetkov:98 suggest, however, that the ILT mechanism alone is not sufficient to account for the large condensation energy $`E_c`$, as extracted experimentally from measurements of the penetration length $`\lambda _c`$ of several single layered compounds, such as Tl<sub>2</sub>Ba<sub>2</sub>CuO<sub>6+δ</sub> Moler:98 ; Tsvetkov:98 , whereas the predictions of the ILT model agrees with the measured value of $`E_c`$ for La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> Panagopoulos:97 ; Panagopoulos:99 ; Leggett:98 ; Anderson:98 . It has been pointed out, however, that while considerable experimental effort has been devoted to the determination of $`\lambda _c`$, extracting $`E_c`$ from existing data on electronic specific heat is by no means straightforward Chakravarty:99 . A direct evaluation of $`E_c`$ from its mean-field expression at $`T=0`$ Schrieffer:64 would relieve the complications arising from thermal fluctuation effects, inherent in the method of integrating specific heat data, from $`T=0`$ through $`T_c`$, recently pointed out in Ref. Chakravarty:99 . By utilizing the gap equation, Eq. (4), within the ILT model, and of the expression relating $`\lambda _c`$ at $`T=0`$ to $`E_c`$ Chakravarty:98 , we find results for $`\lambda _c(T=0)`$ in Bi2212 which are within factors of order unity from the experimental values, rather than factors of 10 to 20 Angilella:00 . The observed doping dependence of $`\lambda _c`$ Panagopoulos:99 is also qualitatively reproduced Angilella:00 .
The emerging scenario suggests therefore that some in-plane effective interaction might co-operate with the ILT mechanism in establishing the superconducting state note:int-kin . One could think of such a mechanism as a seed for the Cooper instability, and the origin of the gap’s dominant $`d`$-wave symmetry. Once Cooper pairs are formed in the appropriate symmetry channel(s) via such in-plane effective interaction, the ILT mechanism would allow the condensate for an additional energy gain, by releasing the constraint of in-plane segregation.
Without explicitly specifying the microscopic origin of the in-plane mechanism, we therefore assume the in-plane pairing potential to be given by $`V_{\stackrel{}{k}\stackrel{}{k}^{}}=Vg(\stackrel{}{k})g(\stackrel{}{k}^{})`$ ($`V<0`$), thus allowing for $`d`$-wave symmetry of the order parameter. The issue of the competition with other subdominant ($`s`$-wave) symmetry channels in the presence of ILT has been addressed in Ref. Angilella:99 , showing that the $`d`$-wave contribution wins out at optimal doping and in the underdoped regime. Despite its kinetic nature, ILT can be absorbed in the interacting part of the Hamiltonian as an effective term $`T_J(\stackrel{}{k})\delta _{\stackrel{}{k}\stackrel{}{k}^{}}`$, whose $`\stackrel{}{k}`$-space locality enforces in-plane momentum conservation during a tunneling process Chakravarty:93 . Following Ref. Chakravarty:93 , we assume $`T_J(\stackrel{}{k})=t_{}^2(\stackrel{}{k})/tT_Jg^4(\stackrel{}{k})`$, being $`t_{}(\stackrel{}{k})`$ the single-particle interlayer hopping amplitude, with $`t_{}(\stackrel{}{k})g^2(\stackrel{}{k})`$, as suggested by ARPES as well as by band structure calculations Chakravarty:93 ; Andersen:96 . A standard mean-field diagonalization technique then yields the following expression for the energy gap Sudbo:95 ; Angilella:99 :
$$\mathrm{\Delta }_\stackrel{}{k}=\frac{\mathrm{\Delta }g(\stackrel{}{k})}{1T_J(\stackrel{}{k})\chi _\stackrel{}{k}},$$
(4)
where $`\chi _\stackrel{}{k}=(2E_\stackrel{}{k})^1\mathrm{tanh}(\beta E_\stackrel{}{k}/2)`$ is the superconducting pair susceptibility, and $`E_\stackrel{}{k}=(\xi _\stackrel{}{k}^2+|\mathrm{\Delta }_\stackrel{}{k}|^2)^{1/2}`$ is the upper branch of the superconducting elementary excitation spectrum.
Along the Fermi line ($`\xi _\stackrel{}{k}=0`$) at $`T=0`$, one immediately sees that:
$$\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }g(\stackrel{}{k})+\frac{1}{2}T_Jg^4(\stackrel{}{k})sgn[g(\stackrel{}{k})].$$
(5)
Such an expression, together with manifestly fulfilling the requirement of $`d`$-wave symmetry, also endows the superconducting gap with a richer structure near the nodal points along the $`k_x=\pm k_y`$ directions. This is probably best seen by considering the Fourier expansions:
$`g(\stackrel{}{k})`$ $`=2{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}J_{4m2}(k)\mathrm{cos}[(4m2)\varphi ],`$ (6a)
$`h(\stackrel{}{k})`$ $`=J_0(k)+2{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}J_{4m}(k)\mathrm{cos}(4m\varphi ),`$ (6b)
$`T_J(\stackrel{}{k})/T_J`$ $`={\displaystyle \frac{9}{64}}+{\displaystyle \frac{a_0}{2}}+{\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}a_{4m}\mathrm{cos}(4m\varphi ),`$ (6c)
with
$`a_{4m}`$ $`={\displaystyle \frac{1}{32}}J_{4m}(4k)+{\displaystyle \frac{1}{2}}J_{4m}(2k)+{\displaystyle \frac{3}{16}}(1)^mJ_{4m}(2k\sqrt{2})`$
$`{\displaystyle \frac{3}{4}}(1)^mJ_{4m}(k\sqrt{2}){\displaystyle \frac{1}{4}}J_{4m}\left({\displaystyle \frac{k}{\mathrm{sin}\varphi _0}}\right)\mathrm{cos}(4m\varphi _0).`$
Here, the generic wavevector $`\stackrel{}{k}`$ is expressed in terms of its modulus $`k`$ and of the angle $`\varphi `$ formed with the $`\mathrm{\Gamma }X`$ direction in the first Brillouin zone (1BZ), $`\stackrel{}{k}=(k\mathrm{cos}\varphi ,k\mathrm{sin}\varphi )`$, $`J_\alpha (x)`$ are Bessel functions of the first kind and order $`\alpha `$, and $`\mathrm{tan}\varphi _0=\frac{1}{3}`$.
Eq. (5) is to be contrasted to the phenomenological fit
$$\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }[B\mathrm{cos}(2\varphi )+(1B)\mathrm{cos}(6\varphi )]$$
(8)
proposed in Ref. Mesot:99 for $`\mathrm{\Delta }_\stackrel{}{k}`$ along the Fermi line: Instead of requiring an in-plane interaction extended to further neighbours, Eq. (4) endows the superconducting gap with the observed flat structure around the nodes, through the ILT term $`T_J(\stackrel{}{k})`$. In Fig. 1, we fit Eq. (5) against Mesot et al.’s experimental data for one of the underdoped Bi2212 samples in Ref. Mesot:99 , having $`T_c=75`$ K. A remarkable agreement follows already by fixing $`\mathrm{\Delta }`$ so that $`|\mathrm{\Delta }_𝐤|`$ reproduces the maximum datum at $`\stackrel{}{k}=(\pi ,\pi )`$, whereas $`T_J`$ is taken to be 0.04 eV Chakravarty:93 . In particular, besides obtaining an enhanced maximum value of $`|\mathrm{\Delta }_𝐤|`$ at $`\stackrel{}{k}=(\pi ,\pi )`$, we are thus able to recover the anomalously flat region around the node at $`\varphi =45^{}`$ in a rather natural way. We note, however, that our fit requires $`\mathrm{\Delta }T_J/2`$ around optimal doping, which will not be without consequences in evaluating other fundamental quantities Angilella:00 .
Eq. (5) already contains the doping dependence of the observed gap anisotropy, although in a hidden way. As pointed out in Ref. Angilella:99 , the auxiliary parameter $`\mathrm{\Delta }`$ is to be self-consistently determined by solving the appropriate gap equation. Besides being intrinsically doping dependent, this equation is unconventionally modified by the presence of a $`\stackrel{}{k}`$-local effective interaction, as induced by the ILT mechanism. Moreover, the role of the contribution $`T_Jg^4(\stackrel{}{k})`$ in Eq. (5) is strongly influenced by the actual location of the Fermi line, as $`g^4(\stackrel{}{k})`$ is sharply peaked at $`\stackrel{}{k}=(0,\pi )`$ (and symmetry related points).
Eq. (5) also facilitates the evaluation of the slope of the superconducting gap $`v_\mathrm{\Delta }=(1/2)d|\mathrm{\Delta }_\stackrel{}{k}|/d\varphi `$ at the nodal point along the Fermi line. Such a quantity is related to the temperature derivative of the superfluid stiffness at $`T=0`$. In particular, it is seen that the ratio $`v_\mathrm{\Delta }/\mathrm{\Delta }_{\mathrm{max}}`$ decreases with underdoping Mesot:99 . From Eq. (5), one derives that $`v_\mathrm{\Delta }`$ is independent of $`T_J`$, and that therefore a doping induced change of $`v_\mathrm{\Delta }`$ through $`\mathrm{\Delta }`$ essentially can be traced back to the actual position of the Fermi line, as discussed above, within the ILT model. The ratio $`v_\mathrm{\Delta }/\mathrm{\Delta }_{\mathrm{max}}`$ will anyway deviate from its value within simple $`d`$-wave (BCS-like) models, as a function of doping, due to the enhancement of $`\mathrm{\Delta }_{\mathrm{max}}`$ induced by ILT.
## 3 Low-temperature asymptotic behaviour of electronic properties
We now address the issue, whether such extended features of the OP near the nodes, as those described in the previous section, induce deviations in the low or intermediate temperature asymptotic behaviour of linear response electronic properties in the superconducting state. In what follows, we shall limit our discussion to clean superconductors, and neglect impurity effects altogether. Mean-field (BCS or BCS-like) expressions for most linear response electronic properties are available also in the case of anisotropic, *i.e.* non $`s`$-wave, superconductors. In particular, we have in mind observable quantities such as the superconducting density Leggett:75 , the electronic specific heat Leggett:75 , the spin susceptibility Leggett:75 , the penetration depth Scalapino:92 , the thermal conductivity Bardeen:59 , and so on. Their expressions basically involve the evaluation of some integral of the kind:
$$[\beta ;\phi _\stackrel{}{k}(\beta )]=\frac{1}{(2\pi )^2}d^2\stackrel{}{k}\phi _\stackrel{}{k}(\beta )e^{\beta E_\stackrel{}{k}},$$
(9)
where $`\beta =(k_\mathrm{B}T)^1`$, $`\phi _\stackrel{}{k}(\beta )`$ is a (dimensional) function of wavevector $`\stackrel{}{k}`$ and temperature, related to the electronic quantity of interest, and the integration is extended to the 1BZ, $`\stackrel{}{k}[\pi ,\pi ]\times [\pi ,\pi ]`$ (see App. A). In the case of $`d`$-wave superconductors, $`E_\stackrel{}{k}`$ is allowed to vanish at the intersection between the Fermi line and the nodal lines of the gap function. Around such points, quasiparticles can be created in large numbers. In the limit of low temperatures ($`\beta \mathrm{}`$), therefore, the value of the integral in Eq. (9) is dominated by the contributions from wavevectors $`\stackrel{}{k}`$ close to such point nodes. Around such nodes, it is useful to introduce the new sets of coordinates $`(k_1,k_2)`$ or $`(ϵ,\theta )`$, defined as Lee:97 :
$`\xi _\stackrel{}{k}`$ $`𝐯_\mathrm{F}\stackrel{}{k}v_\mathrm{F}k_1=ϵ\mathrm{cos}\theta ,`$ (10a)
$`\mathrm{\Delta }g(\stackrel{}{k})`$ $`𝐯_2\stackrel{}{k}v_2k_2=ϵ\mathrm{sin}\theta ,`$ (10b)
in units where $`\mathrm{}=1`$. Here, $`v_\mathrm{F}`$ and $`v_2`$ are the Fermi velocity and a suitable ‘gap’ velocity, respectively, evaluated at $`E_\stackrel{}{k}=0`$, and $`ϵ`$ measures the distance in energy from a given dispersionless point implicitly defined by $`E_\stackrel{}{k}=0`$. In terms of the new coordinates, the superconducting spectrum for a simple $`d`$-wave superconductor near a node therefore looks like an anisotropic Dirac cone Lee:97 ,
$$E_\stackrel{}{k}\left(v_\mathrm{F}^2k_1^2+v_2^2k_2^2\right)^{1/2}=ϵ.$$
(11)
The observation of flatter structures near the nodes Mesot:99 not only implies a more significant anisotropy ratio $`v_\mathrm{F}/v_2`$, but also the possibility that higher order terms in $`ϵ`$ may contribute to $`E_\stackrel{}{k}`$, Eq. (11). Indeed, within the ILT model, from Eq. (4) at $`T=0`$ one obtains
$$E_\stackrel{}{k}ϵ\left[1+\left(\frac{ϵ}{ϵ_{}}\right)^3\mathrm{sin}^6\theta \right],$$
(12)
to lowest order in $`ϵ/ϵ_{}`$, with $`1/ϵ_{}^3=(1/2)(T_J/\mathrm{\Delta }^4)`$ related to the pair-tunneling amplitude $`T_J`$ and to the auxiliary gap parameter $`\mathrm{\Delta }`$ (Figs. 2 and 3).
Other models, based on extended in-plane pairing mechanisms, would in general yield different polynomial corrections in $`ϵ`$ to $`E_\stackrel{}{k}`$. For instance, within the spin fluctuation theory Monthoux:91 , the following phenomenological expansion holds for the momentum distribution of the superconducting energy gap Ghosh:99
$$\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }g(\stackrel{}{k})\underset{n=0}{\overset{N}{}}d_nh^n(\stackrel{}{k}),$$
(13)
with all coefficients $`d_n=1`$. We explicitly observe that for $`N=0`$, $`d_0=1`$, one recovers the simple $`d`$-wave gap $`\mathrm{\Delta }_\stackrel{}{k}g(\stackrel{}{k})`$, while the case $`N=1`$, with the identifications $`\mathrm{\Delta }B\mathrm{\Delta }`$, $`d_0=1`$, $`d_1=4(1B)/B`$, maps to Mesot:99 ; Ghosh:99 :
$$\mathrm{\Delta }_\stackrel{}{k}=\mathrm{\Delta }[Bg(\stackrel{}{k})+(1B)g(2\stackrel{}{k})],$$
(14)
which is compatible with the phenomenological fit Eq. (8) proposed by Mesot *et al.* in Ref. Mesot:99 for their experimental data of $`|\mathrm{\Delta }_\stackrel{}{k}|`$ along the Fermi line Mesot:99 ; Ghosh:99 . In particular, Eq. (14) would follow from a correction $`\delta V_{\stackrel{}{k}\stackrel{}{k}^{}}g(2\stackrel{}{k})g(2\stackrel{}{k}^{})`$ to the in-plane coupling, corresponding to next-nearest neighbours interaction.
In such a particular case, and assuming for simplicity $`t^{}=\mu =0`$ in Eq. (3), one straightforwadly obtains
$$E_\stackrel{}{k}ϵ\left[\mathrm{cos}^2\theta +\mathrm{sin}^2\theta \left(1\frac{ϵ}{\stackrel{~}{ϵ}_{}}\mathrm{cos}\theta \right)^2\right]^{1/2},$$
(15)
where $`\stackrel{~}{ϵ}_{}=tB/(1B)`$ is now related to the ratio of nearest *vs* next-nearest neighbours coupling. Therefore, both within the ILT model and within other models, based on extended in-plane pairing, the additional mechanism responsible for the nonlinear correction to $`E_\stackrel{}{k}`$ away from its nodes introduces new energy scales (here, $`ϵ_{}`$ or $`\stackrel{~}{ϵ}_{}`$, respectively). Fig. 3 depicts the two different ways in which $`E_\stackrel{}{k}`$ deviates from the cone-like shape, Eq. (11), near a node, in the two cases given by Eqs. (12) and (15).
In the absence of any such additional mechanism ($`ϵ_{},\stackrel{~}{ϵ}_{}=0`$), the leading contribution to Eq. (9) for the simplest, reference case $`\phi _\stackrel{}{k}(\beta )1`$ is:
$$_1(\beta )[\beta ;1]\frac{A}{\beta ^2},$$
(16)
where $`A=(2\pi v_\mathrm{F}v_2\mathrm{\Delta })^1`$ is a doping-dependent factor, and $``$ denotes equality up to terms vanishing exponentially with $`\beta `$ at *all* energy scales, as $`\beta \mathrm{}`$ ($`T0`$). Eq. (16) should be regarded as typical of the power-law asymptotic low-temperature behaviour of the superconducting electronic properties within a simple $`d`$-wave BCS-like model.
In order to obtain an asymptotic expansion for $`_1(\beta )`$ as $`\beta \mathrm{}`$ ($`T0`$), including the corrections due to ILT, Eq. (12), we observe that the integration over $`ϵ`$ in Eq. (9) is actually made of two contributions:
$$_0^{\mathrm{}}𝑑ϵ=_0^ϵ_{}𝑑ϵ+_ϵ_{}^{\mathrm{}}𝑑ϵ.$$
(17)
In the first integral, we may safely retain only the linear term $`E_\stackrel{}{k}ϵ`$ in the exponent, since $`ϵϵ_{}`$. In the second contribution, this is no longer possible, and Eq. (12) has to be retained in full. However, since $`ϵϵ_{}>0`$, one can make use of Laplace’s (saddle point) method for the integral over angles around $`\theta =0`$. The final result is:
$`_1(\beta )`$ $`{\displaystyle \frac{A}{\beta ^2}}\left[1(1+\beta ϵ_{})e^{\beta ϵ_{}}+{\displaystyle \frac{1}{3\pi }}\mathrm{\Gamma }\left({\displaystyle \frac{1}{6}}\right)\beta ϵ_{}\mathrm{\Gamma }({\displaystyle \frac{4}{3}},\beta ϵ_{})\right]`$ (18)
$`{\displaystyle \frac{A}{\beta ^2}}\left[1\left(1+\beta ϵ_{}{\displaystyle \frac{1}{3\pi }}\mathrm{\Gamma }\left({\displaystyle \frac{1}{6}}\right)(\beta ϵ_{})^{5/6}\right)e^{\beta ϵ_{}}\right],`$
where $`\mathrm{\Gamma }(x)`$, $`\mathrm{\Gamma }(\alpha ,x)`$ are Euler gamma and incomplete gamma functions, respectively Gradshteyn:94 . A comparison of Eq. (16) and Eq. (18) is provided by Fig. 4, and shows that $`_1(\beta )`$ gets effectively *suppressed* in the presence of flat nodes in the order parameter, as provided by the ILT mechanism, with respect to the simple $`d`$-wave case, at an energy scale $`ϵ_{}`$.
No such simple asymptotic expansion for $`_1(\beta )`$ is available in the extended $`d`$-wave case described by Eq. (15), and the integrations have to be performed numerically. Fig. 4 shows the result, with the identifications $`A\stackrel{~}{A}=(2\pi v_\mathrm{F}v_2B\mathrm{\Delta })^1`$ and $`ϵ_{}\stackrel{~}{ϵ}_{}`$. In this case, Eq. (15) provides $`E_\stackrel{}{k}`$ with a different kind of anisotropy with respect to the simple case, Eq. (11), than Eq. (12) does. While in the latter case one always has $`E_\stackrel{}{k}ϵ`$, here one has $`E_\stackrel{}{k}ϵ`$, depending on the angle $`\theta `$ (cfr. Fig. 3). As a consequence, $`_1(\beta )`$ is *enhanced* with respect to the simple $`d`$-wave case, at an energy scale $`\stackrel{~}{ϵ}_{}`$.
## 4 Conclusions
Motivated by recent experimental findings of extended flat structures in the order parameters of the underdoped $`d`$-wave superconductor Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> Mesot:99 , we have addressed the issue of whether nonlinear, high-energy corrections to the superconducting energy spectrum $`E_\stackrel{}{k}`$ around the gap nodes induce deviations in the predicted power-law behaviour of several electronic properties at low or intermediate temperatures. We have shown that nonlinear corrections to $`E_\stackrel{}{k}`$ in general introduce additional energy scales in the problem. Deviations from the usual power-law behaviour of the superconducting electronic properties are indeed to be expected at such energy scales, but the actual value and *sign* of such deviations are specific to the model under consideration. In particular, within the ILT model, we have explicitly derived the expected corrections to a typical power-law asymptotic behaviour as $`T0`$, showing these to be *negative,* whereas within a phenomenological model of extended $`d`$-wave superconductivity Mesot:99 such corrections are predicted to be *positive.* Whether such deviations will actually be observable in real measurements of superconducting electronic properties, will of course depend on the effective values of the additional energy scales $`ϵ_{}`$ or $`\stackrel{~}{ϵ}_{}`$ in real compounds.
###### Acknowledgements.
G. G. N. A. thanks P. Falsaperla, J. O. Fjærestad, and Ch. Wälti for valuable discussions, and acknowledges the NTNU (Trondheim, Norway) for warm hospitality and financial support during the period in which the present work was brought to completion. A. S. acknowledges support from Norges Forskningsråd through Grants No. 110566/410 and No. 110569/410.
## Appendix A Low-temperature superconducting electronic properties
We now give a sketch of how the low-temperature asymptotic behaviour of several electronic properties of interest can be reduced to that of $`_1(\beta )`$ or its derivatives. Most electronic quantities in the superconducting state are in fact given by Eq. (9), with $`\phi _\stackrel{}{k}(\beta )`$ actually depending on $`\stackrel{}{k}`$ only through $`E_\stackrel{}{k}`$. In what follows, $`f(ϵ)=[1+\mathrm{exp}(\beta ϵ)]^1`$ denotes the Fermi function.
Within BCS theory, the electronic specific heat is given by Leggett:75 :
$`C_V`$ $`={\displaystyle \underset{\stackrel{}{k}}{}}2k_\mathrm{B}\beta E_\stackrel{}{k}\left[E_\stackrel{}{k}+{\displaystyle \frac{E_\stackrel{}{k}}{\beta }}\right]\left({\displaystyle \frac{f}{E_\stackrel{}{k}}}\right)`$ (19)
$`{\displaystyle \underset{\stackrel{}{k}}{}}2k_\mathrm{B}\beta E_\stackrel{}{k}^2\left({\displaystyle \frac{f}{E_\stackrel{}{k}}}\right)`$
$`2k_\mathrm{B}\beta ^2_1^{\prime \prime }(\beta ),`$
where apices denote derivatives with respect to $`\beta `$. Here, $`\phi _\stackrel{}{k}(\beta )=2k_\mathrm{B}\beta ^2E_\stackrel{}{k}^2`$, and we have made use of the fact that $`(f/E_\stackrel{}{k})\beta \mathrm{exp}(\beta E_\stackrel{}{k})`$.
Analogously, the unrenormalized, static, isotropic spin susceptibility $`\chi _0=\chi _0(𝐪0,\omega 0)`$, which is directly related to the Knight shift, is simply given by Schrieffer:64 ; Sudbo:94 :
$$\chi _0=\underset{\stackrel{}{k}}{}\left(\frac{f}{E_\stackrel{}{k}}\right)\beta _1(\beta ).$$
(20)
The expression of the electronic thermal conductivity for an anisotropic $`d`$-wave superconductor also involves an average of $`(f/E_\stackrel{}{k})`$ over the 1BZ Bardeen:59 ; Krishana:97 :
$$\kappa _e=\frac{1}{T}\underset{\stackrel{}{k}}{}\left(\frac{f}{E_\stackrel{}{k}}\right)E_\stackrel{}{k}^2\left(\frac{E_\stackrel{}{k}}{k_x}\right)^2\tau (\stackrel{}{k}),$$
(21)
where $`\tau (\stackrel{}{k})`$ is the superconducting quasiparticles lifetime. Due to the presence of the $`x`$ component of the group velocity $`_\stackrel{}{k}E_\stackrel{}{k}`$, however, its expression in our notation reduces to:
$$\kappa _e=\frac{1}{8\pi ^2}\frac{\mathrm{}_0}{v_2}\frac{1}{T}_0^{\mathrm{}}𝑑ϵϵ^3_0^{2\pi }𝑑\theta \left(\frac{f}{ϵ}\right)\left(\mathrm{cos}\theta +\frac{v_2}{v_\mathrm{F}}\mathrm{sin}\theta \right)^2,$$
(22)
where $`\mathrm{}_0=v_\mathrm{F}\tau (\stackrel{}{k}_\mathrm{F})`$ is the quasiparticle mean free path at the nodes. The final result crucially depends on the anisotropy ratio $`v_2/v_\mathrm{F}`$, and would be different in the two cases given by Eqs. (12) and (15), due to their different $`\theta `$ dependence. This has to be contrasted with the result obtained in the simple $`d`$-wave case, where Krishana:97 :
$$\kappa _e=\eta k_\mathrm{B}^3T^2\frac{\mathrm{}_0}{v_2}\left(1+\frac{v_2^2}{v_\mathrm{F}^2}\right),$$
(23)
with $`\eta =(8\pi )^1_0^{\mathrm{}}𝑑xx^3(f/x)`$.
## Appendix B A limiting case
In the absence of in-plane coupling, a spurious solution of the mean-field gap equation at $`T=0`$ can be implicitly expressed via Angilella:99 :
$$E_\stackrel{}{k}=\frac{1}{2}T_J(\stackrel{}{k})=\frac{1}{2}T_Jg^4(\stackrel{}{k}).$$
(24)
In such a limiting case, the superconducting energy spectrum would have purely kinetic origin, and would be identified with the interlayer pair-tunneling amplitude, divided by two. A closed expression can then be obtained for $`_1(\beta )`$, by utilizing the useful result:
$$\frac{1}{(2\pi )^2}d^2\stackrel{}{k}𝒢[\eta (\stackrel{}{k})]=\frac{2}{\pi ^2}_1^1𝑑x𝒢(x)K(\sqrt{1x^2}),$$
(25)
where $`𝒢[\eta (x)]`$ is any continuous functional of $`\eta (\stackrel{}{k})=h(\stackrel{}{k})`$ or $`g(\stackrel{}{k})`$ alone, and $`K(x^{})`$ ($`x^{}=\sqrt{1x^2}`$) is the complete elliptic integral of first kind Gradshteyn:94 . From Eq. (9), expanding $`K(x^{})`$ around $`x=0`$, one eventually arrives at the closed expression:
$$_1(\beta )\frac{1}{\pi ^2}\mathrm{\Gamma }\left(\frac{1}{4}\right)\frac{1}{\zeta ^{1/4}}\left(2\mathrm{log}2\frac{1}{4}\psi \left(\frac{1}{4}\right)+\mathrm{log}\zeta ^{1/4}\right),$$
(26)
where $`\psi (x)`$ is the digamma function Gradshteyn:94 , and the ILT amplitude $`T_J`$ itself here fixes the appropriate energy scale, through $`\zeta =\frac{1}{2}\beta T_J`$. |
warning/0001/astro-ph0001218.html | ar5iv | text | # Gravitational waves from inspiralling compact binaries with magnetic dipole moments
## 1 INTRODUCTION
Direct detection of gravitational waves (GWs) is one of the most exciting challenges in the history of science. Long baseline interferometers for detection of GWs such as TAMA300 (Kuroda et al. 1997), GEO600 (Hough 1992), VIRGO (Bradaschia et al. 1990), and LIGO (Abramovici et al. 1992) will be in operation within five years. One of the most promising sources of GWs for such detectors is the inspiralling binary neutron stars (BNSs), since we may expect several coalescing events per year within 200Mpc (Phinney 1991, Narayan, Piran, & Shemi 1991, van den Heuvel & Lorimer 1996).
As the orbital radius of BNSs decays due to gravitational radiation reaction, the frequency of GWs sweeps upward in detector’s sensitive bandwidth from $`10`$Hz to $`1000`$Hz. In the early inspiralling phase of BNSs, each neutron star (NS) can be treated as a point particle and the post-Newtonian (PN) expansion will converge (Cutler et al. 1993) since the orbital separation of NSs is much larger than the NS’s radius and the orbital velocity is much smaller than the velocity of light. This means that the theoretical templates of GWs in the inspiralling phase can be calculated with high accuracy by the PN approximation of general relativity using only several parameters: each mass and spin of NSs and the initial orbital elements. By cross correlating the observed noisy signals with the theoretical templates, the binary parameters, such as masses, spins (Cutler et al. 1993, Kidder, Will, & Wiseman 1993, Cutler & Flanagan 1994, Poisson & Will 1995), and cosmological distances (Schutz 1986), are determined. The number of rotation of BNSs is about 16000 in the detector’s sensitive bandwidth. Therefore, the quite accurate theoretical templates are needed in order to extract physical information about BNSs from GWs since any effect that causes only one cycle ambiguity over 16000 accumulated cycles in the theoretical templates will reduce the signal-to-noise ratio (SNR). For the inspiralling BNSs in the sensitive bandwidth, with $`v^2m/r`$ (hereafter $`G=c=1`$) typically around $`10^2`$, the correction of $`1/1600010^4`$ corresponds to second PN (2PN) order, $`(m/r)^210^4`$ (Blanchet et al. 1995). Many efforts are devoted to calculate higher-order PN corrections to theoretical templates (e.g., Blanchet 1996, Jaranowski & Sch$`\ddot{\mathrm{a}}`$fer 1998a, 1998b, Damour, Jaranowski & Sch$`\ddot{\mathrm{a}}`$fer 1999, Tagoshi & Nakamura 1994, Tagoshi & Sasaki 1994, Poisson 1995). However, these studies pay attention only to the gravitational effects on the theoretical templates, and it has not been studied how large corrections to the theoretical templates are caused by the electromagnetic effects, i.e., the magnetic fields of NSs, as far as we know.<sup>1</sup><sup>1</sup>1 Of course, the electromagnetic effects on the GWs from a rotating NS are investigated vigorously (Bocquet et al. 1995, Bonazzola & Gourgoulhon 1996, Konno, Obata, & Kojima 1999).
NSs observed as radio pulsars are believed to have strong magnetic fields, typically $`10^{12}`$G, assuming that the spin-down of pulsars is due to magnetic dipole radiation (e.g. Taylor, Manchester, & Lyne 1993). In addition, it begins to be recognized recently that NSs with superstrong magnetic fields $`10^{14}`$G really exist (see below). We can crudely estimate the correction to the waveform due to magnetic dipole fields of NSs by comparing the gravitational force $`F_Gm_1m_2/r^2`$ and the magnetic force $`F_M3\mu _1\mu _2/r^4`$ between BNSs of masses, $`m_1`$ and $`m_2`$, with magnetic dipole moments, $`\mu _1`$ and $`\mu _2`$, as
$`{\displaystyle \frac{F_M}{F_G}}`$ $`=`$ $`1\times 10^4\left({\displaystyle \frac{r}{6(m_1+m_2)}}\right)^2\left({\displaystyle \frac{H_1}{2\times 10^{16}\mathrm{G}}}\right)\left({\displaystyle \frac{H_2}{2\times 10^{16}\mathrm{G}}}\right)`$ (1)
$`\times \left({\displaystyle \frac{R_1}{10^6\mathrm{cm}}}\right)^3\left({\displaystyle \frac{R_2}{10^6\mathrm{cm}}}\right)^3\left({\displaystyle \frac{1.4M_{}}{m_1}}\right)\left({\displaystyle \frac{1.4M_{}}{m_2}}\right)\left({\displaystyle \frac{2.8M_{}}{m_1+m_2}}\right)^2,`$
where $`H_p=2\mu _p/R_p^3`$ $`(p=1,2)`$ are the magnetic fields at the pole of the star and $`R_p`$ are radii of NSs. This corresponds to the 2PN order correction. Therefore, the magnetic fields of order $`10^{16}`$G might cause about one rotation error. Note that the $`r`$-dependence of the magnetic correction is the same as that of 2PN order, $`v^4r^2`$, so that this argument is independent of the value of the separation $`r`$.
Theoretically, in a new born NS, such superstrong magnetic fields $`10^{16}`$ G can be generated if the initial spin is in millisecond range since the conditions for helical dynamo action are met during the first few seconds after gravitational collapse (Duncan & Thompson 1992, Thompson & Duncan 1993). Observationally, such a superstrongly magnetized NS, or “magnetar”, may be found as the soft gamma repeaters (SGRs) and the anomalous X-ray pulsars (AXPs) (Duncan & Thompson 1992, Thompson & Duncan 1996). The dipolar magnetic fields of SGRs are estimated as $`10^{15}`$G by using the measured spin periods with spin-down rates for SGR1900+14 (Kouveliotou et al. 1999, Woods et al. 1999b) and SGR1806-20 (Kouveliotou et al. 1998), and with the peak luminosity for SGR1627-41 (Woods et al. 1999a). APXs that have measured the spin-down rate (Israel et al. 1999, Mereghetti, Israel, & Stella 1998, Baykal et al. 1998, Vasisht & Gotthelf 1997, Haberl et al. 1997) can be considered to have magnetic fields of $`10^{14}`$$`10^{15}`$G. Although there are some other models for SGRs and AXPs, these objects are best understood within the framework of magnetar (see, e.g., Thompson & Duncan 1995, Kouveliotou et al. 1998, Vasisht & Gotthelf 1997). Therefore, from both theoretical and observational results, it may be possible that a NS has superstrong magnetic fields $`10^{16}`$G.
In this paper we investigate the effects of the magnetic dipole fields of NSs on the frequency evolution of GWs from the inspiralling BNSs, since the magnetic fields of order $`10^{16}`$G might cause about one rotation error. Throughout the paper, we use the relation $`H=2\mu /R^3`$ to connect the magnetic moment $`\mu `$ to the magnetic field at the magnetic pole $`H`$, and $`R=10^6`$cm as the radius of a NS. For later convenience, note that $`\mu =1.4\times 10^9(H/10^{16}\mathrm{G})`$ cm<sup>2</sup> in units of $`G=c=1`$.
## 2 EQUATIONS OF MOTION
We consider a binary system of two compact bodies of masses, $`m_1`$ and $`m_2`$, with magnetic dipole moments, $`𝝁_1`$ and $`𝝁_2`$, respectively. Since we pay particular attention to the effects of magnetic fields, we treat the orbital motion of BNSs in Newtonian gravity. Although a spherical symmetry of the stellar configuration is, in general, incompatible with the presence of magnetic fields (Chandrasekhar 1981), we ignore quadrupole effects which are caused by magnetic fields for a moment (see §5 and §B). We also neglect tidal effects which are expected to be small until pre-merging phase of BNSs (Bildsten and Cutler 1992).
By eliminating the motion of the center of mass of BNSs and setting the origin of the coordinate frame at the center of mass of BNSs, the effective one-body equations of motion can be derived from a Lagrangian,
$`={\displaystyle \frac{1}{2}}\eta mv^2+{\displaystyle \frac{\eta m^2}{r}}+_{DD},`$ (2)
where
$`_{DD}=𝝁_1𝑯_2={\displaystyle \frac{1}{r^3}}\left\{3(\widehat{𝒏}𝝁_1)(\widehat{𝒏}𝝁_2)𝝁_1𝝁_2\right\}.`$ (3)
Here $`m=m_1+m_2`$, $`\eta =m_1m_2/m^2`$, $`r=|𝒙|`$, $`𝒙=𝒙_1𝒙_2`$, $`\widehat{𝒏}=𝒙/r`$, $`𝒗=\dot{𝒙}`$, and $`𝑯_2=\left\{3(\widehat{𝒏}𝝁_2)\widehat{𝒏}𝝁_2\right\}/r^3`$ is the magnetic field at $`𝒙_1`$ produced by the magnetic moment $`𝝁_2`$. By using the Euler-Lagrange equations, we obtain the equations of motion as
$$𝒂=\frac{m}{r^2}\widehat{𝒏}+𝒂_{DD},$$
(4)
where $`𝒂=\ddot{𝒙}`$ and
$$𝒂_{DD}=\frac{3}{\eta mr^4}\left\{(𝝁_1𝝁_2)\widehat{𝒏}+(\widehat{𝒏}𝝁_2)𝝁_1+(\widehat{𝒏}𝝁_1)𝝁_25(\widehat{𝒏}𝝁_1)(\widehat{𝒏}𝝁_2)\widehat{𝒏}\right\}.$$
(5)
From equation (2), the energy of this system is given by
$$E=\frac{1}{2}\eta mv^2\frac{\eta m^2}{r}+E_{DD},$$
(6)
where
$$E_{DD}=\frac{1}{r^3}\left\{3(\widehat{𝒏}𝝁_1)(\widehat{𝒏}𝝁_2)𝝁_1𝝁_2\right\}.$$
(7)
The total angular momentum can be defined as $`𝑳=𝑳_N+𝑺`$ where $`𝑳_N=\eta m(𝒙\mathbf{\times }𝒗)`$ is the Newtonian orbital angular momentum and $`𝑺=𝑺_1+𝑺_2`$ is the total spin angular momentum. We can show explicitly $`\dot{E}=\dot{𝑳}=0`$ with the equations of motion (4) and the evolution equations of the spins,<sup>2</sup><sup>2</sup>2 Note that there are no spin-orbit and spin-spin interactions since we consider Newtonian gravity. The spins are generated by the torque due to the magnetic dipole-dipole interaction.
$$\dot{𝑺}_p=𝝁_p\mathbf{\times }𝑯_q=\frac{1}{r^3}\{3(\widehat{𝒏}𝝁_q)(𝝁_p\mathbf{\times }\widehat{𝒏})𝝁_p\mathbf{\times }𝝁_q\}.(p,q=1,2)$$
(8)
Assuming NSs as spherical compact bodies, the spin angular velocities $`𝛀_p`$ are related to the spins $`𝑺_p`$ as $`𝑺_p=I_p𝛀_p`$ where $`I_p`$ is the principle moment of inertia of the bodies. Since the magnetic moments evolve as $`\dot{𝝁_p}=𝛀_p\mathbf{\times }𝝁_p=𝑺_p\mathbf{\times }𝝁_p/I_p`$, the angular velocities of the magnetic moments $`\mathrm{\Omega }_p`$ will be of order $`\mathrm{\Omega }_p(\mu _1\mu _2/mR^2r^3)^{1/2}`$ from dimensional analysis with equation (8). Note also that the orbital angular velocity $`w`$ is of order $`w(m/r)^{3/2}/m`$, and the orbital inspiral rate $`w_{ins}=(dE/dt)_{GW}/E`$ is of order $`w_{ins}(m/r)^4/m`$.
## 3 GRAVITATIONAL WAVES AND ELECTROMAGNETIC WAVES
As a first step, we use the quadrupole formula to derive the rate of energy loss from a binary system due to GWs (e.g. Thorne 1980). The symmetric, trace-free parts of the quadrupole moments of this system are given by $`I_{ij}=\eta m(x_ix_j\frac{1}{3}r^2\delta _{ij})`$. Taking third time derivatives of these quadrupole moments, we obtain the energy loss rate from the quadrupole formula as
$`\left({\displaystyle \frac{dE}{dt}}\right)_{GW}`$ $`=`$ $`{\displaystyle \frac{1}{5}}\stackrel{}{I}_{ij}\stackrel{}{I}_{ij}`$ (9)
$`=`$ $`{\displaystyle \frac{8}{15}}{\displaystyle \frac{\eta ^2m^4}{r^4}}\{12v^211\dot{r}^2`$
$`+{\displaystyle \frac{1}{\eta m^2r^2}}[6(12v^2+13\dot{r}^2)(𝝁_1𝝁_2)+12(21v^234\dot{r}^2)(\widehat{𝒏}𝝁_1)(\widehat{𝒏}𝝁_2)`$
$`36(𝒗𝝁_1)(𝒗𝝁_2)+87\dot{r}\{(\widehat{𝒏}𝝁_1)(𝒗𝝁_2)+(𝒗𝝁_1)(\widehat{𝒏}𝝁_2)\}]\},`$
where we have used the equations of motion (4) and assumed $`\mu _1\mu _2/m^2r^21`$.
On the other hand, electromagnetic (EM) waves are also emitted from this binary system since the magnetic moments are moving. Using the linearity in the EM fields, the radiation fields $`𝑩_0^{rad}`$ in equation (A4) for this binary system are given by
$$𝑩_0^{rad}=\frac{1}{D}(\widehat{𝒅}𝝁_{eff})\left\{(\widehat{𝒅}\dot{𝒂})\widehat{𝒅}\dot{𝒂}\right\},$$
(10)
where
$$𝝁_{eff}=\frac{1}{m}(m_2𝝁_1m_1𝝁_2).$$
(11)
Therefore, since the radiated power is given by equation (A6), the rate of energy loss due to the EM radiation is calculated as
$$\left(\frac{dE}{dt}\right)_{EM}=\frac{2}{15}\frac{m^2}{r^6}\left[2\mu _{eff}^2\left\{v^26\dot{r}(\widehat{𝒏}𝒗)+9\dot{r}^2\right\}\left\{𝝁_{eff}(𝒗3\dot{r}\widehat{𝒏})\right\}^2\right],$$
(12)
where we have substituted $`𝝁_{eff}`$ into $`𝝁`$ in equation (A6).
Note that the assumption of the constant magnetic moments in equations (9), (10) and (A4) is valid since the angular velocities of the magnetic moments $`\mathrm{\Omega }_p(\mu _1\mu _2/mR^2r^3)^{1/2}`$ are much smaller than the orbital angular velocities $`w(m/r)^{3/2}/m`$ when $`(H_1H_2)^{1/2}10^{18}`$G.
## 4 INSPIRAL OF CIRCULAR ORBITS
For calculational simplicity we assume that the orbital motion of BNSs has decayed to be circular apart from the adiabatic inspiral (Peters & Mathews 1963, Peters 1964). In general, circular orbit solutions of equation (4) do not exit unless the magnetic moments are aligned perpendicular to the orbital plane. However, since the angular velocities of the magnetic moments $`\mathrm{\Omega }_p(\mu _1\mu _2/mR^2r^3)^{1/2}`$ are much smaller than the orbital angular velocity $`w(m/r)^{3/2}/m`$ when $`(H_1H_2)^{1/2}10^{18}`$G, we can regard the magnetic moment vectors and $`\widehat{𝑳}`$ as time-independent ones over an orbit where $`\widehat{𝑳}`$ is a unit vector orthogonal to the orbital plane. Then, after taking an average of the magnetic term in the acceleration (5), we can obtain orbits of constant separation, $`\ddot{r}=\dot{r}=\widehat{𝒏}𝒗=0`$, $`w=v/r`$, $`𝑳_N=\eta mr^2w\widehat{𝑳}`$ (see similar discussions on spin precessions, Kidder, Will, & Wiseman 1993, Kidder 1995). From the equations of motion for circular orbits, $`\widehat{𝒏}𝒂=\ddot{r}rw^2`$, we can calculate the orbital angular velocity as
$$w^2=\frac{m}{r^3}\left[1+\frac{3}{2\eta m^4}\left(\frac{m}{r}\right)^2\left\{𝝁_1𝝁_23(\widehat{𝑳}𝝁_1)(\widehat{𝑳}𝝁_2)\right\}\right],$$
(13)
where we have used the orbit-averaged relation,
$$\overline{(\widehat{𝒏}𝝁_1)(\widehat{𝒏}𝝁_2)}=\frac{1}{2}\left\{𝝁_1𝝁_2(\widehat{𝑳}𝝁_1)(\widehat{𝑳}𝝁_2)\right\}.$$
(14)
The total energy and the energy loss rate for a circular orbit, averaged over an orbit, can be obtained as
$`E`$ $`=`$ $`\eta {\displaystyle \frac{m^2}{2r}}\left[1{\displaystyle \frac{1}{2\eta m^4}}\left({\displaystyle \frac{m}{r}}\right)^2\left\{𝝁_1𝝁_23(\widehat{𝑳}𝝁_1)(\widehat{𝑳}𝝁_2)\right\}\right],`$ (15)
$`{\displaystyle \frac{dE}{dt}}`$ $`=`$ $`\left({\displaystyle \frac{dE}{dt}}\right)_{GW}+\left({\displaystyle \frac{dE}{dt}}\right)_{EM}`$ (16)
$`=`$ $`{\displaystyle \frac{32}{5}}\eta ^2\left({\displaystyle \frac{m}{r}}\right)^5[1+{\displaystyle \frac{9}{2\eta m^4}}\left({\displaystyle \frac{m}{r}}\right)^2\{𝝁_1𝝁_23(\widehat{𝑳}𝝁_1)(\widehat{𝑳}𝝁_2)\}`$
$`+{\displaystyle \frac{1}{96\eta ^2m^4}}\left({\displaystyle \frac{m}{r}}\right)^2\{3\mu _{eff}^2+(𝝁_{eff}\widehat{𝑳})^2\}],`$
by using equations (6), (9), (12) and (13).
Combining equations (13), (15) and (16), we can express the change rate of the orbital angular velocity $`\dot{w}`$ as a function of $`w`$,
$$\dot{w}=\frac{96}{5}\eta m^{5/3}w^{11/3}\left\{1+\sigma _{mag}(mw)^{4/3}\right\},$$
(17)
where
$$\sigma _{mag}=\frac{5}{\eta m^4}\left\{𝝁_1𝝁_23(\widehat{𝑳}𝝁_1)(\widehat{𝑳}𝝁_2)\right\}+\frac{1}{96\eta ^2m^4}\left\{3\mu _{eff}^2+(𝝁_{eff}\widehat{𝑳})^2\right\}.$$
(18)
By using equation (17), we calculate the accumulated number of GW cycles $`N=(f/\dot{f})𝑑f`$, where $`f=w/\pi `$ is the frequency of the quadrupolar waves. For calculational simplicity, we assume that the two magnetic moments are aligned parallel to the orbital axis of the binary system. In this case, $`\sigma _{mag}`$ becomes a constant value. Then, the accumulated number is integrated as
$$N=N_{grav}+N_{mag},$$
(19)
where $`N_{grav}`$ and $`N_{mag}`$ denote the contributions from the Newtonian gravity term and the magnetic term respectively, and are expressed as
$`N_{grav}`$ $`=`$ $`{\displaystyle \frac{1}{32\pi \eta }}(\pi mf)^{5/3}|_{f_{min}}^{f_{max}},`$ (20)
$`N_{mag}`$ $`=`$ $`{\displaystyle \frac{5}{32\pi \eta }}\sigma _{mag}(\pi mf)^{1/3}|_{f_{min}}^{f_{max}}.`$ (21)
Here $`f_{max}`$ is the exit frequency and $`f_{min}`$ is the entering one of the detector’s bandwidth.
## 5 DISCUSSION
Using 10Hz as the entering frequency and 1000Hz as the exit one, we obtain the contribution to the total number of GW cycles from the magnetic term as
$$N_{mag}=5.9\times 10^3\left(\frac{H_1}{10^{16}\mathrm{G}}\right)\left(\frac{H_2}{10^{16}\mathrm{G}}\right)\left(\frac{m}{2\times 1.4M_{}}\right)^{13/3},$$
(22)
where we assume $`𝝁_1𝝁_2<0`$, $`𝝁_1\widehat{𝑳}`$, $`𝝁_2\widehat{𝑳}`$,<sup>3</sup><sup>3</sup>3 This set of configurations, $`𝝁_1𝝁_2<0`$, $`𝝁_1\widehat{𝑳}`$ and $`𝝁_2\widehat{𝑳}`$, is the most stable one, i.e., the EM interaction energy in equation (7) becomes the minimum value. However, note that the alignment rate of the magnetic moments $`w_{ali}`$ is smaller than the orbital inspiral rate $`w_{ins}(m/r)^4/m`$ when $`H10^{16}(r/10^{11}\mathrm{cm})^{1/4}`$ G, since the alignment rate $`w_{ali}`$ can be estimated as $`w_{ali}(\mathrm{energy}\mathrm{loss}\mathrm{rate}\mathrm{due}\mathrm{to}\mathrm{magnetic}\mathrm{dipole}\mathrm{radiation})/E_{DD}(\mu \mathrm{\Omega }_p^2)/(\mu ^2/r^3)\mu ^4/r^3m^2R^4`$ when $`\mu _1\mu _2\mu `$. and $`m_1=m_2`$. Note that the contribution of the EM radiation reaction (the second term in equation (18)) is much less than that of the dipole-dipole interaction (the first term in equation (18)). The maximum magnetic field allowed by the scalar virial theorem is $`10^{18}`$G for NSs (Chandrasekhar 1981, see also Bocquet et al. 1995). If NSs in the inspiralling BNSs have such magnetic fields $`10^{17}\mathrm{G}H10^{18}\mathrm{G}`$, the effects of the magnetic fields can change more than one cycle in the accumulated cycles. However, if we consider that the maximum value of the observed fields $`10^{16}`$ G is the upper limit for the magnetic fields of NSs, the magnetic term will make negligible contributions to the accumulated phase, contrary to the crude estimate in §1 and equation (1). Consequently, the magnetic fields of NSs will not present difficulties for the detection of GWs from BNSs, if the upper limit for the magnetic fields of NSs is less than $`10^{16}`$G.
The magnetic term in equation (17) has the same dependence on the angular velocity $`w`$ as 2PN terms. We can see this dependence from the 2PN expression for the frequency sweep (Blanchet et al. 1995),
$`\dot{w}`$ $`=`$ $`{\displaystyle \frac{96}{5}}\eta m^{5/3}w^{11/3}\{1({\displaystyle \frac{743}{336}}+{\displaystyle \frac{11}{4}}\eta )(mw)^{2/3}+(4\pi \beta )(mw)`$ (23)
$`+({\displaystyle \frac{34103}{18144}}+{\displaystyle \frac{13661}{2016}}\eta +{\displaystyle \frac{59}{18}}\eta ^2+\sigma )(mw)^{4/3}\},`$
where the spin-orbit parameter is $`\beta =\frac{1}{12}\mathrm{\Sigma }_p(113m_p^2/m^2+75\eta )\widehat{𝑳}𝝌_p`$, the spin-spin parameter is $`\sigma =(\eta /48)(247𝝌_1𝝌_2+721\widehat{𝑳}𝝌_1\widehat{𝑳}𝝌_2)`$ and $`𝝌_p=𝑺_p/m_p^2`$. As we can see from equations (17) and (23), we cannot distinguish the magnetic term $`\sigma _{mag}`$ from the spin-spin term $`\sigma `$ at 2PN order. Only the sum of the magnetic term and the spin-spin term can be deduced from the frequency evolution of 2PN order. This degeneracy might be broken by examining the wave-form modulations caused by the spin-induced precession of the orbit (Apostolatos et al. 1994, Kidder 1995). However, the analysis of Poisson & Will (1995) shows that the measurement error on the spin-spin term is $`\mathrm{\Delta }\sigma 17.3`$ assuming that the SNR is 10, $`\beta =0`$, $`\sigma =0`$ and there is no prior information.<sup>4</sup><sup>4</sup>4 Because of some simplifying assumptions in Poisson & Will (1995), the true measurement error is still uncertain. However, this will not affect the discussion since the true measurement error will be within a factor of the order of 2 (Balasubramanian, Sathyaprakash, & Dhurandhar 1996, Balasubramanian & Dhurandhar 1998, Nicholson & Vecchio 1998). On the other hand, the contribution of the magnetic term is estimated as
$$\sigma _{mag}=2.9\times 10^3\left(\frac{H_1}{10^{16}\mathrm{G}}\right)\left(\frac{H_2}{10^{16}\mathrm{G}}\right)\left(\frac{m}{2\times 1.4M_{}}\right)^4,$$
(24)
where we assume $`𝝁_1𝝁_2<0`$, $`𝝁_1\widehat{𝑳}`$, $`𝝁_2\widehat{𝑳}`$ and $`m_1=m_2`$. We can see from this equation that the contribution of the magnetic fields of NSs is much smaller than the measurement error on the spin-spin term. Therefore, the effects of the magnetic fields of NSs are also negligible for parameter estimation with moderate SNR if we consider that the maximum value of the observed fields $`10^{16}`$G is the upper limit for the magnetic fields of NSs.
Conservatively, there are some other reasons to consider that the magnetic fields of inspiralling NSs are smaller than $`10^{16}`$G. There are three observed BNSs in our Galaxy and nearby globular cluster that will merge in less than a Hubble time: PSR B1913+16 (Taylor & Weisberg 1989), PSR B1534+12 (Wolszczan 1991, Stairs et al. 1998) and PSR B2127+11C (Prince et al. 1991). From the spin-down rate, the magnetic field strength is estimated as of order $`10^{10}`$G for all pulsars in these binary systems. If such a value is typical of the magnetic field strength,<sup>5</sup><sup>5</sup>5 However, note that no radio pulsars have magnetic fields above $`10^{14}`$G, and hence there may be a selection bias (Baring & Harding 1998). the magnetic terms will be negligible. Moreover, if the decay time of the magnetic fields is shorter than the coalescence time (Heyl and Kulkarni 1998, Thompson and Duncan 1996, Goldreich & Reisenegger 1992, Shalybkov and Urpin 1995), of course, the magnetic fields will not be concerned, although the magnetic field evolution of isolated NSs is an unresolved issue. It may be also difficult for BNSs including magnetar to be formed because of a large recoil velocity (Duncan & Thompson 1992).
In this paper, we regard NSs as spherical compact bodies. However, when we consider NSs as extended bodies, we have to take into account of the quadrupole effects induced by magnetic fields. The magnetic fields are a source of non-hydrostatic stress in the interiors of NSs. A magnetic dipole moment $`\mu `$ would give rise to moment differences of order
$$ϵ:=\frac{I_cI_a}{I_a}\frac{R^4H^2}{m^2},$$
(25)
where $`I_c`$ and $`I_a`$ refer to the moments of inertia about the dipole axis and about an axis in the magnetic equator. Then, the gravitational potential $`\eta m^2/r`$ between NSs is modified by an amount of order $`mI_aϵ/r^3\mu ^2/r^3`$. This is the same order as the EM interaction term in equation (2) when $`\mu _1\mu _2`$. Therefore, an accurate ellipticity $`ϵ`$ in equation (25) is needed to determine the magnetic effects within a factor. (Bocquet et al. 1995, Bonazzola & Gourgoulhon 1996, Konno, Obata, & Kojima 1999). Note that there are such quadrupole effects even if only one companion has magnetic moment.
Although the effects of the magnetic fields of NSs will be negligible for observations of GWs, they might be concerned with gamma ray bursts (GRBs). The BNS merger is one of models of GRB sources. (see, e.g., Piran 1999 for a review). If a NS in the binary system has strong magnetic fields, the total energy emitted by EM waves until coalescence can be estimated from equations (13), (16) and (17) as $`(dE/dt)_{EM}𝑑t\pi ^2\mu _{eff}^2(f_{max}^2f_{min}^2)/144\eta m10^{46}(H/10^{16}\mathrm{G})^2(f_{max}/10^3\mathrm{Hz})^2`$ ergs. This energy will be radiated at very low frequency $`10^3`$ Hz, which is difficult to be observed by the present radio telescope. Furthermore such radiation cannot propagate a plasma if an electron density is larger than $`0.01\mathrm{cm}^3`$ since the plasma frequency is larger than the radiation frequency (e.g. Spitzer 1962). However, this energy may be converted to the thermal energy of the surrounding plasma efficiently if the electron density $`n_e`$ is sufficiently high since the electron-electron relaxation time is about $`t_{rel}1(n_e/10^{11}\mathrm{cm}^3)^1(kT_e/2\mathrm{k}\mathrm{e}\mathrm{V})^{3/2}`$ s where $`T_e`$ is the electron temperature (e.g. Spitzer 1962). Then, this thermal radiation might explain the precursory X-ray emission $`10`$ s before the onset of the GRB observed by the Ginga satellite, in which the total energy of the X-ray precursor emission is estimated to be about $`10^{46}(d/100\mathrm{M}\mathrm{p}\mathrm{c})^2`$ ergs (Murakami et al. 1991). Even though the strong magnetic fields are not relevant to the X-ray precursor in GRB, the EM radiation can be the EM signature of the coalescing BNSs. Therefore it is an interesting future problem to investigate the conversion of the low frequency EM radiation to the higher frequency one.
We would like to thank H. Sato and T. Nakamura for continuous encouragement and useful discussions. We are also grateful to T. Tanaka, R. Nishi, K. Nakao, T. Harada and K. Omukai for useful discussions. This work was supported in part by Grant-in-Aid for Scientific Research Fellowship (No.9627: KI) and (No.9402: KT) of the Japanese Ministry of Education, Science, Sports and Culture.
## Appendix A ELECTROMAGNETIC RADIATION FROM A MOVING MAGNETIC MOMENT
In this section we review the radiation from a moving magnetic dipole moment. We consider a particle with only a magnetic dipole moment $`𝝁^{}`$ in its rest frame $`K^{}`$. A moving magnetic dipole moment with velocity $`𝒗=\dot{𝒙}`$ relative to an observer frame $`K`$ also has an associated electric dipole moment. The apparent electric dipole moment is
$$𝒑=𝒗\mathbf{\times }𝝁,$$
(A1)
where $`𝝁=𝝁^{}\frac{\gamma }{\gamma +1}(𝒗𝝁^{})𝒗`$ is the magnetic moment observed in K and $`\gamma =(1v^2)^{1/2}`$ (Jackson 1998). Therefore, in the observer frame $`K`$ the magnetization density $`𝑴`$ and electric polarization density $`𝑷`$ are given by
$`𝑴(t,𝒛)`$ $`=`$ $`𝝁(t)\delta [𝒛𝒙(t)],`$
$`𝑷(t,𝒛)`$ $`=`$ $`𝒑(t)\delta [𝒛𝒙(t)],`$ (A2)
where $`𝒑`$ is given by equation (A1). Recalling that the moving magnetic moment is equivalent to a current $`𝑱=\mathbf{\times }𝑴+\dot{𝑷}`$ (Jackson 1998), we can calculate the electric and magnetic fields. For our purpose, it is sufficient to obtain the radiative parts which fall off as the inverse of the distance $`D^1`$. The radiation field $`𝑩^{rad}`$ of this moving magnetic moment is given by (e.g. Heras 1994)
$`𝑩^{rad}(t,𝒛)`$ $`=`$ $`{\displaystyle \frac{3\widehat{𝒅}\mathbf{\times }(\widehat{𝒅}\mathbf{\times }𝝁𝒑)(\widehat{𝒅}𝒂)^2}{DW^5}}+{\displaystyle \frac{3\widehat{𝒅}\mathbf{\times }(\widehat{𝒅}\mathbf{\times }\dot{𝝁}\dot{𝒑})(\widehat{𝒅}𝒂)}{DW^4}}`$ (A3)
$`+{\displaystyle \frac{\widehat{𝒅}\mathbf{\times }(\widehat{𝒅}\mathbf{\times }𝝁𝒑)(\widehat{𝒅}\dot{𝒂})}{DW^4}}+{\displaystyle \frac{\widehat{𝒅}\mathbf{\times }(\widehat{𝒅}\mathbf{\times }\ddot{𝝁}\ddot{𝒑})}{DW^3}},`$
where $`𝑫(t)=𝒛𝒙(t)`$, $`\widehat{𝒅}=𝑫(t)/|𝑫(t)|=𝑫(t)/D`$, $`W=1𝒗\widehat{𝒅}`$, $`𝒂=\dot{𝒗}`$ and the right-hand side of this equation is evaluated at the retarded time $`t^{}`$, i.e., $`t^{}+D(t^{})=t`$. When we assume that the magnetic dipole moment vector $`𝝁(t)`$ is a constant one, the above equation (A3) can be calculated as
$`𝑩_0^{rad}`$ $`=`$ $`{\displaystyle \frac{1}{D}}\left\{\widehat{𝒅}\mathbf{\times }(\widehat{𝒅}\mathbf{\times }𝝁)(\widehat{𝒅}\dot{𝒂})\widehat{𝒅}\mathbf{\times }\ddot{𝒑}\right\}`$ (A4)
$`=`$ $`{\displaystyle \frac{1}{D}}(\widehat{𝒅}𝝁)\left\{(\widehat{𝒅}\dot{𝒂})\widehat{𝒅}\dot{𝒂}\right\},`$
up to the leading order term of $`𝒗`$ and $`|𝒙|/D`$. Note that the radiation field $`𝑩_0^{rad}`$ becomes the same equation as (A4) up to leading order even if we use $`𝝁^{}`$ instead of $`𝝁`$. The power radiated per unit solid angle is given by (e.g. Landau & Lifshitz 1975)
$$\frac{dP_0}{d\mathrm{\Omega }}=\frac{1}{4\pi }\left(B_0^{rad}\right)^2D^2=\frac{1}{4\pi }(\widehat{𝒅}𝝁)^2\left\{|\dot{𝒂}|^2(\widehat{𝒅}\dot{𝒂})^2\right\}.$$
(A5)
The total instantaneous power is obtained by integrating equation (A5) over all solid angle as
$$P_0=\frac{2}{15}\left\{2\mu ^2|\dot{𝒂}|^2(𝝁\dot{𝒂})^2\right\}.$$
(A6)
## Appendix B THE CONTRIBUTION OF ELECTROMAGNETIC FIELDS TO THE MASS
Thus far we have ignored the contribution of the EM fields to the mass of a compact body in a binary system. In this section, we estimate the correction to the mass by the EM fields. We have implicitly defined the mass $`m`$ as that of the isolated spherical body, i.e., the orbital separation of the binary system is infinity. Therefore, the mass $`m`$ includes the self-energy of the EM fields. The self-energy $`m_{mag}`$ of the magnetic fields outside the compact body with a magnetic dipole moment $`𝝁`$ is obtained by
$$m_{mag}=\frac{H^2}{8\pi }d^3x^{}=\frac{1}{8\pi }\frac{3\mu ^2\mathrm{cos}^2\theta ^{}+\mu ^2}{r^6}d^3x^{}=\mu ^2\left[\frac{1}{3r^3}\right]_R^{\mathrm{}}=\frac{\mu ^2}{3R^3},$$
(B1)
where $`R`$ is the radius of the compact body. (The effects of the magnetic fields inside the compact body are discussed in §5.)
Next, we consider the gravitational interaction between the above compact bodies with magnetic moments. Here we note that at a finite separation the gravitational field of a mass with magnetic dipole moment $`𝝁`$ is different from that of a point mass without magnetic dipole moment even if the masses are the same value since EM fields have an extent. Therefore, the mass $`m`$ in the gravitational potential term $`\eta m^2/r`$ in equation (2) suffers a small correction. This correction can be estimated by evaluating the “gravitational potential” produced by the energy of the EM fields of the magnetic dipole moment $`𝝁`$. The “gravitational potential” $`\varphi `$ at $`𝒙`$ is obtained by
$`\varphi (𝒙)={\displaystyle \frac{H^2}{8\pi |𝒙𝒙^{}|}d^3x^{}}={\displaystyle \frac{\mu ^2}{8\pi }}{\displaystyle \frac{3\mathrm{cos}^2\theta ^{}+1}{r^6|𝒙𝒙^{}|}d^3x^{}}.`$ (B2)
Here we expand $`1/|𝒙𝒙^{}|`$ by the spherical harmonics $`Y_{lm}(\theta ,\phi )`$,
$$\frac{1}{|𝒙𝒙^{}|}=4\pi \underset{l=0}{\overset{\mathrm{}}{}}\underset{m=l}{\overset{l}{}}\frac{1}{2l+1}\frac{r_<^l}{r_>^{l+1}}Y_{lm}^{}(\theta ^{},\phi ^{})Y_{lm}(\theta ,\phi ),$$
(B3)
where $`r_<(r_>)`$ is the smaller (larger) one of $`|𝒙|`$ and $`|𝒙^{}|`$. Noting the orthonormality of the spherical harmonics, $`Y_{l^{}m^{}}^{}(\theta ,\phi )Y_{lm}(\theta ,\phi )𝑑\mathrm{\Omega }=\delta _{l^{}l}\delta _{m^{}m}`$, and a relation, $`3\mathrm{cos}^2\theta ^{}+1=2\sqrt{4\pi /5}Y_{20}(\theta ^{},\phi ^{})+2\sqrt{4\pi }Y_{00}(\theta ^{},\phi ^{})`$, the equation (B2) is calculated as
$`\varphi (𝒙)={\displaystyle \frac{1}{r}}\left[{\displaystyle \frac{\mu ^2}{3R^3}}+{\displaystyle \frac{\mu ^2(3\mathrm{cos}\theta 1)}{10Rr^2}}{\displaystyle \frac{\mu ^2\mathrm{cos}^2\theta }{4r^3}}\right].`$ (B4)
The first term in the bracket on the right-hand side of the above equation comes from the total self-energy of the EM fields in equation (B1), and the last two terms are the corrections due to the extent of the EM fields. The order of the correction to the mass is estimated as $`\delta m/m(\mu ^2/Rr^2)/m\mu ^2/mRr^2`$. On the other hand, the correction due to the dipole-dipole interaction is of order $`\mu _1\mu _2/m^2r^2`$ from equation (17). Therefore, the correction due to the extent of the magnetic fields are smaller than that due to the dipole-dipole interaction when $`\mu _1\mu _2`$. We have also confirmed that the contribution of the EM fields to the quadrupole moments $`I_{ij}`$ is of order $`\mu ^2/mRr^2`$. |
warning/0001/hep-ph0001265.html | ar5iv | text | # A Phenomenological Description of 𝜋⁻Δ⁺⁺ Photo- and Electroproduction in Nucleon Resonance Region
## 1 Introduction.
In the total photoabsorption cross section on the proton above 400 MeV the two pion production becomes possible, as it is clearly shown by several data sets from bubble chamber real photon experiments , as well as from exclusive electron scattering experiments and recent measurements with real photon beams with the DAPHNE and SAPHIR detectors. From the data one can see that the cross section for this process shows a steep rise at photon energies above threshold and then exhibits a smooth decrease, although still representing a large fraction of the total hadron photoproduction at $`W>1.8`$ GeV. It is well established that a large contribution to this process is given by the intermediate formation of $`\mathrm{\Delta }`$(1236), $`\rho `$(770) and their subsequent decay in the cascade processes
$$\gamma _{r,v}p\mathrm{\Delta }^{++}\pi ^{}(p\pi ^+)\pi ^{}$$
(1)
$$\gamma _{r,v}p\rho ^0pp(\pi ^+\pi ^{})$$
(2)
where indices $`r,v`$ stand for real and virtual photons.
An important role in quasi two-body reactions (1),(2) is played by nucleon resonances, whose excitation is clearly seen in the total photohadronic cross section as well as in exclusive channels like single pion photoproduction . Many of such nucleon resonances have an appreciable branching ratio in the $`\mathrm{\Delta }\pi `$ and $`\rho N`$ channels (see Table I) and for $`N^{}`$ heavier than 1.7 GeV multipion decay channels become dominant. Therefore reactions (1),(2) offer a promising opportunity to study the structure of high lying resonances composed by light $`u`$ and $`d`$ quarks. Moreover, quark model calculations based on the three constituent quark picture predict more states than those found in experimental searches. Some parameters for such ”missing” states are shown in Table II: many of them are expected to have a strong branching ratio in the two-pion final state and weak or absent coupling to the single pion: this could be a reason for them having escaped detection in experiments with pion beams and in single pion photoproduction. An additional prediction for these unobserved states is that some of them could have sizeable electromagnetic couplings, similar to other states observed in electromagnetic production. Therefore, sizeable electromagnetic couplings and strong decays through $`\mathrm{\Delta }\pi `$ or $`\rho N`$ channels could make poorly known states more visible as well as make some of the ”missing” states appear in measurements of photo- or electroproduction of the two-pion final states. It is the aim of a wide experimental program at new facilities like TJNAF to investigate $`N^{}`$ structure and to perform a search for missing states by measuring the multipion production exclusive channels produced via electromagnetic interaction. In this paper we will focus on the particular two-pion production channel that proceeds through the $`\mathrm{\Delta }^{++}\pi ^{}`$ intermediate state.
An important feature in the study of resonances in reactions (1),(2) is the strong contribution from non-resonant processes , creating a ”continuum” in which resonant states are embedded. For this reason, common methods for resonance investigation such as multipole amplitude analysis , which essentially use resonance dominance, may be not very effective for reactions (1),(2).
Different approaches to describe double pion photoproduction have been presented in a few papers, where models based on a variety of tree-level diagrams and a few baryon resonances have been used to calculate the two pion production. The limited number of resonances included however makes them applicable only for W lower than about 1.6 GeV; moreover non-resonant terms are not always corrected for unitarity absorption effects: actually, a very important problem in the investigation of $`N^{}`$ in exclusive meson production by photons is the description of non-resonant processes at $`W>1.6`$ GeV, where many competitive channels open up and coupled channel calculations appear to be very difficult due to restricted knowledge of hadron couplings as well as to computational problems.
A phenomenological description of the interaction with open inelastic channels in the initial and final states (ISI and FSI) has been proposed first in . In this approach particle absorption in the initial and final states accounts for the complexity of interactions with inelastic channels; the absorptive coefficients have been usually obtained assuming a diffractive character of the interaction of ingoing and outgoing particles; this assumption could be justified at W values above $`N^{}`$ excitation region, while at $`W<2.0`$ GeV other mechanisms like especially s-channel processes should contribute in ISI and FSI significantly. Moreover, the high value of the effective coupling coefficients, close to the physical limit, adopted in , at relatively small $`W1.61.7`$ GeV looks somewhat artificial and not physically justified.
The need to develop a method for $`N^{}`$ electromagnetic form factor investigation from the experimental two-pion photo- and electroproduction data led us to the development of a phenomenological approach based on minimal model ingredients. This was accomplished parametrizing, under simple assumptions, the main two-pion photo- and electroproduction mechanisms and using experimental data to determine the corresponding parameters.
This paper is our first step in establishing such phenomenological approach: we focused our attention on the particular quasi-two-body channel (1) and set up a simple model including all known resonant states contributing to this reaction and non-resonant processes starting from a minimal set of mechanisms proposed in . We used existing data for the pion electromagnetic form factor to provide a description of non-resonant terms for $`Q^2>0`$, as well as data about the strong form factor in $`\pi N\mathrm{\Delta }`$ vertex to take into account the particle size in hadronic vertices. Unitarity effects manifesting in the competition of many hadronic channels in non-resonant processes were taken into account effectively, implementing initial and final state interactions (ISI and FSI), in form of absorptive corrections; a new feature of our approach is that we developed a specific parametrisation of ISI and FSI mechanisms responsible for channel coupling at W $`2`$ GeV, based on experimental information about hadronic scattering amplitudes in this energy region. In this regard, we examined some aspects of the $`N^{}`$ electromagnetic vertex dressing as related to our implementation of higher order corrections in form of ISI and FSI. We studied also the transition to the higher energy region, W $`>`$ 2 GeV, where the pure tree-level scheme corrected for finite size (strong form factor) and unitarity effects (ISI-FSI) fails in reproducing the experimental cross section: we found that a Regge trajectory exchange, in the particular form recently proposed, can be a valid continuation into the higher energy region.
Our model relates $`N^{}`$ electromagnetic helicity couplings $`A_{1/2}`$, $`A_{3/2}`$, $`C_{1/2}`$ and measured cross sections for reaction (1) induced both by real and virtual photons, therefore offering a way to attempt the measurement of $`N^{}`$ contributions from a comparison to experimental data or a fit.
## 2 Helicity amplitudes and differential cross section.
The helicity representation has been chosen to describe the amplitudes for reaction (1). The differential cross section for reaction (1) induced by real or virtual photons in the one-photon exchange approximation can be expressed as :
$$\frac{d\sigma }{d\mathrm{\Omega }_\pi ^{}}=\frac{1}{4K_LM_N}\left[(4\pi \alpha )\frac{1\epsilon }{Q^2}\frac{1}{2}L_{\mu \nu }W^{\mu \nu }\right]\frac{1}{(2\pi )^2}\frac{p_\pi ^{}}{4W}$$
(3)
$`K_L={\displaystyle \frac{W^2M_N^2}{2M_N}},`$ $`\alpha ={\displaystyle \frac{1}{137}}`$
where $`W`$ is CM total energy, $`p_\pi ^{}`$ is the pion three momentum modulus in the CM frame, $`K_L`$ is the ”equivalent” photon energy, $`M_N`$ is the nucleon mass, $`Q^2`$ = - $`q^2`$ where $`q^\mu `$ is the real-virtual photon four-momentum, $`\epsilon `$ is the photon polarization parameter, $`L_{\mu \nu }`$ is the leptonic tensor well-known from QED. The information about hadronic production is contained in the hadronic tensor $`W^{\mu \nu }`$ which is a bilinear combination of hadronic currents:
$$W^{\mu \nu }=\frac{1}{2}\underset{\lambda _p\lambda _\mathrm{\Delta }}{}J_\mu ^{}J_\nu $$
(4)
related to the helicity amplitudes $`\lambda _\mathrm{\Delta }\left|T\right|\lambda _\gamma \lambda _p`$ according to:
$$\epsilon _\mu (\lambda _\gamma )J^\mu (\lambda _p\lambda _\mathrm{\Delta })=\lambda _\mathrm{\Delta }\left|T\right|\lambda _\gamma \lambda _p$$
(5)
where $`\lambda _\mathrm{\Delta }`$, $`\lambda _p`$, $`\lambda _\gamma `$ are helicities of $`\mathrm{\Delta }`$, proton and photon, respectively, $`\epsilon ^\mu `$ is the four – vector associated to the photon polarization state $`\lambda _\gamma `$ and all variables are evaluated in the reaction CM frame.
Being the $`\mathrm{\Delta }`$ an unstable particle, it is necessary to fold (3) on its mass distribution as follows:
$$\frac{d\sigma }{d\mathrm{\Omega }_\pi ^{}}=𝑑M^2\frac{1}{\pi }\frac{M_\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{\Delta }}{(M^2M_\mathrm{\Delta }^2)^2+M_\mathrm{\Delta }^2\mathrm{\Gamma }_\mathrm{\Delta }^2}\frac{d\sigma }{d\mathrm{\Omega }_\pi ^{}}(M^2)$$
(6)
where $`M^2`$ is the squared running invariant mass of $`\mathrm{\Delta }`$, while $`M_\mathrm{\Delta }`$, $`\mathrm{\Gamma }_\mathrm{\Delta }`$ are $`\mathrm{\Delta }`$ mass and width, respectively.
## 3 The tree-level diagrams.
The mechanisms of reaction (1) were described by a minimal set of Feynman tree-level diagrams presented in Fig. 1. These diagrams can be subdivided into $`N^{}`$ contributions (Fig.1a) and a group of non-resonant processes, or Born terms (Fig.1b-e). The resonant part (Fig.1a) involves all relevant $`N^{}`$ and $`\mathrm{\Delta }^{}`$ excitations in s-channel. The non-resonant amplitudes correspond to the same set of mechanisms considered in . New features of our approach in this respect are:
1. implementation of electromagnetic vertex functions to describe the behaviour of non-resonant amplitudes at $`Q^2>0`$ and
2. implementation of a strong $`\pi N\mathrm{\Delta }`$ form factor, relativistically invariant and based on the $`NN`$ scattering analysis, to take into account the finite size of hadrons involved in non-resonant mechanisms.
### 3.1 Non-resonant processes.
The non-resonant processes are composed by the ”seagull” or ”contact” term (Fig. 1b) (which also naturally arises when considering the $`\pi N\mathrm{\Delta }`$ vertex as described by an effective meson-baryon Lagrangian and then introducing the interaction with the electromagnetic field by the minimal coupling prescription), the t-channel pion-in-flight diagram (Fig.1c), the u-channel delta-in-flight diagram (Fig.1e) and by the s-pole nucleon term (Fig.1d). To treat photon, pion and delta off-shell we introduced form factors in the corresponding vertices. Therefore Born terms (Fig.1b-e) are functions of $`Q^2`$ and Mandelstam t variables. For the vertex function evaluation we used a compilation of experimental data about electromagnetic and strong form factors ; these data provide a reliable vertex function evaluation in the relevant region of $`Q^2`$ and t ($`Q^2<2`$ $`GeV^2`$, $`t<2`$ $`GeV^2`$).
In momentum space, helicity amplitudes corresponding to the contact term are given by
$$f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}^c(W,\theta )=g_c(Q^2,t)\overline{u}_\mu (p_2,\lambda _\mathrm{\Delta })u(p_1,\lambda _p)\epsilon _\mu (q,\lambda _\gamma )$$
(7)
where W is the usual invariant CM energy, $`\theta `$ is the pion production angle in the hadronic CMS, $`p_1`$ and $`p_2`$ are the target nucleon and $`\mathrm{\Delta }`$-particle four momenta, $`q`$ is the photon four momentum and $`u_\mu `$, $`\epsilon _\mu `$ and $`u`$ are Rarita-Schwinger spinor-tensor for $`\mathrm{\Delta }`$-particle with $`\lambda _\mathrm{\Delta }`$ helicity, polarisation vector for photon with $`\lambda _\gamma `$ helicity and spinor for target nucleon with $`\lambda _p`$ helicity, respectively; $`g_c(Q^2,t)`$ is an effective ”seagull” term vertex function.
The pion-in-flight contribution in momentum representation reads:
$`f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}^{pif}(W,\theta )=`$
$`g_\pi (Q^2,t){\displaystyle \frac{(2p_\pi ^\mu q^\mu )\epsilon _\mu (q,\lambda _\gamma )}{tm_\pi ^2}}\overline{u}_\nu (p_2,\lambda _\mathrm{\Delta })u(p_1,\lambda _p)(q^\nu p_\pi ^\nu )`$ (8)
where $`p_\pi `$ is the pion momentum, $`m_\pi `$ is the pion mass, $`g_\pi `$ is the product of strong and electromagnetic vertex functions:
$$g_\pi (Q^2,t)=G_{\pi ,em}(Q^2)G_{\pi N\mathrm{\Delta }}(t)$$
(9)
The electromagnetic vertex function in this case was described by the well-known pole fit of pion form factor
$$G_{\pi ,em}(Q^2)=\frac{1}{\left(1+\frac{Q^2(GeV^2)}{\mathrm{\Lambda }_\pi ^2}\right)}\frac{1}{G_{\pi N\mathrm{\Delta }}(t_{min})}$$
(10)
where $`t_{min}`$ corresponds to pion production in hadronic CMS at zero degree angle; the factor $`\frac{1}{G_{\pi N\mathrm{\Delta }}(t_{min})}`$ reflects the way of $`G_{\pi em}(Q^2)`$ form factor extraction from single pion electroproduction data presented in . Analysis in yielded $`\mathrm{\Lambda }_\pi ^2=0.462`$ $`GeV^2`$ and we used this value in calculations. However considering uncertainties in as well as other data about pion electromagnetic form factor , we could allow for a variation of $`\mathrm{\Lambda }_\pi ^2`$ cut-off parameter within $`0.40.5`$ GeV<sup>2</sup>. Concerning the $`\pi N\mathrm{\Delta }`$ t-dependence, we introduce as vertex function a strong form factor successfully applied in $`NNN\mathrm{\Delta }`$ relativistic transition potentials
$$G_{\pi N\mathrm{\Delta }}(t)=g_0\frac{\mathrm{\Lambda }^2m_\pi ^2}{\mathrm{\Lambda }^2t}$$
(11)
The interaction constant $`g_0`$ and cut-off parameter $`\mathrm{\Lambda }`$ are: $`g_0=2.1/m_\pi `$ and $`0.6<\mathrm{\Lambda }<0.9`$ GeV. We found that best data reproduction corresponds to $`\mathrm{\Lambda }=0.75`$ GeV. Gauge invariance implies equality of the ”seagull” term coupling $`g_c(Q^2,t)`$ and the pion-in-flight term coupling $`g_\pi (Q^2,t)`$; we stress that, instead of applying specific electromagnetic and strong vertex functions to the contact term (which would imply to introduce some free parameters to be fitted), we rather fixed the $`\mathrm{p}\mathrm{r}\mathrm{o}\mathrm{d}\mathrm{u}\mathrm{c}\mathrm{t}`$ of them $`g_c(Q^2,t)`$ requiring gauge invariance.
The s-channel nucleon contribution (diagram of Fig.1d) is given by
$$f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}^N(W,\theta )=g_N(Q^2)g_0\frac{(2p_1^\mu +q^\mu )\epsilon _\mu (q,\lambda _\gamma )}{sm_N^2}\overline{u}_\nu (p_2,\lambda _\mathrm{\Delta })u(p_1,\lambda _p)p_\pi ^\nu $$
(12)
where s is Mandelstam invariant $`s=(q+p_1)^2`$, $`m_N`$ is the nucleon mass, $`g_N`$ is the electromagnetic vertex function (we put the strong vertex function for s-channel equal to unity), described by the well-known dipole fit:
$$g_N(Q^2)=\frac{1}{\left(1+\frac{Q^2(GeV^2)}{0.71}\right)^2}$$
(13)
The last contribution to the Born terms that we considered was the $`\mathrm{\Delta }`$-in-flight (diagram of Fig.1e). The expression for such process is:
$$f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}^\mathrm{\Delta }(W,\theta )=2g_\mathrm{\Delta }(Q^2,t)\frac{(2p_2^\mu q^\mu )\epsilon _\mu (q,\lambda _\gamma )}{um_\mathrm{\Delta }^2}\overline{u}_\nu (p_2,\lambda _\mathrm{\Delta })p_\pi ^\nu u(p_1,\lambda _p)$$
(14)
where u is Mandelstam variable corresponding to the crossed invariant momentum transfer $`u=(p_1p_\pi )^2`$; factor 2 accounts for the double electric charge of the $`\mathrm{\Delta }^{++}`$ particle, while again $`g_\mathrm{\Delta }(Q^2,t)`$ is the product of electromagnetic and strong vertex functions. Due to lack of information about $`\mathrm{\Delta }`$’s elastic electromagnetic form factors we followed two ways to evaluate $`g_\mathrm{\Delta }(Q^2,t)`$ vertex function: first, it was derived from gauge invariance of the total Born amplitude:
$$g_\mathrm{\Delta }(Q^2,t)=\frac{g_\pi (Q^2,t)+g_N(Q^2)}{2};$$
(15)
second, we used function (11) for $`\pi N\mathrm{\Delta }`$ vertex, substituting mass $`m_\pi `$ by $`m_\mathrm{\Delta }`$. In the latter choice we found the $`\mathrm{\Delta }`$-in-flight contribution to be negligible, because according to $`\mathrm{\Lambda }`$ cut-off parameter does not exceed 1.3 GeV. Of course in the second evaluation of $`\mathrm{\Delta }`$-in-flight term gauge invariance of the total Born amplitude is lost but, considering that our description of non-resonant terms aims to be phenomenological, we can assume that gauge invariance could be restored by other mechanisms not contributing significantly to the cross section. Moreover, a negligible $`\mathrm{\Delta }`$-in-flight contribution seems to be supported by the comparison with the experimental data.
### 3.2 The absorptive corrections.
It is well known that the tree-level calculation is not able to reproduce the data extracted from experimental analysis of $`\pi ^{}\mathrm{\Delta }^{++}`$ production. One reason is that the lowest order mechanisms for reaction (1) do not incorporate unitarity and hence interaction with other open channels in initial and final states must be taken into account. Direct description of the unitarity effects in the initial and final states would require a simultaneous treatment of all relevant hadronic channels in a coupled channel calculation using hadronic elastic and inelastic scattering amplitudes data. The present status of the hadronic amplitudes knowledge as well as computational capabilities restrict coupled channel calculations to W lower than 1.6 GeV, while investigation of the cross section for process (1) above 1.6 GeV is very important for $`N^{}`$ physics.
The description of open channels competition in non-resonance processes remains a challenging open problem of $`N^{}`$ structure investigation in exclusive meson photo- and electroproduction. The modern status of strong interaction theory does not allow to evaluate hadronic amplitudes and couplings from fundamental principles (i.e. QCD Lagrangian). Therefore we chose to adopt an effective treatment of unitarity effects in initial and final state; the initial and final state interactions responsible for the modification of the amplitude for reaction (1) are schematically depicted in Fig. 2. In the spirit of Vector Meson Dominance (VDM), the ISI was assumed to proceed through a transition between photon and vector meson $`\rho `$. The black blobs describe all transition processes, between $`\rho p`$ initial and $`\rho ^{^{}}p^{^{}}`$ intermediate states as well as between $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}`$ intermediate and $`\pi ^{}\mathrm{\Delta }^{++}`$ final states.
We evaluated ISI and FSI effects in $`\gamma p\pi ^{}\mathrm{\Delta }^{++}`$ reaction in the frame of a simple phenomenological recipe described in , where ingoing and outgoing particle interaction is described by penetration factors in the initial and final states $`f_{ISI}^j`$, $`f_{FSI}^j`$, which depend on the reaction total angular momentum $`j`$; they determine the amplitude for $`\rho p`$ initial state to be transformed into intermediate $`\rho ^{^{}}p^{^{}}`$ state and the amplitude for $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}`$ intermediate state to be transformed into final $`\pi ^{}\mathrm{\Delta }^{++}`$ state. According to , the total Born amplitude $`f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}(\theta ,\phi )`$ was numerically decomposed in partial waves with total angular momentum $`j`$:
$`f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}(\theta ,\phi )`$ $`=`$ $`{\displaystyle \underset{j}{}}f_{\lambda \mu }^jd_{\lambda \mu }^j(\theta )e^{i(\lambda \mu )\phi }`$
$`f_{\lambda \mu }^j`$ $`=`$ $`{\displaystyle 𝑑\mathrm{\Omega }\frac{2j+1}{4\pi }f_{\lambda _\mathrm{\Delta }\lambda _\gamma \lambda _p}^{}(\theta ,\phi )d_{\lambda \mu }^j(\theta )e^{i(\lambda \mu )\phi }}`$ (16)
$`\lambda =\lambda _\mathrm{\Delta },`$ $`\mu =\lambda _\gamma \lambda _p`$
The ISI & FSI modify the $`f_{\lambda \mu }^j`$ decomposition coefficients as:
$`f_{\lambda \mu }^{jcorr}=f_{ISI}^jf_{\lambda \mu }^jf_{FSI}^j`$ (17)
$`\lambda =\lambda _\gamma \lambda _p,\mu =\lambda _\mathrm{\Delta }`$
where $`f_{\lambda \mu }^{jcorr}`$ represent decomposition coefficients for the Born term amplitude modified by ISI & FSI. $`f_{ISI}^j`$ and $`f_{FSI}^j`$ factors, under the assumptions of ref. , can be related with $`S^j`$-matrix elements of $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes at total angular momentum $`j`$ as:
$`f_{ISI}^j`$ $`=`$ $`\lambda _\rho \lambda _p|S^j|\lambda _\rho \lambda _p^{1/2}`$
$`f_{FSI}^j`$ $`=`$ $`\pi \lambda _\mathrm{\Delta }|S^j|\lambda _\mathrm{\Delta }\pi ^{1/2}`$ (18)
where $`\lambda _\rho `$, $`\lambda _p`$, $`\lambda _\mathrm{\Delta }`$ are helicities for $`\rho `$-meson (equal to photon helicity in the VDM picture), for proton and $`\mathrm{\Delta }`$-isobar, respectively; S and T matrix elements are related as:
$$S=1+2iT$$
(19)
The penetration coefficients $`f_{ISI}^j`$ and $`f_{FSI}^j`$ are unambiguously determined by $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes, therefore these amplitudes effectively describe all mechanisms (resonant and non-resonant) responsible for transitions between the initial $`\rho p`$ and intermediate $`\rho ^{^{}}p^{^{}}`$ states as well as between the intermediate $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}`$ and final $`\pi ^{}\mathrm{\Delta }^{++}`$ states (blobs in Fig. 2).
As mentioned in the introduction, a new feature of our approach is that, instead of using a diffractive ansatz like in the previous literature, better suited for higher energies, T-matrix elements for $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering were more appropriately evaluated using data about hadronic scattering in the resonance region, as briefly described in the following paragraphs.
T-matrix elements for $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering were evaluated in an isobar model assuming them to be a superposition of relevant $`N^{}`$ contributions and a smooth background as depicted in Fig. 3. The resonant part of the corresponding amplitudes was evaluated in Breit – Wigner ansatz described in details in sect. 3.3. To obtain $`T_{res}`$ amplitude normalization, we considered a schematical situation, when complexity of all transitions $`\rho p\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}\pi ^{}\mathrm{\Delta }^{++}`$ processes is limited to a single $`N^{}`$ excitation with elastic scattering as the only open channel; for this case, ingoing and outgoing particle absorption should be absent at resonant point ($`W=M_N^{}`$) and the moduli of $`f_{ISI}^j`$ and $`f_{FSI}^j`$ coefficients should be equal to unity; from unitarity conditions and from the standard resonance phase behavior, we determined $`T_{res}^j`$ amplitude normalization as $`f|T_{res}^j|i=i`$ at $`W=M_N^{}`$, therefore
$$f_{ISI}^j,f_{FSI}^j=1+2if|T_{res}^j|i=1$$
(20)
This way we obtained the following relations between elastic $`\pi \mathrm{\Delta }`$ and $`\rho p`$ scattering resonant amplitudes $`\pi \lambda _\mathrm{\Delta }|T_{res}^j|\pi \lambda _\mathrm{\Delta }`$, $`\lambda _\rho \lambda _p|T_{res}^j|\lambda _\rho \lambda _p`$ and $`N^{}`$ decay helicity amplitudes $`\pi \lambda _\mathrm{\Delta }|T|N^{}`$, $`\lambda _\rho \lambda _p|T|N^{}`$:
$`\pi \lambda _\mathrm{\Delta }(\lambda _\rho \lambda _p)|T_{res}^j|\pi \lambda _\mathrm{\Delta }(\lambda _\rho \lambda _p)=`$
$`{\displaystyle \underset{N^{}}{}}\left[{\displaystyle \frac{a_{\lambda _\mathrm{\Delta }(\lambda _\rho \lambda _p)}^{j2}}{M_N^{}^2W^2i\mathrm{\Gamma }_{M_N^{}}(W)M_N^{}}}\right]\left[{\displaystyle \frac{P_{\pi (\rho )}^c}{8\pi (2j+1)W}}\right]`$ (21)
where $`M_N^{}`$, $`\mathrm{\Gamma }_N^{}(W)`$ are masses and W – dependent $`N^{}`$ decay widths, $`P_\pi ^c`$, $`P_\rho ^c`$ are three – momenta moduli of pion and $`\rho `$. Summation in (21) is performed over all $`N^{}`$ contributing to the partial wave of total angular momentum $`j`$ and presented in table III; $`a_{\lambda _\mathrm{\Delta }}^j`$, $`a_{\lambda _\rho \lambda _p}^j`$ are decomposition coefficients of $`\pi \lambda _\mathrm{\Delta }|T|N^{}`$ and $`\lambda _\rho \lambda _p|T|N^{}`$ $`N^{}`$ decay amplitudes through the states of total angular momentum j and related with $`N^{}`$ partial decay widths in helicity representation $`\mathrm{\Gamma }_{\lambda _\mathrm{\Delta }}`$ and $`\mathrm{\Gamma }_{\lambda _\rho \lambda _p}`$ according to (31). Partial decay widths were taken from analysis and transformed from LS to helicity representation.
Our approach should be somewhat modified if dressed $`\gamma pN^{}`$ electromagnetic vertices are used for the resonant part of reaction (1) amplitude (sect 3.3). Indeed, in the complexity of transition mechanisms in ISI & FSI shown in Fig. 2 by black blobs, also the sequence of processes reported in Fig. 4 is effectively taken into account; that clearly represents a dressing mechanism of the $`\gamma pN^{}`$ electromagnetic vertex. But for instance, $`N^{}`$ electromagnetic helicity amplitudes $`A_{1/2}`$, $`A_{3/2}`$ extracted from pion photoproduction data analysis are usually interpreted as corresponding to dressed $`\gamma pN^{}`$ verticies. Therefore, a double counting problem can arise in this case, as dressing effects of the electromagnetic $`N^{}`$ vertex can be sizeable, as discussed in . To avoid double counting, one possibility is to exclude $`\rho pN^{}\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}N^{}\pi ^{}\mathrm{\Delta }^{++}`$ mechanisms from ISI & FSI. To work out a phenomenological prescription for such exclusion we considered the above mentioned schematical situation, where the complexity of $`\rho p\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}\pi ^{}\mathrm{\Delta }^{++}`$ transition mechanisms is restricted to single $`N^{}`$ excitation with only elastic scattering as open channel. Unitarity for Breit – Wigner formula can be expressed as
$$\mathrm{\Gamma }_{tot}=\underset{i}{}\mathrm{\Gamma }_i,$$
(22)
where $`\mathrm{\Gamma }_{tot}`$, $`\mathrm{\Gamma }_i`$ are total and partial $`N^{}`$ decay widths. For single $`N^{}`$ with only elastic open channel
$$\mathrm{\Gamma }_{tot}=\mathrm{\Gamma }_{el},$$
(23)
where $`\mathrm{\Gamma }_{tot}`$, $`\mathrm{\Gamma }_{el}`$ are total and elastic $`N^{}`$ decay width; our assumption insures that all transition mechanisms in $`\rho p\rho ^{^{}}p^{^{}}`$, $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}\pi ^{}\mathrm{\Delta }^{++}`$ are represented by $`\rho pN^{}\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}N^{}\pi ^{}\mathrm{\Delta }^{++}`$ processes; therefore exclusion of $`N^{}`$ excitation from ISI & FSI should lead to complete absorption of ingoing and outgoing particle or to zero values for $`f_{ISI}^j`$, $`f_{FSI}^j`$ penetration factors. As it follows from (19) the substitution
$$T_{res}^j\frac{1}{2}T_{res}^j$$
(24)
would correspond to vanishing $`f_{ISI}^j`$, $`f_{FSI}^j`$ coefficients and, therefore, could be considered as an empirical prescription for excluding $`\rho pN^{}\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}N^{}\pi ^{}\mathrm{\Delta }^{++}`$ mechanisms from ISI & FSI treatment.
The complete absorption of Born term amplitudes takes place only for such schematical situation (single $`N^{}`$ excitation with only elastic open channel) on which we based our exclusion prescription. In actual situation prescription (24) gives only partial absorption of ingoing and outgoing particles, since $`N^{}`$ contributing to $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering have many open decay channels, while relations (21) – (24) give zero $`f_{ISI}^j`$, $`f_{FSI}^j`$ coefficients only in our schematic assumption. Inelasticities ($`\mathrm{\Gamma }_{tot}\mathrm{\Gamma }_{el}`$) provide instead non-zero $`f_{ISI}^j`$, $`f_{FSI}^j`$ coefficients even at resonant point; $`N^{}`$ off-shell excitations as well as non-resonant processes also give non-zero values of $`f_{ISI}^j`$ and $`f_{FSI}^j`$ coefficients. We assumed that this partial absorption of Born terms can represent ISI & FSI corrections, being the contribution from $`\rho pN^{}\rho ^{^{}}p^{^{}}`$ and $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}N^{}\pi ^{}\mathrm{\Delta }^{++}`$ mechanisms excluded.
Of course our exclusion procedure (substitution (24)) should be applied only if for the resonant part of reaction (1) dressed electromagnetic $`\gamma pN^{}`$ vertices are used. Actually, $`\gamma pN^{}`$ vertices calculated in quark models are generally assumed to be “bare”, i.e. free from higher order corrections (see for instance ): in this case our calculation should not be affected by any double counting and therefore no exclusion procedure is necessary. The information on $`N^{}`$ electromagnetic form factors from the measured cross section could be extracted both with and without application of exclusion procedure: in the first case we would obtain information about dressed $`\gamma pN^{}`$ vertices; in the second about $`\gamma pN^{}`$ vertices without contribution from higher order corrections depicted in Fig. 4. Therefore our approach provides the flexibility of choosing one procedure or the other.
Evaluation of non-resonant part of $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ scattering amplitude is described in more detail in , but we want to report here the main features of our evaluation. We used data on partial-wave $`\pi N`$ elastic cross-sections with definite total angular momentum and isospin; the amplitudes for these partial waves were described by a superposition of all relevant $`N^{}`$ (table 3) and a particular background for each orbital momentum L and total spin S. The resonant part of $`\pi N`$ elastic scattering amplitudes was treated in the same manner as $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering processes; data were used for $`N^{}N\pi `$ couplings. Background for each LS wave was parametrized as a function of W by a simple linear dependence
$$T_{backgrLS}^j=A_{LS}W+B_{LS}$$
(25)
Parameters $`A_{LS}`$ and $`B_{LS}`$ were determined from our fit of $`\pi N`$ elastic scattering data. To calculate background amplitudes for $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering starting from $`\pi N`$ elastic scattering background we used SU(3) flavour symmetry relations. The fit results for $`\pi N`$ partial elastic scattering cross – section are shown in Fig. 5. Dotted curves represent resonant contributions, dashed lines are background contributions, while solid lines correspond to complete amplitudes. As follows from Fig. 5, in $`N^{}`$ excitation region nucleon resonances provide the main contribution in $`\pi N`$ elastic amplitudes, therefore also in penetration coefficients (17). This contribution was completely neglected in previous evaluations . Hence, implementation of s – channel $`N^{}`$ excitation amplitude performed in our approach is particularly important for ISI & FSI treatment in $`N^{}`$ excitation region. For W above 2.0 GeV the contribution of resonant part in $`\pi N`$ elastic scattering partial waves (Fig. 5) falls down drastically and non – resonant processes provide the main contribution for most partial waves. Therefore a diffractive approximation for ISI and FSI absorptive corrections could be justified above $`N^{}`$ excitation region.
To check the reliability of our approach to evaluate ISI and FSI absorptive factors, we also performed calculations of such effects implementing $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes directly provided by the authors of analysis , based on a global unitary fit of $`\pi N`$ scattering T-matrix. The gross features of ISI and FSI absorptive factors as obtained from these two procedures are in reasonable coincidence . It is also worth to note that using the above mentioned description of Born terms, together with our ISI-FSI description, non-resonant processes in our approach do not have any free parameters to be determined from reaction (1).
Actually, our approach for ISI & FSI treatment is mostly phenomenological and our assumptions can be mainly justified by comparison with experimental data (we will discuss it in more detail in the next sections): pure Born terms give rise to a cross section that does not reproduce real photon data at high W by a large factor, while absorption-corrected calculations are able to give a good account of data up to W around 2 GeV. Implementation of pion Regge trajectory exchange according to prescription in Born terms allows to obtain a satisfactory data description up to W=3 GeV as discussed in section 4.
Assuming vector dominance as main mechanism responsible for $`Q^2`$ evolution of coupling with other hadronic channels in the initial state we obtained the following expression for the $`f_{ISI}^j`$ absorptive factor $`Q^2`$ dependence:
$$f_{ISI}^j(Q^2)=\frac{\mathrm{\Lambda }_\pi ^2f_{ISI}^j(Q^2=0)+Q^2}{\mathrm{\Lambda }_\pi ^2+Q^2}$$
(26)
where $`\mathrm{\Lambda }_\pi ^2=0.46`$ GeV<sup>2</sup>.
### 3.3 Resonance contribution.
A simple Breit-Wigner ansatz was chosen to describe the coherent superposition of all relevant $`N^{}`$, $`\mathrm{\Delta }^{}`$ resonant amplitudes (see Table I, Fig. 1a).
$`\lambda _\mathrm{\Delta }\left|T_{res}\right|\lambda _\gamma \lambda _p=`$
$`{\displaystyle \underset{N^{},\mathrm{\Delta }^{}}{}}\pi \lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R{\displaystyle \frac{1}{M_{res}^2W^2i\mathrm{\Gamma }_{res}(W)M_{res}}}\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$ (27)
where $`M_{res}`$, $`\mathrm{\Gamma }_{res}`$ are resonance mass and energy-dependent total width and $`\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$, $`\pi \lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R`$ are electromagnetic production and strong decay amplitudes of $`N^{}`$ with helicity $`\lambda _R=\lambda _\gamma \lambda _p`$, respectively. The $`N^{}`$ off-shell effects were taken into account by the Breit-Wigner propagator in (23) as well as by the W-dependence of resonance width and strong decay amplitudes.
Following a phenomenological approach we related $`\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$ and $`\pi \lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R`$ amplitudes with observables extracted from experimental data analysis. Electromagnetic $`N^{}`$ amplitudes $`\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$ were expressed in terms of commonly used helicity couplings $`A_{1/2}`$, $`A_{3/2}`$ and $`C_{1/2}`$. Comparison of the cross section for only one isolated $`N^{}`$ state calculated according to (3)-(6) with Breit-Wigner formula gives:
$`\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$ $`=`$ $`{\displaystyle \frac{W}{M_{res}}}\sqrt{{\displaystyle \frac{8M_NM_{res}p_{\gamma _R}^{}}{4\pi \alpha }}}\sqrt{{\displaystyle \frac{p_{\gamma _R}^{}}{p_\gamma ^{}}}}A_{1/2,3/2}(Q^2);\left|\lambda _\gamma \lambda _p\right|={\displaystyle \frac{1}{2}},{\displaystyle \frac{3}{2}}`$
$`fortransversephotons`$
$`\lambda _R\left|T_{em}\right|\lambda _\gamma \lambda _p`$ $`=`$ $`{\displaystyle \frac{W}{M_{res}}}\sqrt{{\displaystyle \frac{8M_NM_{res}p_{\gamma _R}^{}}{4\pi \alpha }}}\sqrt{{\displaystyle \frac{p_{\gamma _R}^{}}{p_\gamma ^{}}}}C_{1/2}(Q^2)`$
$`forlongitudinalphotons`$
where $`p_{\gamma _R}^{}`$ and $`p_\gamma ^{}`$ are photon three-momentum moduli in the CM frame at resonance point ($`W=M_{res}`$) and at running W, respectively. Couplings $`A_{1/2}(Q^2)`$, $`A_{3/2}(Q^2)`$, $`C_{1/2}(Q^2)`$ completely describe $`N^{}`$ electromagnetic excitation and their values, calculated in a model or taken from some experimental analysis, can be used to calculate the resonant part of cross section for reaction (1). On the other hand, considering these couplings as parameters our approach could be used to attempt their extraction from a fit of the measured cross section.
The $`N^{}`$ strong decay amplitudes $`\pi \lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R`$ were parametrised through the projection on the set of states with definite total angular momentum j:
$$\pi :\lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R=a_{\lambda _\mathrm{\Delta }}^jd_{\lambda _R\lambda _\mathrm{\Delta }}^j(\theta ^{})\sqrt{\frac{p_{\pi _R}^{}}{p_\pi ^{}}}e^{\lambda _\mathrm{\Delta }i\phi }$$
(29)
where $`a_{\lambda _\mathrm{\Delta }}^j`$ is the decomposition coefficient, that does not depend on the resonance helicity state due to rotational invariance, $`\theta ^{}`$ is the CM pion emission angle, $`p_{\pi _R}^{}`$ and $`p_\pi ^{}`$ are CM pion three-momentum moduli at resonance and running W, respectively. We then related the $`a_{\lambda _\mathrm{\Delta }}^j`$ parameters with the partial decay widths $`\sqrt{\mathrm{\Gamma }_{ls}}`$ of $`N^{}\mathrm{\Delta }^{++}+\pi ^{}`$ in LS- representation extracted in analysis .
This way we obtained the decay amplitudes at the resonance point. General requirements for the amplitude threshold behaviour as well as data analysis of single-pion production by photon and pion beams suggest the introduction of a W dependence of $`N^{}`$ decay amplitudes in (27). Assuming as usually that barrier penetration effects provide a good description of $`W`$\- evolution in $`\pi \mathrm{\Delta }(ls)|T_{dec}|R`$ decay amplitudes and using the parametrization from we arrived to:
$`\pi \mathrm{\Delta }:(ls)\left|T_{dec}\right|R(W)=`$
$`\pi \mathrm{\Delta }:(ls)\left|T_{dec}\right|R(M_{res})\left[{\displaystyle \frac{M_{res}}{W}}{\displaystyle \frac{J_l^2(p_{\pi _R}^{}R)+N_l^2(p_{\pi _R}^{}R)}{J_l^2(p_\pi ^{}R)+N_l^2(p_\pi ^{}R)}}\right]^{1/2}`$ (30)
where $`J_l`$, $`N_l`$ are Bessel’s and Neumann’s functions and the factor in square brackets in (30) appears from barrier penetration ratio for a particle emitted with relative orbital momentum $`l`$, while R is an interaction radius whose value was set to 1 fm. The decay amplitude (30) was then transformed into the helicity representation $`\lambda _\mathrm{\Delta }\left|T_{dec}\right|\lambda _R`$ and using the relationship between two-body decay amplitude and width $`\mathrm{\Gamma }_{\lambda _\mathrm{\Delta }}`$ we obtained:
$`\left|a_{\lambda _\mathrm{\Delta }}^j\right|`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2\pi }M_{res}\sqrt{2j+1}\sqrt{\mathrm{\Gamma }_{\lambda _\mathrm{\Delta }}}}{\sqrt{p_\pi }}}`$
$`p_\pi `$ $`=`$ $`{\displaystyle _{(m_\pi +m_N)^2}^{(Wm_\pi )^2}}𝑑M^2{\displaystyle \frac{1}{\pi }}{\displaystyle \frac{M_\mathrm{\Delta }\mathrm{\Gamma }_\mathrm{\Delta }}{(M^2M_\mathrm{\Delta }^2)^2+M_\mathrm{\Delta }^2\mathrm{\Gamma }_\mathrm{\Delta }^2}}p_\pi (M^2)`$ (31)
$`p_\pi (M^2)`$ $`=`$ $`{\displaystyle \frac{W^2+m_\pi ^2M^2}{2W}}`$
where the integration over running $`\mathrm{\Delta }`$ mass squared $`M^2`$ takes into account the unstable character of $`\mathrm{\Delta }`$ particle and the quantity $`\sqrt{\mathrm{\Gamma }_{\lambda _\mathrm{\Delta }}}`$ contains the above mentioned W evolution according to (30).
The total $`N^{}`$ decay width ($`\mathrm{\Gamma }_{res}(W)`$ in (27)) was assumed to be a sum over all partial widths presented in . The W evolution of each partial width $`\mathrm{\Gamma }_i(W)`$ was evaluated again based on a barrier penetration ansatz:
$$\mathrm{\Gamma }_i(W)=\mathrm{\Gamma }_i(W=M_{res})\frac{M_{res}}{W}\frac{J_l^2(p_{\pi _R}^{}R)+N_l^2(p_{\pi _R}^{}R)}{J_l^2(p_\pi ^{}R)+N_l^2(p_\pi ^{}R)}$$
(32)
where $`p_\pi ^{}`$ and $`p_{\pi _R}^{}`$ are three-momenta moduli of meson at running W and at resonance point respectively.
The total amplitude of $`\gamma p\mathrm{\Delta }^{++}\pi ^{}`$ process was then evaluated as a coherent superposition of resonant and total Born amplitude $`f_{\lambda _\gamma \lambda _p\lambda _\mathrm{\Delta }}^{B^{corr}}`$ , corrected for ingoing and outgoing channel absorption:
$$\pi \lambda _\mathrm{\Delta }\left|T\right|\lambda _\gamma \lambda _p=\pi \lambda _\mathrm{\Delta }\left|T_{res}\right|\lambda _\gamma \lambda _p+f_{\lambda _\gamma \lambda _p\lambda _\mathrm{\Delta }}^{B^{corr}}$$
(33)
## 4 Results and discussion.
Using the above described approach we performed a cross section calculation for reaction (1). The resonances included in the evaluation are listed in Table I. The relative contribution of a particular $`N^{}`$ in the cross section can be described by the factor of merit $`\sqrt{\mathrm{\Gamma }_{\gamma p}\mathrm{\Gamma }_{\mathrm{\Delta }\pi }}/\mathrm{\Gamma }_{tot}`$<sup>1</sup> <sup>1</sup><sup>1</sup>footnotetext: Proportional to peak amplitude for Breit-Wigner curve. also reported in Table I. All three- and four-star resonances with $`\sqrt{\mathrm{\Gamma }_{\gamma p}\mathrm{\Gamma }_{\mathrm{\Delta }\pi }}/\mathrm{\Gamma }_{tot}>0.1\%`$ were included. As follows from Table I, $`F_{15}(1680)`$, $`D_{33}(1700)`$ and $`F_{37}(1950)`$ resonances give maximum contribution and reaction (1) presents a promising opportunity for investigation of their structure.
As mentioned in sect. 3.1, we adopted two ways for $`\mathrm{\Delta }`$-in-flight term evaluation: a) from gauge invariance requirements; b) neglecting this term due to proximity of $`\mathrm{\Delta }`$ mass and $`\mathrm{\Lambda }`$ cut-off parameter.
Neglecting $`\mathrm{\Delta }`$-in-flight term we obtained a better data description at high W value and pion emission angles. The results are presented in Fig. 6 by solid lines. The data are reasonably reproduced in the overall W region. A remarkable point is that our calculations predict cross sections at $`\theta _\pi ^{}>120^0`$ and $`W>1.8`$ GeV lower than 1.2 $`\mu `$b/sr; such strong cross section suppression namely results from absorption in the initial and final states due to interactions with open hadronic channels. For W = 1.62 GeV the calculated cross sections are a bit below the data points; this discrepancy could be ascribed to simultaneous uncertainties in the cut-off parameter appearing in the strong vertex function used for the Born terms, in the $`\pi ^{}\mathrm{\Delta }^{++}\pi ^{}\mathrm{\Delta }^{++}`$ , $`\rho ^0p\rho ^0p`$ elastic amplitudes appearing in the absorption parameterisation, as well as to uncertainties in some s-channel $`N^{}`$ electromagnetic and strong parameters. The Born term contribution in angular distributions for reaction (1) is shown by dashed lines in Fig. 6, therefore the maximum $`N^{}`$ contribution takes place for CM pion emission angle above $`90^0`$. The evaluation of angular distributions performed by varying $`N^{}`$ strong decay couplings inside uncertainties of analysis demonstrated that the cross section variation does not exceed a few percent, being negligible compared with data uncertainties: therefore, our approach for the non-resonant background appears to be quite stable with respect to $`N^{}`$ strong decay parameter variation; this is an important feature for the extraction of $`N^{}`$ electromagnetic form factors from measured cross sections.
The comparison between the calculated total cross section for reaction (1) with recent SAPHIR data as well as with old ABBHHM Collaboration data is presented in Fig. 7. Decomposition of the total cross-section for reaction (1) in resonant, non-resonant processes and interference terms is also shown in Fig. 7. The maximum contribution of $`N^{}`$ (at level $`2030\%`$) is found at $`W<1.6`$ GeV and decreases as W increases. The region $`W<1.6`$ GeV also corresponds to maximum contribution of interference effects in coincidence with results presented in .
To demonstrate the effects of $`\pi N\mathrm{\Delta }`$ form factor implementation as well as ISI and FSI absorptive corrections, we also reported in Fig. 8 calculation results assuming: a) $`\pi N\mathrm{\Delta }`$ form factor equal to unity and absence of absorptive corrections (dashed line in Fig. 8); b) $`\pi N\mathrm{\Delta }`$ form factor from and absence of absorptive corrections (dotted line in Fig. 8); c) complete model (solid line). Calculation a) is not able to reproduce data at all: at $`W>1.6`$ GeV the difference between calculated and measured cross-section marks the complete failure of the pure Born terms calculation. Implementation of $`\pi N\mathrm{\Delta }`$ strong form factor turns out to be very important (dotted line of Fig. 8): we stress that such strong vertex function was taken from N – N scattering analysis without further tuning and assumed to be appropriate to describe the $`\pi N\mathrm{\Delta }`$ vertex in the t-channel pion exchange diagram, while for the contact term which is of similar magnitude as the pion exchange, we used gauge invariance as guidance for the vertex functions evaluation. However without initial and final state absorptive corrections it is not possible to reproduce data at W above 1.7 GeV. On the other hand, even neglecting ISI and FSI effects our approach is able to describe reasonably the data at $`W<1.6`$ GeV. Implementation of ISI and FSI absorptive correction factors provides a reasonable agreement between measured and calculated cross sections for W below 1.8 GeV, while for W above 1.9 GeV calculated cross sections are systematically higher than data. The reason for such deviation could be connected to the description of $`\pi \mathrm{\Delta }`$, $`\rho p`$ elastic scattering amplitudes in ISI-FSI: this aspect could be improved by adding contributions from isobar states with higher mass and spin; another explanation could be that for W above 1.9 GeV we may start to observe a transition from the corrected Born terms picture to other processes, as those discussed in ; in this case, the framework for non-resonant reaction mechanisms in $`\gamma p\pi ^{}\mathrm{\Delta }^{++}`$ reaction should be modified at W above 1.9 GeV. We discuss these aspects in the next paragraph, to examine the limits of applicability of our meson-baryon approach and study how to continue the description of reaction (1) in the higher energy region.
Following the recipe of , we replaced the pion exchange amplitude in Born term (Fig. 1c) by a pion Regge trajectory exchange. The results are presented in Fig. 9 (dotted line) and compared with the non-reggeized pion exchange (solid line). Regge trajectory exchange implementation provides a better data description at high W values. However, even in this case the calculated cross section appears to be systematically higher than the data. The possible reason could be the need to modify also the contact term in the way proposed in , thereby restoring gauge invariance. We found that such contact term modification led to a sizeable cross section reduction at W below 1.6 GeV, where descriptions in the frame of pion and Regge trajectory exchange have to coincide. To provide this coincidence we had to require $`N\mathrm{\Delta }`$ – (pion Regge trajectory) effective coupling to be a factor 1.2 higher than $`\pi N\mathrm{\Delta }`$ coupling. The cross section calculated under this assumption is shown in Fig. 9 with a dashed line and reproduces reasonably the data up to 2.5 GeV. As mentioned above, non – resonant part for $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes in ISI-FSI was estimated from $`\pi N`$ scattering partial wave data fit, containing data up to W $`=`$ 2.1 GeV. This is the main reason to restrict our calculation to W $`=`$ 2.5 GeV (3 GeV photon energy in lab. frame). Implementation of any data on $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes above 2 GeV would allow a model extension toward higher W.
To investigate the influence of $`\gamma pN^{}`$ vertex dressing on our cross section, we reported in Fig. 10 the calculations with $`N^{}`$ electromagnetic helicity amplitudes taken from PDG and from quark models. In the PDG case, we assumed the effects of $`\gamma pN^{}`$ vertex dressing at level of meson – baryon degrees of freedom to be included in the numbers quoted, while in the quark models case we assumed such effects not to be accounted for (”bare” $`\gamma pN^{}`$ vertices). Therefore, using the PDG amplitudes we applied our prescription (24) for removal of $`N^{}`$ dressing in ISI-FSI, while in cross section calculations with $`\gamma pN^{}`$ vertices from quark models we kept $`N^{}`$ excitation terms in ISI & FSI mechanisms with no modification. In principle, one would expect these two results approximately to coincide, assuming that the same dressing effects are introduced one way or another. Our results show that this is not the case, being the cross section calculated using the ”dressed” PDG values systematically higher than the quark model results; moreover, if we removed completely the $`N^{}`$ from the ISI-FSI description, the discrepancy would be even bigger; we also found that the difference between the two quark model cross sections is significantly smaller than the deviation between quark model and PDG results. Actually, the difference between the cross section evaluated with ”bare” and ”dressed” vertices could be due to specific approximations of quark model approaches; or it could be an indication that additional dressing effects are present in the experimentally extracted photocouplings from . This difference does not exceed 30 %, with a maximum at W between 1.5 – 1.7 GeV. This W range actually corresponds to maximum $`N^{}`$ contribution. The discrepancy appears to be negligible at W above 2.0 GeV and below 1.4 GeV, where $`N^{}`$ contributions are less pronunced.
To see the influence of particular dressing effects – $`N^{}`$ excitation in $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering processes shown in Fig. 4c – we performed calculations using ”dressed” $`N^{}`$ electromagnetic vertices, but keeping the $`N^{}`$ excitation in ISI & FSI mechanisms for Born terms. The results are shown in Fig. 10 by dotted line; of course in this case some double counting of the $`N^{}`$ electromagnetic vertex dressing takes place, but surprisingly the cross section seems to be in better agreement with the ”bare” quark model results. In any case, the difference between solid and dotted curves could be considered as estimation of this particular $`\gamma pN^{}`$ vertex dressing contribution in the cross section. This contribution is lower than 20 % and vanishes at W above 2.0 GeV, where $`N^{}`$ excitation is negligible.
To extract $`N^{}`$ helicity amplitudes for $`Q^2>0`$ it is important to have a good description of the $`Q^2`$ dependence for non-resonant processes. We checked this point by comparing the total cross section for reaction (1) calculated in our approach with data reported in . For the resonant part we used results from a Single Quark Transition Model (SQTM) fit about the $`Q^2`$ evolution of $`N^{}`$ photocouplings. The comparison between data and our calculations is presented in Fig. 11 where dashed lines correspond to the contribution of non – resonant processes only, while the complete evaluation with $`N^{}`$ included is shown by solid lines. Considering this comparison, two important points must be stressed: first, bins for experimental data, both in W and $`Q^2`$ were very wide, implying a strong cross section averaging over the measured kinematical range; second, experimental data about high-lying $`N^{}`$ photocouplings are rather scarce and do not allow to check the validity of SQTM predictions. However, our approach reproduces the measured cross-section for W in the 1.3-1.5 GeV and in the 1.5-1.7 GeV interval; for W in the 1.7-2.0 GeV bin our calculated cross-section is higher than the measurements; the reasons for such deviation could be namely the poor knowledge of the $`N^{}`$ contributions or the kinematical average effects, as well as the above mentioned limitations in the treatment of Born terms absorption in the higher W region. New high precision data upcoming from new facilities like TJNAF are therefore definitely needed for a better understanding of the $`Q^2`$ behaviour of $`N^{}`$ electromagnetic form factors. An important prediction of our evaluations is a strong $`N^{}`$ contribution in the W = 1.5 – 1.7 GeV interval, representing over 60 % of cross section at $`Q^21.0`$ GeV<sup>2</sup>. Therefore the two pion exclusive channel seems to present a promising opportunity for $`N^{}`$ structure investigation by photons with high virtuality.
To investigate $`N^{}`$ electromagnetic vertex dressing effects on the cross section at $`Q^2>0`$, we performed the calculation with $`N^{}`$ electromagnetic form factors taken from different approaches. Form factors in approach were determined from exclusive single pion production data analysis, imposing symmetry relations between form factors of $`N^{}`$’s belonging to a particular SU(6) multiplet: therefore they can be assumed to represent ”dressed” vertices. Instead, quark models should not contain higher order corrections and we considered their results as representing ”bare” vertices. Again we expected to have approximate coincidence of results obtained using these different ingredients. Actually in this case, according to Fig. 12, the difference between cross sections calculated with the two ”bare” vertices and is higher than the difference between results obtained using form factors from and quark models results. Also in this case we calculated the cross section using ”dressed” vertices from , with and without prescription (24) for $`N^{}`$ removal in ISI-FSI (solid and dotted curve in Fig. 12, respectively); in the latter case the result is supposed to be affected by double counting, but the difference can give an indication of the extent of dressing introduced by our effective unitarity implementation; as shown in Fig. 12, such effect appears to be negligible, thereby indicating that, for the $`Q^2>0`$ evolution in our framework, quark model theoretical uncertainties could be more important than dressing effects. Therefore, comparison of cross sections calculated in our approach with forthcoming precise experimental data at $`Q^2>0`$ can provide a promising opportunity to select between model approaches for $`N^{}`$ structure description in non – perturbative QCD region.
As mentioned above, we also compared the cross section calculated using ”dressed” $`N^{}`$ form factors with/without implementation of $`N^{}`$ excitation in ISI & FSI mechanisms (dotted and solid lines in Fig. 12). The contribution of this particular dressing mechanism is negligible for W$``$ 1.4 and 1.85 GeV and vanishes at $`Q^2>0.5`$ GeV<sup>2</sup> for W $`1.6`$ GeV. There are actually reasons for such behaviour: a) the relative contribution of non – resonant processes drastically falls down as $`Q^2`$ increases (from more than 80 % at the photon point to lower than 50 % at $`Q^2`$ above 1 GeV<sup>2</sup>; b) the ISI effects also go down as $`Q^2`$ increases due to suppression of transitions between photon and vector meson.
## 5 Conclusions.
We have extensively studied a phenomenological model for the two pion photo- and electroproduction through the intermediate $`\mathrm{\Delta }^{++}\pi ^{}`$ channel, as part of a broader effort in establishing a basis for analysis and interpretation of the upcoming data from TJNAF, where this as well as other exclusive electromagnetic production channels will be studied with unprecedented accuracy.
Our calculation included a minimal set of Born amplitudes, with appropriate absorptive corrections to effectively take into account interaction with open channels in the initial and final states, plus a large number of nucleon resonances that, according to the existing data, are thought to give a sizeable contribution to the cross section.
A strong form factor for $`\pi N\mathrm{\Delta }`$ vertex was introduced according to NN scattering experiments. Data about pion and proton electromagnetic form factors were also used to evaluate the pion-in-flight term and the s-channel nucleon pole term, respectively, for $`Q^2>0`$, while the $`Q^2`$ behaviour of the contact term was obtained imposing gauge invariance; the delta-in-flight term was calculated both from gauge invariance and using a $`\pi N\mathrm{\Delta }`$ vertex determined from NN scattering analysis. Strong and electromagnetic $`N^{}`$ couplings were related with experimental observables: electromagnetic helicity couplings $`A_{1/2}`$, $`A_{3/2}`$, $`C_{1/2}`$ and partial hadronic decay widths $`\mathrm{\Gamma }_{ls}`$.
We found that strong absorptive corrections were essential to reproduce experimental data for CM pion emission angle above 30<sup>0</sup> at W above 1.6 GeV. We developed a specific approach for ISI and FSI absorptive corrections, relating absorptive factors in ingoing and outgoing channels with $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering amplitudes. These elastic strong amplitudes were evaluated in a simple isobaric model, so that our approach did not have any free parameters to be determined from reaction (1). Any other approach for $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic amplitude evaluation can be easily implemented in our calculation. Reasonable reproduction of pion angular distributions in the photon point as well as good agreement between measured and calculated total cross sections for $`Q^2>0`$ demonstrate the ability of our approach to reproduce the main features of $`\pi ^{}\mathrm{\Delta }^{++}`$ production by real and virtual photons in the resonance region. Moreover the proposed method for ISI and FSI description can be considered as a reasonable phenomenological way to describe non-resonant processes in electromagnetic meson production for $`W>1.6`$ GeV, where the competition of many open hadronic channels makes a rigorous background evaluation very difficult. Our unitarity corrections effectively introduce a dressing of the resonance amplitudes; therefore we elaborated an empirical prescription to remove the dressing from ISI-FSI when using experimental resonance photocouplings, assumed to be already “dressed”; we also performed calculations without this prescription when using resonance amplitudes from quark models, asssumed to be “bare”. In the real photon case, we found that the calculations performed with two different quark models show a good agreement, while the results obtained using the experimental PDG amplitudes show a systematic disagreement with the previous ones; the reason is unclear and could be connected to other dressing effects than $`N^{}`$ excitation in $`\rho p\rho ^{^{}}p^{^{}}`$ $`\pi ^{^{}}\mathrm{\Delta }^{++^{}}\pi ^{}\mathrm{\Delta }^{++}`$ transitions as well as to quark model uncertainties and also to model-dependence in the experimental resonance amplitude extraction. We performed the same analysis for the virtual photon cross section, but in this case we found a disagreement between the two quark models adopted, while the calculation based on resonance amplitudes extracted from an analysis of experimental data showed less sensitivity to our “dressing” effects: the quark models discrepancies appeared to be more important, leading us to believe that in the virtual photon case the sensitivity to different quark model ingredients could be more pronunced, although the reason for the different behaviour of real and virtual photon is not evident.
The calculation presented could serve as a first basis for interpreting the data coming from new facilities like Jefferson Lab and moreover, allowing the resonance parameters to vary, it could be the starting point of a fitting procedure with the goal of investigating electromagnetic form factors for high mass $`N^{}`$’s ($`M_{res}>1.5`$ GeV), as well as for attempting to discover new states. The model presented is currently being used as a foundation in our development of a full three-body final state description for the two-pion production off the nucleon by real and virtual photons.
Acknowledgements Our special thanks to Dr. V. Burkert from Jefferson Laboratory, for continuos interest, useful discussions and support. Particular thanks to Prof. S. Dytman from the University of Pittsburgh, for kindly providing the elastic hadronic amplitudes from his global fit and to Prof. M. Giannini from the University of Genova, Italy, for kindly providing the $`N^{}`$ electromagnetic transition amplitudes from his quark model. We also want to thank Prof. N.C. Mukhopadhyay from Rensselaer Polytechnic Institute, for useful discussions and suggestions.
Figure captions
Fig 1. Tree-level diagrams for the $`\pi ^{}\mathrm{\Delta }^{++}`$ electromagnetic production on proton.
Fig2. ISI and FSI mechanisms.
Fig3. Description of $`\pi \mathrm{\Delta }`$ and $`\rho N`$ elastic scattering amplitudes related to ISI and FSI.
Fig 4. The effective $`\gamma pN^{}`$ vertex (a), the bare $`\gamma pN^{}`$ vertex (b), the dressing vertex correction, containing $`N^{}`$ excitation in $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho p`$ elastic scattering.
Fig 5a. $`\pi N`$ partial wave cross section from decomposed in resonant (dotted lines) background (dashed lines) parts according to procedure described in sect. 3.5. The fit results are shown by solid lines.
Fig 5b. (continued)
Fig 5c. (continued)
Fig 6. The calculated and measured angular distributions for $`\gamma p\mathrm{\Delta }^{++}\pi ^{}`$ reaction assuming negligible contribution for $`\mathrm{\Delta }`$-in- flight term (see text sect. 3). Dashed lines correspond to Born terms alone, while solid lines represent the complete calculation with $`N^{}`$ contribution included.
Fig 7. Decomposition of $`\gamma p\mathrm{\Delta }^{++}\pi ^{}`$ reaction cross section in Born terms (dashed line), $`N^{}`$ terms (dotted line) and interference term (dash-dotted line) contributions. Data are from (squares) and from (circles and triangles).
Fig 8. Total cross section for $`\gamma p\mathrm{\Delta }^{++}\pi ^{}`$ reaction at the photon point. Dashed line represents calculation results with no $`\pi N\mathrm{\Delta }`$ form factor and no ISI and FSI absorptive corrections. Dotted line corresponds to $`\pi N\mathrm{\Delta }`$ strong form factor taken from and no ISI and FSI absorptive corrections. Solid line represents cross section evaluation in the complete model. Data as in Fig. 8.
Fig 9. The calculation with pion Regge trajectory exchange. Solid line represents calculations performed with Born terms evaluated in single pion exchange picture. Dotted line represents the results obtained after substitution of pion exchange by pion Regge trajectory exchange, dashed line corresponds to an additional modification of contact term according to prescription of . Data as in Fig. 8.
Fig 10. The influence of $`\gamma pN^{}`$ vertex dressing effects at the photon point. Our calculations with PDG $`\gamma pN^{}`$ vertices (“dressed”) and $`N^{}`$ excitations excluded (solid line) and included (dotted line) in ISI & FSI treatment for Born terms are presented. The calculation results with $`\gamma pN^{}`$ vertices calculated in quark models (”bare” vertices) and with $`N^{}`$ excitation included in ISI & FSI are shown by: dashed line for model , dot – dashed line for model . Data as in Fig. 8.
Fig 11. $`Q^2`$ dependence of $`\gamma _vp\mathrm{\Delta }^{++}\pi ^{}`$ total virtual photon cross section in comparison with data from .
Fig 12. $`Q^2`$ dependence of $`\gamma pN^{}`$ vertex dressing effects. Our calculations with $`\gamma pN^{}`$ vertices from (“dressed”) and $`N^{}`$ excitations excluded (solid line) and included (dotted line) in ISI & FSI treatment are presented. The calculation results with $`\gamma pN^{}`$ vertices from quark models (”bare” vertices) and with $`N^{}`$ excitation included in ISI & FSI treatment are shown by: dashed line for model , dot – dashed line for model . Data from .
Table captions
Table 1. List of resonances included in our calculation.
Table 2. List of the resonances predicted by quark models that are weakly coupled to the $`N\pi `$ channel but should be strongly coupled to the $`N\pi \pi `$ channels.
Table 3. Nucleon resonances contributing in $`\pi ^{}\mathrm{\Delta }^{++}`$ and $`\rho ^0p`$ elastic scattering amplitudes. |
warning/0001/hep-th0001096.html | ar5iv | text | # 1 Introduction
## 1 Introduction
The essential feature of scalar tensor theories, such as Brans-Dicke theory, is to generalize general relativity and bring it into accord with Mach’s principle (the origin of physical properties of space is in the matter contained therein ). These theories are not completely geometrical since the gravitational effects are described by a scalar field as well as a metric tensor. In fact, the global distribution of matter affects the local gravitational properties through the emergence of a scalar field. The implementation of this interrelation between global and local properties of matter in quantum field theory, as demanded by Mach’s principle, is a complicated problem. Some ideas in this direction can be found in .
In a simplified picture one can expect that the role of a scalar tensor theory may be of importance for improving our knowledge on the local properties of a linear quantum field propagating in a gravitational background, in particular the local properties of the quantum stress-tensor induced by the two-point function of the quantum field. The present paper deals with the consideration of this issue. In specific terms, we study a model in which the local properties of a linear quantum field conformally coupled to a gravitational background is affected both by the local geometry and a conformal invariant scalar field derived from the state (boundary)-dependent part of the two-point function. To arrive at this model we basically take into account the local constraints imposed on the two-point function by the Hadamard state condition. In this context there is a problem concerning the specification of the state-dependent part of the two-point function. In our presentation we establish a connection between this problem and the problem of the determination of a conformal frame .
To avoid any confusion at the outset, we should note that the scalar tensor theory we wish to consider is meant only to provide an analytical mean to determine the general properties of a quantum stress-tensor that can consistently be coupled to conformally related background metrics, and in this respect its interpretation differs from the standard interpretation of such theories as alternative theories of gravitation.
The organization of this paper is as follows: In section 2 we present the Hadamard prescription and review the derivation of the local constraints on the state-dependent part of two-point function of a linear scalar quantum field conformally coupled to gravity. In section 3, we present a way to use a conformally invariant scalar field for analyzing the state-dependent part of the two-point function. It is shown that the implications of the resulting scalar tensor theory for the stress-tensor are in accord with the standard predictions of the renormalization theory. In section 4, we make some general remarks on the existence of an alternative frame in which the trace of the stress-tensor is determined by a cosmological constant rather than the usual anomalous trace. The existence of this frame indicates that the state-dependent part of the two-point functions may have some large scale characteristics which are basically not present in the conformal frame determined by the local characteristics. Similar arguments were discussed previously in a different context .
## 2 Hadamard state condition
We consider a free scalar quantum field $`\varphi (x)`$ propagating in a curved background spacetime with the action functional (We use the conventions of Hawking and Ellis for the signature and the sign of curvature)
$$S[\varphi ]=\frac{1}{2}d^4xg^{1/2}(g_{\alpha \beta }^\alpha \varphi ^\beta \varphi +\xi R\varphi ^2+m^2\varphi ^2)$$
(1)
where $`m`$ and $`\xi `$ are parameters, and $`R`$ is the scalar curvature (In the following the semicolon and $``$ indicate covariant differentiation). This gives rise to the field equation
$$(\mathrm{}m^2\xi R)\varphi (x)=0$$
(2)
The choice of the parameters $`m`$ and $`\xi `$ depends on the particular type of coupling. For example, the minimal coupling corresponds to $`m`$=0, $`\xi =0`$ and the conformal coupling (in four dimensions) corresponds to $`m`$=0, $`\xi =1/6`$ . A state of $`\varphi (x)`$ is characterized by a hierarchy of Wightman-functions (n-point functions)
$$\varphi (x_1),\mathrm{},\varphi (x_n)$$
(3)
We are primarily interested in those states which reflect the intuitive notion of a ”vacuum”. For this aim, we may restrict ourselves basically to quasi-free states, i.e. states for which the truncated n-point functions vanish for $`n>2`$ (In a linear theory this property is shared by the vacuum state of Minkowski space). Such states may be characterized by their two-point functions. In a linear theory the antisymmetric part of the two-point function is common to all states in the same representation. It is just the universal commutator function. Thus, in our case all the relevant informations about the state-dependent part of the two-point function are encoded in its symmetric part, denoted in the following by $`G^+(x,x^{})`$, which satisfies Eq.(2) in each argument. Equivalence principle suggests that the leading singularity of $`G^+(x,x^{})`$ should have a close correspondence to the singularity structure of the two-point function of a free massless field in Minkowski space. In general the entire singularity of $`G^+(x,x^{})`$ may have a more complicated structure. Usually one assumes that $`G^+(x,x^{})`$ has a singular structure represented by the Hadamard expansions. This means that in a normal neighborhood of a point $`x`$ the function $`G^+(x,x^{})`$ can be written as
$$G^+(x,x^{})=\frac{1}{8\pi ^2}\{\frac{\mathrm{\Delta }^{1/2}(x,x^{})}{\sigma (x,x^{})}+V(x,x^{})\mathrm{ln}\sigma (x,x^{})+W(x,x^{})\}$$
(4)
where $`2\sigma (x,x^{})`$ is the square of the distance along the geodesic joining $`x`$ and $`x^{}`$ and $`\mathrm{\Delta }(x,x^{})`$ is the Van Vleck determinant
$$\begin{array}{cc}\mathrm{\Delta }(x,x^{})=g^{1/2}(x)Det\{\sigma _{;\mu \nu _{^{}}}\}g^{1/2}(x^{})\hfill & \\ g(x)=Detg_{\alpha \beta }\hfill & \end{array}$$
(5)
The functions $`V(x,x^{})`$ and $`W(x,x^{})`$ are regular and have the following representations as power series
$$V(x,x^{})=\underset{n=0}{\overset{+\mathrm{}}{}}V_n(x,x^{})\sigma ^n$$
(6)
$$W(x,x^{})=\underset{n=0}{\overset{+\mathrm{}}{}}W_n(x,x^{})\sigma ^n$$
(7)
in which the coefficients are determined by applying Eq.(2) to $`G^+(x,x^{})`$, yielding the recursion relations
$$(n+1)(n+2)V_{n+1}+(n+1)V_{n+1;\alpha }\sigma ^{;\alpha }(n+1)V_{n+1}\mathrm{\Delta }^{1/2}\mathrm{\Delta }_{;\alpha }^{1/2}\sigma ^{;\alpha }+\frac{1}{2}(\mathrm{}m^2\xi R)V_n=0$$
(8)
$$(n+1)(n+2)W_{n+1}+(n+1)W_{n+1;\alpha }\sigma ^{;\alpha }(n+1)W_{n+1}\mathrm{\Delta }^{1/2}\mathrm{\Delta }_{;\alpha }^{1/2}\sigma ^{;\alpha }+\frac{1}{2}(\mathrm{}m^2\xi R)W_n$$
$$+(2n+3)V_{n+1}+V_{n+1;\alpha }\sigma ^{;\alpha }V_{n+1}\mathrm{\Delta }^{1/2}\mathrm{\Delta }_{;\alpha }^{1/2}\sigma ^{;\alpha }=0$$
(9)
together with the boundary condition
$$V_0+V_{0;\alpha }\sigma ^{;\alpha }V_0\mathrm{\Delta }^{1/2}\mathrm{\Delta }_{;\alpha }^{1/2}\sigma ^{;\alpha }+\frac{1}{2}(\mathrm{}m^2\xi R)\mathrm{\Delta }^{1/2}=0$$
(10)
From these relations one can determine the function $`V(x,x^{})`$ uniquely in terms of local geometry. Therefore it takes the same universal form for all states. But the biscalar $`W_0(x,x^{})`$ remains arbitrary. Its specification depends significantly on the choice of a state and may be regarded as the imposition of a boundary condition. However there is a general constraint on $`W_0(x,x^{})`$ which can be obtained from the symmetry condition of $`G^+(x,x^{})`$ together with the following dynamical equation which can be obtained using (2),(4) and (6)
$$(\mathrm{}m^2\xi R)W(x,x^{})=6v_1(x)+2v_{1;\alpha }\sigma ^{;\alpha }+0(\sigma )$$
(11)
where
$$v_1(x)=\underset{x^{}x}{lim}V_1(x,x^{})=\frac{1}{720}\{\mathrm{}RR_{\alpha \beta }R^{\alpha \beta }+R_{\alpha \beta \delta \gamma }R^{\alpha \beta \delta \gamma }\}$$
(12)
To get this constraint we first expand the symmetric function $`W(x,x^{})`$ into a covariant power series, namely
$$W(x,x^{})=W(x)\frac{1}{2}W_{;\alpha }(x)\sigma ^{;\alpha }+\frac{1}{2}W_{\alpha \beta }(x)\sigma ^{;\alpha }\sigma ^{;\beta }+\frac{1}{4}\{\frac{1}{6}W_{;\alpha \beta \gamma }(x)W_{\alpha \beta ;\gamma }(x)\}\sigma ^{;\alpha }\sigma ^{;\beta }\sigma ^{;\gamma }+0(\sigma ^2)$$
(13)
We may insert this into Eq.(11) and compare term by term up to the third order in $`\sigma ^{;\alpha }`$ to obtain
$$W_\gamma ^\gamma (x)=(\xi R+m^2)W(x)6v_1(x)$$
(14)
$$[W_{\alpha \beta }(x)\frac{1}{2}g_{\alpha \beta }W_\gamma ^\gamma (x)]^{;\alpha }=\frac{1}{4}(\mathrm{}W(x))_{;\beta }\frac{1}{2}m^2W_{;\beta }(x)+2v_1(x)_{;\beta }+\frac{1}{2}R_{\alpha \beta }W^{;\alpha }(x)\frac{1}{2}\xi RW_{;\beta }(x)$$
(15)
Then using the covariant expansion of the symmetric function $`W_0(x,x^{})`$
$$W_0(x,x^{})=W_0(x)\frac{1}{2}W_{0;\alpha }(x)\sigma ^{;\alpha }+\frac{1}{2}W_{0\alpha \beta }(x)\sigma ^{;\alpha }\sigma ^{;\beta }+0(\sigma ^{3/2})$$
(16)
together with Eqs.(7),(9) and (13), we get
$$W(x)=W_0(x)$$
(17)
$$W_{\alpha \beta }(x)=(W_{0\alpha \beta }(x)\frac{1}{4}g_{\alpha \beta }W_{0\gamma }^\gamma (x))+\frac{1}{4}[(m^2+\xi R)W_0(x)6v_1(x)]g_{\alpha \beta }$$
(18)
Substituting (17) and (18) into (15) leads to
$$[W_{0\alpha \beta }(x)\frac{1}{4}g_{\alpha \beta }W_{0\gamma }^\gamma (x)]^{;\alpha }=\frac{1}{2}v_{1;\beta }(x)+\frac{1}{4}(\mathrm{}W_0(x))_{;\beta }\frac{1}{4}m^2W_{0;\beta }(x)+\frac{1}{2}R_{\alpha \beta }W_0^{;\alpha }(x)$$
$$+\frac{1}{4}\xi [R_{;\beta }W_0(x)RW_{0;\beta }(x)]$$
(19)
This equation is a general constraint imposed on the state-dependent part of the two-point function. The function $`W_0(x)`$ may be considered as arbitrary, but once a specific assumption has been made on the form of $`W_0(x)`$, the equation (19) acts as a constraint on $`W_{0\alpha \beta }(x)`$.
We should note that the constraint (19) is, in principle, the first member of a hierarchy of constraints, because we have used the covariant expansion $`W_0(x,x^{})`$ only up to the second order in $`\sigma ^{;\alpha }`$. Thus, in general, there are some additional constraints on the higher order expansion terms. In our analysis we shall neglect these higher order constraints. Such a limitation is suggested by dimensional arguments because the second order expansion terms of $`W_0(x,x^{})`$ has already the physical dimension of a stress-tensor.
## 3 The conformally invariant scalar field
In the case of conformal coupling a local Hilbert space would in general exhibit an essential sensitivity to the pre-existing local causal structure of space-time which in the present case is determined by the conformal class of the background metric. By implication, this causal structure should act as the basic input for the characterization of the local states. Since the conformal transformations leave the causal structure unchanged we expect, in particular, that an essential ambiguity, related to conformal transformations, should enter the dynamical specification of the state-dependent part of the two-point function. Thus, in the case of the conformal coupling, it is suggestive to develop a dynamical model in which the conformal symmetry acts as a fundamental symmetry in the specification of the two-point function, in particular the function $`W_0(x)`$. In the following we shall use the constraint (19) to develop a dynamical model along this line. We first start with the explicit form of the constraint (19) in the case of conformal coupling, namely
$$[W_{0\alpha \beta }(x)\frac{1}{4}g_{\alpha \beta }W_{0\gamma }^\gamma (x)\frac{1}{2}g_{\alpha \beta }v_1(x)\frac{1}{4}g_{\alpha \beta }\mathrm{}W_0(x)]^{;\alpha }$$
$$=\frac{1}{2}R_{\alpha \beta }W_0^{;\alpha }(x)+\frac{1}{24}(R_{;\beta }W_0(x)RW_{0;\beta }(x)).$$
(20)
One can use the Bianchi identity
$$R_{\alpha \beta }^{;\alpha }=\frac{1}{2}R_{;\beta }$$
(21)
and the differential identity
$$\mathrm{}(W_{0;\beta }(x))=(\mathrm{}W_0(x))_{;\beta }+R_{\alpha \beta }W_0^{;\alpha }(x)$$
(22)
to show that (20) can be written as a total divergence
$$\mathrm{\Sigma }_{\alpha \beta }^{;\alpha }=0$$
(23)
where
$$\mathrm{\Sigma }_{\alpha \beta }=(W_{0\alpha \beta }(x)\frac{1}{4}g_{\alpha \beta }W_{0\gamma }^\gamma (x))\frac{1}{6}(R_{\alpha \beta }\frac{1}{4}Rg_{\alpha \beta })W_0(x)\frac{1}{3}(W_{0;\beta \alpha }(x)\frac{1}{4}g_{\alpha \beta }\mathrm{}W_0(x))$$
$$\frac{1}{2}g_{\alpha \beta }v_1(x)$$
(24)
Now, the basic input is to subject in (23) the choice of $`W_0(x)`$ to the condition
$$W_0(x)=\psi ^2(x)$$
(25)
where $`\psi (x)`$ is taken to be a conformally invariant scalar field coupled to the gravitational background, so that its dynamical equation is
$$(\mathrm{}\frac{1}{6}R)\psi =0.$$
(26)
For a given Hadamard state the field $`\psi `$ may be interpreted as measuring the one-point function of the quantum field $`\varphi `$. The conformal invariance of $`\psi `$ ensures that there exists no pre-assigned dynamical configuration for the one-point function in a local Hilbert space. This is indeed a desirable characteristic of a linear theory.
Technically, the merit of introducing the field $`\psi `$ is that the tensor $`\mathrm{\Sigma }_{\alpha \beta }+\frac{1}{2}g_{\alpha \beta }v_1(x)`$, which is traceless due to (24), may now be related to the conformal stress-tensor of $`\psi `$, namely
$$\mathrm{\Sigma }_{\alpha \beta }+\frac{1}{2}g_{\alpha \beta }v_1(x)=T_{\alpha \beta }[\psi ]$$
(27)
where the conformal stress-tensor $`T_{\alpha \beta }[\psi ]`$ is given by
$$T_{\alpha \beta }[\psi ]=(\frac{2}{3}_\alpha \psi _\beta \psi \frac{1}{6}g_{\alpha \beta }_\gamma \psi ^\gamma \psi )\frac{1}{3}(\psi _\alpha _\beta \psi g_{\alpha \beta }\psi \mathrm{}\psi )+\frac{1}{6}\psi ^2G_{\alpha \beta }$$
(28)
in which $`G_{\alpha \beta }`$ is the Einstein tensor. The tensor $`T_{\alpha \beta }`$ is traceless due to the dynamical equation (26). The meaning of the relation (27) is that it defines a formal prescription which allows us to relate the tensor $`W_{0\alpha \beta }(x)`$ in (24) to the function $`W_0(x)`$ and the metric tensor $`g_{\alpha \beta }`$, so it characterizes a criterion to select the class of admissible Hadamard states. Taking into account (28) we can write this criterion as
$$G_{\alpha \beta }3\psi ^2g_{\alpha \beta }v_1(x)=6\psi ^2(\mathrm{\Sigma }_{\alpha \beta }+\tau _{\alpha \beta }(\psi )).$$
(29)
Here $`\tau _{\alpha \beta }(\psi )`$, is equal to $`T_{\alpha \beta }[\psi ]`$ without the $`G_{\alpha \beta }`$-term, so it coincides up to a sign with the so called modified energy-momentum (stress-) tensor . Now, the basic strategy is to consider the tensor $`\mathrm{\Sigma }_{\alpha \beta }`$ as the quantum stress-tensor induced by the two-point function. Our criterion can then be interpreted as a rule for relating the latter tensor to the local background geometry, as reflected in (29). The essential point is that this rule is expressed in the form of a scalar tensor theory in which the dynamics of the scalar field $`\psi `$ makes substantially no distinction between different frames in the conformal class of the background metric. The implication is that at the dynamical level all conformal frames may be considered as equivalent.
This conformal invariance reflects a basic connection between the state-dependent part of the two-point function and the pre-existing causal structure determined by the background metric. In particular, it establishes a basic connection between the properties of a given physical state in a local Hilbert space and those of a corresponding conformal frame. To see this in explicit terms let us consider a conformal transformation
$$\overline{g}_{\alpha \beta }=\mathrm{\Omega }^2(x)g_{\alpha \beta }$$
$$\overline{\psi }(x)=\mathrm{\Omega }^1(x)\psi (x)$$
(30)
Due to (25), $`W_0(x)`$ would then transform as
$$\overline{W}_0(x)=\mathrm{\Omega }^2(x)W_0(x)$$
(31)
It is now clear from (31) that a given conformal frame may be characterized by the particular configuration of $`W_0(x)`$ (or alternatively $`\psi `$) in that frame. Therefore the problem of specification of $`W_0(x)`$ for a given physical state is basically connected with the problem of determination of a conformal frame. In particular, different states characterized by conformally related configurations of $`W_0(x)`$ should principally be supported on different conformally related metrics. The same conclusion holds for their stress-tensors.
At this point we make a general remark concerning the consistency of our results with the standard prediction of the renormalization theory. Focusing ourselves to the two-point function on the background metric we can take the trace of (27), to obtain
$$\mathrm{\Sigma }_\alpha ^\alpha =2v_1(x)$$
(32)
This together with (23) characterize the general properties of the quantum stress-tensor on the background metric. These properties are consistent with the well-known results of the renormalization theory and $`v_1(x)`$ is actually the function that determines what is commonly known as the trace anomaly. In our presentation this quantum anomaly requires a somewhat distinct behavior of the scalar field $`\psi `$. In fact, according to (23) and (27) and due to the nonvanishing trace anomaly, the tensor $`T_{\alpha \beta }[\psi ]`$, which may be considered as the stress-tensor of the field $`\psi `$, appears not to be conserved on the background metric, requiring the dynamical properties of $`\psi `$ on the background metric not to fit in with the properties of a diffeomorphism invariant action characterizing a C-number (classical) field. But it is necessary to stress that this behavior does not appear to be a physical contradiction in the present case. Actually, the scalar field $`\psi `$ which characterizes the local property of the two-point function may in general change its configuration if one varies the two-point function within a local Hilbert space. Therefore, in general it may not act as a C-number field within a local Hilbert space. By implication, the standard results of a diffeomorphism invariant action may not be applied to $`\psi `$. We note that a similar process of assigning non-diffeomorphism invariant properties to a local Hilbert space has been previously discussed in the context of generally covariant quantum field theory .
## 4 $`\mathrm{\Lambda }`$-frame
The conformal symmetry which was established in the local specification of $`W_0(x)`$ would imply that locally the stress-tensor $`\mathrm{\Sigma }_{\alpha \beta }`$ can be related to different conformal frames. Thus the question arises as to which frame should be considered as a physical frame. To deal with this question it is necessary to emphasize the role of the superselection rules which characterize the boundary conditions imposed on the physically realizable states and the corresponding Hilbert spaces. In general, the identification of a conformal frame as a physical frame depends on the particular superselection rule one wishes to apply. Of direct physical significance, in the present case, is a superselection rule that tells us how a local Hilbert space is linked to the large scale boundary conditions imposed on physical states. If the latter conditions correspond to the presence of large scale distribution of matter whose energy density is measured by a cosmological constant, one may subject the determination of a conformal frame (alternatively a local Hilbert space) to the asymptotic correspondence between the anomalous trace and a nonvanishing cosmological constant at sufficiently large spacelike distances. In general, this condition may not be realized in the underlying background frame, so in this case the physical frame is expected to be different from the background frame.
This observation opens a way to study the transition from the local characteristics of physical states in a local Hilbert space to the large scale characteristics, which is expected to be of particular importance for establishing the large scale gravitational coupling of physical states in a local Hilbert space. Since by such a transition the small distance properties are no more important, we may take the overall correspondence between the anomalous trace and a nonvanishing cosmological constant everywhere as the defining characteristic of a local conformal frame which, by implication, acts as the physical frame if one focuses on large scale characteristics of physical states in the presence of large scale distribution of matter. For the construction of this frame one needs only to apply a conformal transformation to the background frame which establishes the correspondence between the trace anomaly and a nonvanishing cosmological constant. Denoting the cosmological constant by $`\mathrm{\Lambda }`$, the corresponding conformal factor may be taken to satisfy the equation
$$3\mathrm{\Omega }^2(x)\psi ^2v_1(\mathrm{\Omega }^2(x)g_{\alpha \beta })=\mathrm{\Lambda }$$
(33)
Under this conformal transformation the equation (27) transforms to
$$\overline{\mathrm{\Sigma }}_{\alpha \beta }\frac{1}{6}\mathrm{\Lambda }\overline{g}_{\alpha \beta }\overline{\psi }^2=T_{\alpha \beta }[\overline{\psi }]$$
(34)
or, equivalently
$$G_{\alpha \beta }(\overline{g}_{\alpha \beta })+\mathrm{\Lambda }\overline{g}_{\alpha \beta }=6\overline{\psi }^2(\overline{\mathrm{\Sigma }}_{\alpha \beta }+\tau _{\alpha \beta }(\overline{\psi })).$$
(35)
Therefore in the new frame, which we call the $`\mathrm{\Lambda }`$-frame, a scalar tensor theory with a cosmological constant is obtained together with Eq.(33) which is a complicated constraint on the conformal factor. In the $`\mathrm{\Lambda }`$-frame, contrary to the background frame, the stress-tensor $`\overline{\mathrm{\Sigma }}_{\alpha \beta }`$ may not be conserved. However Eq.(34) implies that one can establish a conserved stress-tensor by replacing $`\overline{\psi }`$ by a constant average value $`\overline{\psi }`$=const. In this case the usual features of general relativity can be established in the $`\mathrm{\Lambda }`$-frame. In particular, the tensor $`T_{\alpha \beta }[\overline{\psi }]`$, which was found to be non-conserved in the back-ground frame, becomes a multiple of the Einstein tensor, so a conserved tensor.
For further investigation of the constraint (33), we write its explicit form on the background metric. Using the conformal transformation of the function $`v_1(x)`$ we find
$$e^{2\omega }\psi ^2\{3v_1(g_{\alpha \beta })+\frac{1}{240}[2R\mathrm{}\omega +2R_{;\alpha }\omega ^{;\alpha }+6\mathrm{}(\mathrm{}\omega )+8[(\mathrm{}\omega )^2$$
$$\omega _{;\alpha \beta }\omega ^{;\alpha \beta }R_{\alpha \beta }\omega ^{;\alpha }\omega ^{;\beta }\omega ^{;\gamma }\omega _{;\gamma }\mathrm{}\omega 2\omega _{;\alpha \beta }\omega ^{;\alpha }\omega ^{;\beta }]]\}=\mathrm{\Lambda }$$
(36)
where $`\omega =\mathrm{ln}\mathrm{\Omega }`$. As an illustration we shall now apply (36) to study an asymptotic relation between the $`\mathrm{\Lambda }`$-frame and a specific background metric which we take to be described by a Schwarzschild black hole. In this case the function $`v_1(x)`$, which determines the trace anomaly, reduces to
$$v_1(g_{\alpha \beta })=\frac{1}{720}R_{\alpha \beta \delta \gamma }R^{\alpha \beta \delta \gamma }=\frac{1}{15}\frac{M^2}{r^6}$$
(37)
where $`M`$ is the mass of the black hole. Since the trace anomaly vanishes for $`r\mathrm{}`$, one may generally expect that for a sufficiently small $`\mathrm{\Lambda }`$ there should be no distinction between the background and the $`\mathrm{\Lambda }`$-frame in a region far from the black hole event horizon. That this behavior is dynamically allowed follows from the equation (36) as we briefly demonstrate: Let us restrict ourselves to the static case and assume that $`\omega `$ is only a function of $`r`$. For $`r>>2M`$ the equation (36) takes then the form
$$\omega ^{^{\prime \prime \prime \prime }}4\omega ^{}_{}{}^{}2\omega {}_{}{}^{\prime \prime }=40\mathrm{\Lambda }\psi ^2e^{2\omega }8\frac{M^2}{r^6},r>>2M$$
(38)
where prime indicates differentiation with respect to $`r`$. This equation reveals that $`\omega `$ as a slowly varying function would be a solution for a large value of $`r`$ and a sufficiently small cosmological constant. In particular, for large values of $`r`$ an almost constant conformal factor (close to one) can be used to establish the correspondence between the background frame and the $`\mathrm{\Lambda }`$-frame.
## 5 Summary and outlook
For a quantum field conformally coupled to a gravitational background we have presented a model in which the role of a scalar tensor theory is emphasized for studying the local constraints imposed on physical states by the Hadamard state condition. The corresponding scalar field is conformally invariant and controls the coupling of the stress-tensor to the conformal class of the background metric. The predictions of this theory are in accord with the standard results of the stress-tensor renormalization if one chooses a conformal frame corresponding to the background metric. We have emphasized that the choice of a specific conformal frame as a physical frame must, in general, be subjected to the superselection rules regulating the coupling of a local Hilbert space to the physical conditions at distant regions. In this context we have discussed the possibility to consider the theory in a distinguished frame, namely the $`\mathrm{\Lambda }`$-frame, which may act as the physical frame for the establishing the large scale gravitational coupling of physical states in the presence of large scale distribution of matter. It is suggestive to link this large scale gravitational coupling of physical states, reflected in the $`\mathrm{\Lambda }`$-frame, with their cut-off property in the short distance scaling. A dynamical cut-off theory of this type, if properly formulated, would reflect one of the characteristic implication of Mach’s principle in quantum field theory. |
warning/0001/physics0001073.html | ar5iv | text | # 1 Introduction
## 1 Introduction
So far, extensive work has been performed on the linear stability analysis of collective motion in particle accelerators . Nonlinear theories of wave interaction and formation of patterns and coherent structures in intense beams are however less prevalent, in part, due to the mathematical complexity of the subject, but also because of the commonly spread opinion that highly nonlinear regime is associated with poor machine performance that is best to be avoided.
Nevertheless, nonlinear wave interaction is a well observed phenomenon , in present machines, complete and self-consistent theory explaining the processes, leading to the formation of self-organized structures mentioned above is far from being established. The present paper is aimed as an attempt in this direction.
The problem addressed here (perhaps, the simplest one) is the evolution of a beam in longitudinal direction under the influence of a resonator voltage induced by the beam itself. Linear theory is obviously unable to explain bunch (droplet) formation and bunch breakoff (especially in the highly damped regime), phenomena that have been observed by numerical simulations , , , but it should be considered as the first important step towards our final goal – nonlinear model of wave interaction developed in Section 3.
It is well-known that within the framework of linear stability analysis the solution of the original problem is represented as a superposition of plane waves with constant amplitudes, while the phases are determined by the spectrum of solutions to the dispersion equation. Moreover, the wave amplitudes are completely arbitrary and independent of the spatial and temporal variables. The effect of nonlinearities is to cause variation in the amplitudes in both space and time. We are interested in describing these variations, since they govern the relatively slow process of formation of self-organized patterns and coherent structures.
The importance of the linear theory is embedded in the dispersion relation and the type of solutions it possesses. If the dispersion relation has no imaginary parts (no dissipation of energy occurs and no pumping from external energy sources is available) and its solutions, that is the wave frequency as a function of the wave number are all real, then the corresponding amplitude equations describing the evolution of the wave envelopes will be of nonlinear Schrödinger type. Another possibility arises for conservative systems when some of the roots of the dispersion equation appear in complex conjugate pairs. Then the amplitude equations can be shown to be of the so called AB–type . For open systems (like the system studied here) the dispersion relation is in general a complex valued function of the wave frequency and wave number and therefore its solutions will be complex. It can be shown that the equation governing the slow evolution of the wave amplitudes in this case will be the Ginzburg–Landau equation.
Based on the renormalization group approach we have recently derived a Ginzburg–Landau equation for the amplitude of the resonator voltage in the case of a coasting beam . The derivation has been carried out under the assumption that the spatial evolution of the system is much slower compared to the temporal one. This restriction has been removed here, and the present paper may be considered as an extension of .
Using the method of multiple scales we derive a set of coupled amplitude equations for the slowly varying part of the longitudinal distribution function and for the intensity of a single resonator wave with an arbitrary wave number (and wave frequency, specified as a solution to the linear dispersion equation). The equation governing the evolution of the voltage envelope is shown to be of Ginzburg–Landau type.
## 2 Formulation of the Problem
It is well-known that the longitudinal dynamics of an individual beam particle is governed by the set of equations
$`{\displaystyle \frac{dz_1}{dt}}=k_0\mathrm{\Delta }E;{\displaystyle \frac{d\mathrm{\Delta }E}{dt}}={\displaystyle \frac{e\omega _sV_{RF}}{2\pi }}(\mathrm{sin}\varphi \mathrm{sin}\varphi _s)+{\displaystyle \frac{e\omega _s}{2\pi }}V_1,`$ (2.1)
where
$`k_0={\displaystyle \frac{\eta \omega _s}{\beta _s^2E_s}}`$ (2.2)
is the proportionality constant between the frequency deviation of a non synchronous particle with respect to the frequency $`\omega _s`$ of the synchronous one, and the energy deviation $`\mathrm{\Delta }E=EE_s`$. The quantity $`k_0`$ also involves the phase slip coefficient $`\eta =\alpha _M\gamma _s^2`$, where $`\alpha _M`$ is the momentum compaction factor . The variables
$`z_1=\theta \omega _st;\varphi =\varphi _shz_1.`$ (2.3)
are the azimuthal displacement of the particle with respect to the synchronous one, and the phase of the RF field, respectively. Here $`V_{RF}`$ is the amplitude of the RF voltage and $`h`$ is the harmonic number. Apart from the RF field we assume that beam motion is influenced by a resonator voltage $`V_1`$ due to a broad band impedance
$`{\displaystyle \frac{^2V_1}{z_1^2}}2\gamma {\displaystyle \frac{V_1}{z_1}}+\omega ^2V_1={\displaystyle \frac{2\gamma e}{\omega _s}}{\displaystyle \frac{I_1}{t}},`$ (2.4)
where
$`\omega ={\displaystyle \frac{\omega _r}{\omega _s}};\gamma ={\displaystyle \frac{\omega }{2Q}};I_1(\theta ;t)={\displaystyle }d\mathrm{\Delta }E(\omega _s+k_0\mathrm{\Delta }E)f_1(\theta ,\mathrm{\Delta }E;t),`$ (2.5)
$`f_1(\theta ,\mathrm{\Delta }E;t)`$ is the longitudinal distribution function, $`\omega _r`$ is the resonant frequency, $`Q`$ is the quality factor of the resonator and $``$ is the resonator shunt impedance.
It is convenient to pass to a new independent variable (“time”) $`\theta `$ and to the new dimensionless variables , :
$`\tau =\nu _s\theta ;z=z_1\sqrt{\nu _s};u={\displaystyle \frac{1}{\sqrt{\nu _s}}}{\displaystyle \frac{k_0\mathrm{\Delta }E}{\omega _s}},`$ (2.6)
$`f_1(\theta ,\mathrm{\Delta }E;t)={\displaystyle \frac{\rho _0\left|k_0\right|}{\omega _s\sqrt{\nu _s}}}f(z,u;\theta );V_1=\lambda _1V;I_1=\omega _s\rho _0I,`$ (2.7)
where
$`\nu _s^2={\displaystyle \frac{ehk_0V_{RF}\mathrm{cos}\varphi _s}{2\pi \omega _s}};\lambda _1=2\gamma _0e\omega _s\rho _0.`$ (2.8)
In the above expressions the quantity $`\rho _0`$ is the uniform beam density in the thermodynamic limit. The linearized equations of motion (2.1) and equation (2.4) in these variables read as:
$`{\displaystyle \frac{dz}{d\tau }}=u;{\displaystyle \frac{du}{d\tau }}=z+\lambda V,`$ (2.9)
$`{\displaystyle \frac{^2V}{z^2}}2\gamma _0{\displaystyle \frac{V}{z}}+\omega _0^2V={\displaystyle \frac{I}{z}};I(z;\theta )={\displaystyle }du(1+u\sqrt{\nu _s})f(z,u;\theta ),`$ (2.10)
where
$`\gamma _0={\displaystyle \frac{\gamma }{\sqrt{\nu _s}}};\omega _0={\displaystyle \frac{\omega }{\sqrt{\nu _s}}};\lambda ={\displaystyle \frac{e^2\gamma _0k_0\rho _0}{\pi \nu _s\sqrt{\nu _s}}}.`$ (2.11)
We can now write the Vlasov equation for the longitudinal distribution function $`f(z,u;\theta )`$, which combined with the equation for the resonator voltage $`V(z;\theta )`$
$`{\displaystyle \frac{f}{\tau }}+u{\displaystyle \frac{f}{z}}z{\displaystyle \frac{f}{u}}+\lambda V{\displaystyle \frac{f}{u}}=0,`$ (2.12)
$`{\displaystyle \frac{^2V}{z^2}}2\gamma _0{\displaystyle \frac{V}{z}}+\omega _0^2V={\displaystyle \frac{I}{z}},`$ (2.13)
$`I(z;\theta )={\displaystyle 𝑑u\left(1+u\sqrt{\nu _s}\right)f(z,u;\theta )},`$ (2.14)
comprises the starting point for our subsequent analysis.
## 3 Derivation of the Amplitude Equations for a Coasting Beam
In this Section we analyze the simplest case of a coasting beam. The model equations (2.12) and (2.13) acquire the form
$`{\displaystyle \frac{f}{\theta }}+u{\displaystyle \frac{f}{z}}+\lambda V{\displaystyle \frac{f}{u}}=0,`$ (3.1)
$`{\displaystyle \frac{^2V}{z^2}}2\gamma {\displaystyle \frac{V}{z}}+\omega ^2V={\displaystyle 𝑑u\left(\frac{f}{\theta }\frac{f}{z}\right)},`$ (3.2)
where the parameter $`\lambda `$ should be calculated for $`\nu _s=1`$. In what follows it will be convenient to write the above equations more compactly as:
$`\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}},u)f+\lambda V{\displaystyle \frac{f}{u}}=0,`$ (3.3)
$`\widehat{𝒱}({\displaystyle \frac{}{z}},\omega )V=\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}})f,`$ (3.4)
where we have introduced the linear operators
$`\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}},u)={\displaystyle \frac{}{\theta }}+u{\displaystyle \frac{}{z}},`$ (3.5)
$`\widehat{𝒱}({\displaystyle \frac{}{z}},\omega )={\displaystyle \frac{^2}{z^2}}2\gamma {\displaystyle \frac{}{z}}+\omega ^2,`$ (3.6)
$`\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}})={\displaystyle \frac{}{\theta }}{\displaystyle \frac{}{z}},`$ (3.7)
$`𝒢(z,u;\theta )={\displaystyle 𝑑u𝒢(z,u;\theta )}.`$ (3.8)
To obtain the desired amplitude equation for nonlinear waves we use the method of multiple scales , . The key point of this approach is to introduce slow temporal as well as spatial scales according to the relations:
$`\theta ;T_1=ϵ\theta ;T_2=ϵ^2\theta ;\mathrm{};T_n=ϵ^n\theta ;\mathrm{}`$ (3.9)
$`z;z_1=ϵz;z_2=ϵ^2z;\mathrm{};z_n=ϵ^nz;\mathrm{}`$ (3.10)
where $`ϵ`$ is a formal small parameter. Next is to utilize the perturbation expansion of the longitudinal distribution function $`f`$, the resonator voltage $`V`$
$`f=f_0\left(u\right)+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}ϵ^kf_k;V={\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}ϵ^kV_k,`$ (3.11)
and the operator expansions
$$\widehat{}(\frac{}{\theta }+\underset{k=1}{\overset{\mathrm{}}{}}ϵ^k\frac{}{T_k},\frac{}{z}+\underset{k=1}{\overset{\mathrm{}}{}}ϵ^k\frac{}{z_k},u)=$$
$`=\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}},u)+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}ϵ^k\widehat{}({\displaystyle \frac{}{T_k}},{\displaystyle \frac{}{z_k}},u),`$ (3.12)
$$\widehat{}(\frac{}{\theta }+\underset{k=1}{\overset{\mathrm{}}{}}ϵ^k\frac{}{T_k},\frac{}{z}+\underset{k=1}{\overset{\mathrm{}}{}}ϵ^k\frac{}{z_k})=$$
$`=\widehat{}({\displaystyle \frac{}{\theta }},{\displaystyle \frac{}{z}})+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}ϵ^k\widehat{}({\displaystyle \frac{}{T_k}},{\displaystyle \frac{}{z_k}}),`$ (3.13)
$`\widehat{𝒱}\left({\displaystyle \frac{}{z}}+{\displaystyle \underset{k=1}{\overset{\mathrm{}}{}}}ϵ^k{\displaystyle \frac{}{z_k}}\right)=\widehat{𝒱}+ϵ\widehat{𝒱}_z{\displaystyle \frac{}{z_1}}+{\displaystyle \frac{ϵ^2}{2}}\left(\widehat{𝒱}_{zz}{\displaystyle \frac{^2}{z_1^2}}+2\widehat{𝒱}_z{\displaystyle \frac{}{z_2}}\right)+\mathrm{}`$ (3.14)
where $`\widehat{𝒱}_z`$ implies differentiation with respect to $`/z`$. Substituting them back into (3.3) and (3.4) we obtain the corresponding perturbation equations order by order. It is worth noting that without loss of generality we can miss out the spatial scale $`z_2`$, because it can be transformed away by a simple change of the reference frame. For the sake of saving space we will omit the explicit substitution and subsequent calculations and state the final result order by order.
First order $`O(ϵ)`$:
$`\widehat{}f_1+\lambda V_1{\displaystyle \frac{f_0}{u}}=0,`$ (3.15)
$`\widehat{𝒱}V_1=\widehat{}f_1.`$ (3.16)
Second order $`O(ϵ^2)`$:
$`\widehat{}f_2+\lambda V_2{\displaystyle \frac{f_0}{u}}=\widehat{}_1f_1\lambda V_1{\displaystyle \frac{f_1}{u}},`$ (3.17)
$`\widehat{𝒱}V_2=\widehat{}f_2+\widehat{}_1f_1\widehat{𝒱}_z{\displaystyle \frac{V_1}{z_1}}.`$ (3.18)
Third order $`O(ϵ^3)`$:
$`\widehat{}f_3+\lambda V_3{\displaystyle \frac{f_0}{u}}=\widehat{}_1f_2\widehat{}_2f_1\lambda V_1{\displaystyle \frac{f_2}{u}}\lambda V_2{\displaystyle \frac{f_1}{u}},`$ (3.19)
$`\widehat{𝒱}V_3=\widehat{}f_3+\widehat{}_1f_2+\widehat{}_2f_1\widehat{𝒱}_z{\displaystyle \frac{V_2}{z_1}}{\displaystyle \frac{\widehat{𝒱}_{zz}}{2}}{\displaystyle \frac{^2V_1}{z_1^2}},`$ (3.20)
where $`\widehat{}_n`$ and $`\widehat{}_n`$ are the corresponding operators, calculated for $`T_n`$ and $`z_n`$.
In order to solve consistently the perturbation equations for each order we need a unique equation for one of the unknowns; it is more convenient to have a sole equation for the distribution functions $`f_n`$ alone. This will prove later to be very efficient for the removal of secular terms that appear in higher orders. By inspecting the above equations order by order one can catch their general form:
$`\widehat{}f_n+\lambda V_n{\displaystyle \frac{f_0}{u}}=\alpha _n;\widehat{𝒱}V_n=\widehat{}f_n+\beta _n,`$ (3.21)
where $`\alpha _n`$ and $`\beta _n`$ are known functions, determined from previous orders. Eliminating $`V_n`$ we obtain:
$`\widehat{𝒱}\widehat{}f_n+\lambda {\displaystyle \frac{f_0}{u}}\widehat{}f_n=\lambda {\displaystyle \frac{f_0}{u}}\beta _n+\widehat{𝒱}\alpha _n.`$ (3.22)
Let us now proceed with solving the perturbation equations. The analysis of the first order equations (linearized equations) is quite standard, and for the one-wave solution we readily obtain:
$`V_1=E(z_n;T_n)e^{i\phi }+E^{}(z_n;T_n)e^{i\phi ^{}},`$ (3.23)
$`f_1=\lambda {\displaystyle \frac{f_0}{u}}\left[{\displaystyle \frac{E(z_n;T_n)}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}e^{i\phi }+{\displaystyle \frac{E^{}(z_n;T_n)}{\stackrel{~}{}^{}(i\mathrm{\Omega },ik,u)}}e^{i\phi ^{}}\right]+F(z_n,u;T_n),`$ (3.24)
with
$`\phi =\mathrm{\Omega }\theta kz,`$ (3.25)
where given the wave number $`k`$, the wave frequency $`\mathrm{\Omega }(k)`$ is a solution to the dispersion equation:
$`\stackrel{~}{𝒟}(k,\mathrm{\Omega }(k))0.`$ (3.26)
The dispersion function $`\stackrel{~}{𝒟}(k,\mathrm{\Omega })`$ is proportional to the dielectric permittivity of the beam and is given by the expression
$`\stackrel{~}{𝒟}(k,\mathrm{\Omega })=\stackrel{~}{𝒱}\left(ik\right)+\lambda \stackrel{~}{}(i\mathrm{\Omega },ik){\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{f_0}{u}},`$ (3.27)
where
$`\widehat{}e^{i\phi }=\stackrel{~}{}(i\mathrm{\Omega },ik,u)e^{i\phi };\widehat{𝒱}e^{i\phi }=\stackrel{~}{𝒱}(ik)e^{i\phi };\widehat{}e^{i\phi }=\stackrel{~}{}(i\mathrm{\Omega },ik)e^{i\phi }.`$ (3.28)
Note that the wave frequency has the following symmetry property:
$`\mathrm{\Omega }^{}(k)=\mathrm{\Omega }(k).`$ (3.29)
The functions $`E(z_n;T_n)`$ and $`F(z_n,u;T_n)`$ in equations (3.23) and (3.24) are the amplitude function we wish to determine. Clearly, these functions are constants with respect to the fast scales, but to this end they are allowed to be generic functions of the slow ones.
In order to specify the dependence of the amplitude functions on the slow scales, that is to derive the desired amplitude equations one need to go beyond the first order. The first step is to evaluate the right hand side of equation (3.22) corresponding to the second order with the already found solution (3.23) and (3.24) for the first order. This yields terms (proportional to $`e^{i\phi }`$) belonging to the kernel of the linear operator on the left hand side of equation (3.22), which consequently give rise to the so called secular contributions to the perturbative solution. If the spectrum of solutions to the dispersion equation (3.26) is complex (as is in our case), terms proportional to $`e^{2Im(\mathrm{\Omega })\theta }`$ appear on the right hand side of (3.22). Since, the imaginary part of the wave frequency we consider small, the factor $`e^{2Im(\mathrm{\Omega })\theta }`$ is slowly varying in $`\theta `$ and we can replace it by $`e^{2Im(\mathrm{\Omega })T_n}`$, where the slow temporal scale $`T_n`$ is to be specified later. This in turn produces additional secular terms, which need to be taken care of as well. (Note that exactly for this purpose we have chosen two amplitude functions at first order). The procedure to avoid secular terms is to impose certain conditions on the amplitudes $`E(z_n;T_n)`$ and $`F(z_n,u;T_n)`$, that guarantee exact cancellation of all terms proportional to $`e^{i\phi }`$ and terms constant in the fast scales $`z`$ and $`\theta `$ (containing $`e^{2Im(\mathrm{\Omega })T_n}`$) on the right hand side of equation (3.22). One can easily check by direct calculation that the above mentioned conditions read as:
$`{\displaystyle \frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}}{\displaystyle \frac{E}{T_1}}{\displaystyle \frac{\stackrel{~}{𝒟}}{k}}{\displaystyle \frac{E}{z_1}}=i\lambda \stackrel{~}{}{\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{F}{u}}E,`$ (3.30)
$`\widehat{}_1F+2\lambda ^2Im(\mathrm{\Omega }){\displaystyle \frac{}{u}}\left({\displaystyle \frac{1}{\left|\stackrel{~}{}\right|^2}}{\displaystyle \frac{f_0}{u}}\right)\left|E\right|^2e^{2Im(\mathrm{\Omega })T_n}={\displaystyle \frac{\lambda }{\omega ^2}}{\displaystyle \frac{f_0}{u}}\widehat{}_1F.`$ (3.31)
Noting that the group velocity of the wave $`\mathrm{\Omega }_g=d\mathrm{\Omega }/dk`$ is given by
$`{\displaystyle \frac{\stackrel{~}{𝒟}}{k}}+{\displaystyle \frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}}{\displaystyle \frac{d\mathrm{\Omega }}{dk}}=0\mathrm{\Omega }_g={\displaystyle \frac{\stackrel{~}{𝒟}}{k}}\left({\displaystyle \frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}}\right)^1`$ (3.32)
we get
$`{\displaystyle \frac{E}{T_1}}+\mathrm{\Omega }_g{\displaystyle \frac{E}{z_1}}=i\lambda \left({\displaystyle \frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}}\right)^1\stackrel{~}{}{\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{F}{u}}E.`$ (3.33)
The above equations (3.31) and (3.33) are the amplitude equations to first order. Note that if $`Im(\mathrm{\Omega })=0`$ we could simply set $`F`$ equal to zero and then equation (3.33) would describe the symmetry properties of the original system (3.1) and (3.2) with respect to a linear plane wave solution. However, we are interested in the nonlinear interaction between waves (of increasing harmonicity) generated order by order, and as it can be easily seen the first nontrivial result taking into account this interaction will come out at third order. To pursue this we need the explicit (non secular) second order solutions for $`f_2`$ and $`V_2`$.
Solving the second order equation (3.22) with the remaining non secular part of the second order right hand side and then solving equation (3.18) with the already determined $`f_2`$ we find
$`f_2=S_F(k,\mathrm{\Omega },u)E^2e^{2i\phi }+c.c.+F_2(z_n,u;T_n),`$ (3.34)
$`V_2=S_V(k,\mathrm{\Omega })E^2e^{2i\phi }+f_Ve^{i\phi }+c.c.+G_V(z_n,T_n;\left[F\right]),`$ (3.35)
where $`c.c.`$ denotes complex conjugation. Without loss of generality we can set the generic function $`F_2(z_n,u;T_n)`$ equal to zero. Note that, in case $`Im(\mathrm{\Omega })=0`$ we could have set $`F=0`$, as mentioned earlier, but we should keep the function $`F_2`$ nonzero in order to cancel third order secular terms depending on the slow scales only. Moreover, the functions $`S_F`$, $`S_V`$, $`f_V`$ and the functional $`G_V\left(\left[F\right]\right)`$ of the amplitude $`F`$ are given by the following expressions:
$`S_F(k,\mathrm{\Omega },u)={\displaystyle \frac{\lambda ^2}{2}}{\displaystyle \frac{\stackrel{~}{𝒱}(2ik)}{\stackrel{~}{𝒟}(2k,2\mathrm{\Omega })}}{\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{}{u}}\left[{\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{f_0}{u}}\right],`$ (3.36)
$`S_V(k,\mathrm{\Omega })=\lambda ^2{\displaystyle \frac{\stackrel{~}{}(i\mathrm{\Omega },ik)}{\stackrel{~}{𝒟}(2k,2\mathrm{\Omega })}}{\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{}{u}}\left[{\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{f_0}{u}}\right],`$ (3.37)
$`f_V={\displaystyle \frac{i}{\stackrel{~}{𝒱}(ik)}}\left[i\lambda \stackrel{~}{}\widehat{}_1E\stackrel{~}{𝒱}_k(ik){\displaystyle \frac{E}{z_1}}\right],`$ (3.38)
$`G_V(z_n,T_n;\left[F\right])={\displaystyle \frac{1}{\omega ^2}}\widehat{}_1F,`$ (3.39)
$`\stackrel{~}{}(k,\mathrm{\Omega })={\displaystyle \frac{1}{\stackrel{~}{}(i\mathrm{\Omega },ik,u)}}{\displaystyle \frac{f_0}{u}},`$ (3.40)
where the $`k`$-index implies differentiation with respect to $`k`$.
The last step consists in evaluating the right hand side of equation (3.22), corresponding to the third order with the already found first and second order solutions. Removal of secular terms in the slow scales leads us finally to the amplitude equation for the function $`F(z_n,u;T_n)`$, that is
$$\frac{}{T_2}\left(\omega ^2F+\lambda F\frac{f_0}{u}\right)+\frac{2\lambda \gamma }{\omega ^2}\frac{f_0}{u}\frac{}{z_1}\widehat{}_1F+\lambda \frac{F}{u}\widehat{}_1F=$$
$`=\lambda ^2\omega ^2\left[{\displaystyle \frac{}{u}}\left({\displaystyle \frac{1}{\stackrel{~}{}^{}}}{\displaystyle \frac{f_0}{u}}\right)f_VE^{}+{\displaystyle \frac{}{u}}\left({\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{f_0}{u}}\right)f_V^{}E\right]e^{2Im(\mathrm{\Omega })T_2}.`$ (3.41)
Elimination of secular terms in the fast scales leads us to a generalized cubic Ginzburg–Landau type of equation for the amplitude $`E(z_n,T_n)`$:
$$i\frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}\frac{E}{T_2}=𝒜\frac{^2E}{z_1^2}+\lambda a\frac{}{z_1}\left\{𝒢\left(\left[F\right]\right)E\right\}+\lambda \left|E\right|^2Ee^{2Im(\mathrm{\Omega })T_2}$$
$`\lambda ^2𝒞G_V\left(\left[F\right]\right)E+\lambda \stackrel{~}{}𝒢\left(\left[F\right]\right)f_V,`$ (3.42)
where the coefficients $`a(k)`$, $`𝒜(k)`$, $`(k)`$ and $`𝒞(k)`$ are given by the expressions:
$`a(k)=\stackrel{~}{𝒱}_k\left({\displaystyle \frac{\stackrel{~}{𝒟}}{\mathrm{\Omega }}}\right)^1,`$ (3.43)
$`𝒜(k)=1+{\displaystyle \frac{\stackrel{~}{𝒱}_k}{\stackrel{~}{𝒱}}}\left[\stackrel{~}{𝒱}_k+i\lambda \stackrel{~}{}\left(1+\mathrm{\Omega }_g\right)\right],`$ (3.44)
$`(k)=\stackrel{~}{}{\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{S_F}{u}}\lambda \stackrel{~}{}S_V{\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{}{u}}\left({\displaystyle \frac{1}{\stackrel{~}{}^{}}}{\displaystyle \frac{f_0}{u}}\right),`$ (3.45)
$`𝒞(k)=\stackrel{~}{}{\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{}{u}}\left({\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{f_0}{u}}\right),`$ (3.46)
and the functional $`𝒢\left([F]\right)`$ of the amplitude $`F`$ can be written as
$`𝒢\left([F]\right)={\displaystyle \frac{1}{\stackrel{~}{}}}{\displaystyle \frac{F}{u}}.`$ (3.47)
Equations (3.41) and (3.42) comprise the system of coupled amplitude equations for the intensity of a resonator wave with a wave number $`k`$ and the slowly varying part of the longitudinal distribution function. Note that the dependence on the temporal scale $`T_1`$ (involving derivatives with respect to $`T_1`$) in equations (3.41) and (3.42) through the operator $`\widehat{}_1`$ and the function $`f_V`$ can be eliminated in principle by using the first order equations (3.31) and (3.33). As a result one obtains a system of coupled second order partial differential equations for $`F`$ and $`E`$ with respect to the variables $`T_2`$ and $`z_1`$.
## 4 Concluding Remarks
We have studied the longitudinal dynamics of particles moving in an accelerator under the action of a collective force due to a resonator voltage. For a sufficiently high beam density (relatively large value of the parameter $`\lambda `$) the nonlinear wave coupling, described by the nonlinear term in the Vlasov equation becomes important, and has to be taken into account. This is manifested in a spatio-temporal modulation of the wave amplitudes in unison with the slow process of particle redistribution. As a result of this wave-particle interaction (coupling between resonator waves and particle distribution modes) coherent, self-organized patterns can be formed in a wide range of relevant parameters.
We have analyzed the slow evolution of the amplitude of a single resonator wave with an arbitrary wave number $`k`$ (and wave frequency $`\mathrm{\Omega }(k)`$ defined as a solution to the dispersion relation). Using the method of multiple scales a system of coupled amplitude equations for the resonator wave envelope and for the slowly varying part of the longitudinal distribution function has been derived. As expected, the equation for the resonator wave envelope is a generalized cubic Ginzburg–Landau (GCGE) equation. We argue that these amplitude equations govern the (relatively) slow process of formation of coherent structures and establishment of wave-particle equilibrium.
### Acknowledgments
The author wishes to thank Y. Oono and C. Bohn for careful reading of the manuscript and for making valuable comments.
This work was supported by the US Department of Energy, Office of Basic Energy Sciences, under contract DE-AC03-76SF00515. |
warning/0001/hep-th0001153.html | ar5iv | text | # Baby Skyrme models for a class of potentials
## 1 Introduction
The Skyrme model describes a non-linear theory for SU(2) valued fields which has soliton solutions. Although the potential term is optional in the (3+1) dimensional nuclear Skyrme model its presence is necessary in (2+1) dimensions to ensure the stability of these solitonic solutions. In this paper we discuss multisolitons in a two-dimensional version of the model for a class of potentials which generalises the previously studied cases. In the literature we can find three specific potentials that have been studied in some detail.
The holomorphic model has only one stable solution (describing one skyrmion) and this solution has a simple analytical form . The old baby Skyrme model <sup>3</sup><sup>3</sup>3this name was introduced in has stable solutions for any number of skyrmions. In they were shown to lead to a crystalline lattice of skyrmions. The new baby Skyrme model was studied in detail in where it was shown that its solutions corresponded to field configurations with radially symmetrical energy densities - which correspond to many skyrmions lying “on top of each other”. We generalise these studies by considering potentials with a more general dependence on $`\varphi _3`$ thus leading to models with more vacua.
## 2 A Skyrme model for a class of potentials
The Skyrme model was first proposed by T.H.R Skyrme in 1960. Its classical solutions fall into various classes characterised by a topological number to be identified with the baryon number. Thus the model can describe mesons as well as various baryon configurations. The baby Skyrme models are (2+1) dimensional versions of the Skyrme model and its topologically nontrivial solutions are called baby skyrmions.
The Lagrangian density of the model contains three terms, from left to right: the pure $`S^2`$ sigma model, the Skyrme and the potential terms:
$$=_\mu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi }\theta _S\left[(_\mu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })^2(_\mu \stackrel{}{\varphi }_\nu \stackrel{}{\varphi })(^\mu \stackrel{}{\varphi }^\nu \stackrel{}{\varphi })\right]\theta _VV(\stackrel{}{\varphi }).$$
(1)
The vector $`\stackrel{}{\varphi }`$ is restricted to lie on a unit sphere $`𝒮^2`$ hence $`\stackrel{}{\varphi }\stackrel{}{\varphi }=1`$.
Note that to have finite potential energy the field at spatial infinity cannot depend on the polar angle $`\theta .`$
$$\underset{r\mathrm{}}{lim}\stackrel{}{\varphi }(r,\theta )=\stackrel{}{\varphi }^{(0)}.$$
(2)
Hence this boundary condition defines an one-point compatification of $`R_2`$, allowing us to consider $`\stackrel{}{\varphi }`$ on the extended plane $`R_2\mathrm{}`$ topologically equivalent to $`𝒮^2`$. In consequence, the field configurations are maps
$$:𝒮^2𝒮^2.$$
(3)
which can be labeled by an integer valued topological index $`Q`$:
$$Q=\frac{1}{4\pi }ϵ^{abc}𝑑x𝑑y\varphi _a\left(_x\varphi _b\right)\left(_y\varphi _c\right).$$
(4)
As a result of this non-trivial mapping the model has topologically nontrivial solutions which describe “extended structures”, namely, baby skyrmions . The different choices of the potential term lead to various shapes of the energy density of these baby skyrmions.
The equation of motion for a general potential that depends on $`\varphi _3`$ takes the form:
$`_\mu ^\mu \varphi _a(\stackrel{}{\varphi }_\mu ^\mu \varphi )\varphi _a2\theta _S[(_\nu \stackrel{}{\varphi }^\nu \stackrel{}{\varphi })_\mu ^\mu \varphi _a+(_\mu ^\nu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })_\nu \varphi _a`$
$`(^\nu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })_\nu _\mu \varphi _a(_\nu ^\nu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })_\mu \varphi _a+(_\mu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })(_\nu \stackrel{}{\varphi }^\nu \stackrel{}{\varphi })\varphi _a`$
$`(_\nu \stackrel{}{\varphi }_\mu \stackrel{}{\varphi })(^\nu \stackrel{}{\varphi }^\mu \stackrel{}{\varphi })\varphi _a]+{\displaystyle \frac{1}{2}}\theta _V{\displaystyle \frac{dV}{d\varphi _3}}(\delta _{a3}\varphi _a\varphi _3)=0.`$ (5)
It is convenient to rewrite this equation as:
$$_{tt}\varphi _a=K_{ab}^1_b(\stackrel{}{\varphi },_t\stackrel{}{\varphi },_i\stackrel{}{\varphi })$$
(6)
with
$$K_{ab}=(1+2\theta _S_i\stackrel{}{\varphi }_i\stackrel{}{\varphi })\delta _{ab}2\theta _S_i\varphi _a_i\varphi _b$$
(7)
and a rather complicated expression for $`_b`$.
Then to find a solution to this equation we invert the matrix $`K`$ and simulate the time evolution by a 4th order Runge Kutta method supplemented by the imposition of a correction due to the constraint $`\stackrel{}{\varphi }\stackrel{}{\varphi }=1`$.
The kinetic and potential energy densities of the baby Skyrme model are:
$$𝒦=\left(_t\stackrel{}{\varphi }_t\stackrel{}{\varphi }\right)\left(1+2\theta _S\left(_i\stackrel{}{\varphi }_i\stackrel{}{\varphi }\right)\right)2\theta _S\left(_t\stackrel{}{\varphi }_i\stackrel{}{\varphi }\right)^2$$
(8)
$$𝒱=\left(_i\stackrel{}{\varphi }_i\stackrel{}{\varphi }\right)+\theta _S\left[\left(_i\stackrel{}{\varphi }_i\stackrel{}{\varphi }\right)^2\left(_i\stackrel{}{\varphi }_j\stackrel{}{\varphi }\right)\left(_i\stackrel{}{\varphi }_j\stackrel{}{\varphi }\right)\right]+\theta _VV(\varphi _3).$$
(9)
The potentials we want to consider in this paper are of the form
$$V(\varphi _3)=(1\varphi _3)G(\varphi _3),$$
(10)
where $`G(\varphi _3)=1`$ or $`(1+\varphi _3)`$ or $`\varphi _3^2`$ or $`(1+\varphi _3)\varphi _3^2`$. Thus this potential vanishes at $`\varphi _3=1`$ and also at $`\varphi _3=1`$ and/or $`\varphi _3=0`$.
The total energy takes the form:
$$E=𝑑x𝑑y\left(𝒱+𝒦\right).$$
## 3 Static Solutions
In this paper we are primarily interested in finding minimal energy configurations corresponding to many skyrmions. We want to see, for different forms of the potential, what such configurations are like and what properties they have.
We will consider a numerically found configuration to be a multiskyrmion if it is satisfies certain numerical checks for a local minimum and if, moreover, its energy satisfies
$`E_n<E_k+E_l`$for all integers$`1<`$ $`\mathrm{}`$ and $`k<n`$such that $`k+\mathrm{}=n`$
We have included this condition in our definition because we want multisoliton to be stable with respect to decay into multiskyrmions of smaller degree.
There are many ways of finding static solutions. First, as the potential is only a function of $`\varphi _3`$, there is a symmetry corresponding to rotations around the $`\varphi _3`$ axis. So choosing the spatial dependence conveniently it is clear that fields which correspond to a generalised multihedgehog ansatz will be static solutions of the equations of motion . Such fields are given by ie
$$\stackrel{}{\varphi }=\left(\begin{array}{c}\mathrm{sin}[f(r)]\mathrm{cos}(n\theta \chi )\\ \mathrm{sin}[f(r)]\mathrm{sin}(n\theta \chi )\\ \mathrm{cos}[f(r)]\end{array}\right).$$
(11)
where (r,$`\theta `$) are polar coordinates in the $`xy`$-plane and the function $`f(r)`$, called the profile function, is required to satisfy certain boundary conditions to be specified below. The angle $`\chi `$ is arbitrary, but fields with different $`\chi `$ are related by an iso-rotation and are therefore degenerate in energy. $`n`$ is a non-zero integer and equals to the topological charge.
For the $`\varphi `$ field be regular at the origin the profile function $`f(r)`$ has to be satisfy :
$$f(0)=m\pi ,$$
(12)
where $`m`$ is an integer.
We choose the vacuum at infinity to be $`\stackrel{}{\varphi }^0=(0,0,1)`$ and this results in another boundary condition, namely:
$$\underset{r\mathrm{}}{lim}f(r)=0.$$
(13)
The total energy of the field configuration then takes the form:
$$E=(4\pi )\frac{1}{2}_0^{\mathrm{}}r𝑑r\left(f_{}^{^{}}{}_{}{}^{2}+n^2\frac{\mathrm{sin}^2f}{r^2}(1+2\theta _Sf_{}^{^{}}{}_{}{}^{2})+\theta _V\text{}(f)\right),$$
(14)
where $`f^{^{}}=\frac{df}{dr}`$ and $`\text{}(f)=V(\varphi _3)`$. To determine the profile function $`f(r)`$ we treat (14) as a variational problem and we get a second-order ODE for $`f`$:
$`\left(r+{\displaystyle \frac{2\theta _Sn^2\mathrm{sin}^2f}{r}}\right)f^{^{\prime \prime }}+\left(1{\displaystyle \frac{2\theta _Sn^2\mathrm{sin}^2f}{r^2}}+{\displaystyle \frac{2\theta _Sn^2\mathrm{sin}f\mathrm{cos}ff^{}}{r}}\right)f^{^{}}`$
$`{\displaystyle \frac{n^2\mathrm{sin}f\mathrm{cos}f}{r}}r{\displaystyle \frac{\theta _V}{2}}{\displaystyle \frac{d\text{}(f)}{df}}=0.`$ (15)
This equation then has to be solved numerically (ie via the shooting method).
The topological charge takes the form:
$$Q=\frac{n}{2}_0^{\mathrm{}}r𝑑r\left(\frac{f^{^{}}\mathrm{sin}f}{r}\right)=\frac{n}{2}[\mathrm{cos}f(\mathrm{})\mathrm{cos}f(0)].$$
(16)
This equation shows that to have an integer value for the topological charge, $`m`$ in (12) must be an odd number. In this paper we consider the solutions with $`m=1`$.
The behaviour of solutions of this equation near the origin and for large $`r`$ can be deduced analytically and was also discussed in . The result is that
* For small $`r`$, the profile function behaves as
$$f\pi +C_nr^n$$
(17)
and so
$$f^{^{}}nC_nr^{n1}$$
(18)
as long as $`\frac{dV(f)}{df}`$ tends to zero at this point.
* At large r, the ODE reduces to
$$f^{^{\prime \prime }}+\frac{1}{r}f^{^{}}\frac{n^2}{r^2}f\frac{\theta _V}{2}\frac{d\text{}(f)}{df}|_{\text{large r}}=0.$$
(19)
Using the boundary conditions, we see that if $`\frac{1}{f}\frac{d\text{}}{df}|_{\text{large r}}1`$ the profile function decreases exponentially
$$f(r)\frac{1}{\sqrt{\theta _Vr}}\mathrm{exp}(\theta _Vr).$$
(20)
This means that the potential localizes the skyrmion exponentially.
Of course we do not know whether such fields are the global minima of the energy. To check this we can perform a numerical simulation of (6) taking our derived field configurations (11) as initial conditions. We perturb them a little and then evolve them according to (6) with an extra dissipative term - $`\gamma _t\varphi _a`$ added to the right hand side of (6). We vary our perturbation and look at the configurations and their energies that our fields finally settle at. When the original configurations are the global minima of the energy the perturbed fields evolved back to them; otherwise (for sufficiently large perturbations) the fields evolve to other static solutions and, hopefully, the global minima. Another possibility, in particular to check whether these solutions are global minima, is to start with a more general field configuration corresponding to $`n`$ skyrmions, say, $`n`$ skyrmions located on a circle, and then evolve it using dissipative dynamics. In this case the circular set-up of $`n`$ 1-skyrmions with relative iso-orientation $`\delta \chi =\frac{2\pi }{n}`$ was obtained by using the combination of 1-skyrmions determined by the hedghog field (11) and combined together as discussed in the appendix.
The actual outcome (ie which fields are the global minima) depends on the form of the potential. In the next section we discuss the potentials that we have used in our simulations.
## 4 Different potentials
### 4.1 General comments
A skyrmion is a $`S^2S^2`$ map and so as the field at spatial $`\mathrm{}`$ corresponds to $`\varphi _3=+1`$ (ie the “North pole” of the field $`S^2`$) and the skyrmion’s position is at the point for which $`\varphi _3=1`$. If $`V(\varphi _3=1)`$ does not vanish it costs energy for the field to take this value; hence in this case we expect the skyrmions to repel when they are brought “on top of each other”. Thus the minimal energy multiskyrmion configurations should be different when $`V(\varphi _3=1)=0`$ and $`0`$. In fact, this was seen in the earlier studies; in the “new baby Skyrme” model ($`V(\varphi _3)=1\varphi _3^2`$) the skyrmions lie on “top of each other” while for the “old baby Skyrme” model ($`V(\varphi _3)=1\varphi _3`$) they are separated from each other by finite distances. We have repeated these simulations and have extended them to a larger number of skyrmions. Below, we present our results:
### 4.2 $`V=(1+\varphi _3)^4`$
The first potential studied in the baby Skyrme model was $`V=(1+\varphi _3)^4`$ (, , ). This potential was chosen because the equations of motion have an analytic static solution (for $`n=1`$) of the form $`W=\lambda (x+iy)`$, where $`\lambda `$ is related to $`\theta _S`$ and $`\theta _V`$. The resultant soliton is only polynomially localised. Two such skyrmions repel each other and scatter at 90 degrees when sent towards each other with sufficiant speed. The model has further solutions corresponding to several solitons “on top of each other”; but these solutions are unstable. When perturbed the solitons separate and move away to infinity. Thus this model have no stable multi-skyrmion solutions.
### 4.3 $`V=1\varphi _3`$
This potential was studied in some detail in and . As the value of the potential at $`r=0(\varphi _3=1)`$ is $`V(\varphi _3)=2`$ we expect some repulsion of the skyrmions when they are close together. Thus a field configuration of $`n`$ skyrmions on top of each other is unstable. This was studied a little in where it was shown that as a result of this repulsion the minimal energy multi-skyrmion field configurations produce nice lattice-like patterns. We have repeated these studies and extended them further to larger number of skyrmions. As there appear to be many solutions we have performed many simulations starting with different initial conditions and different perturbations.
Our results show a rather complicated pattern of minimal energy field configurations. The results are presented in table 1. We use the coefficients of $`\mathrm{𝑖𝑒}`$ add a factor $`\frac{1}{2}`$ to the sigma model term, $`\theta _S=0.25`$ and $`\theta _V=0.1`$.
All configurations are built out of 1, 2 and 3 skyrmions. In fig. 1 we show the energy densities of a pair and a triple.
A pair of skyrmions is very bound, more bound than a triple, and the pairs and triples are further bound for larger values of $`n`$. Thus the binding per skyrmion of a configuration of $`4`$ and $`5`$ skyrmions is less that of a single pair or a single triple, respectively.
The first more interesting case is of $`6`$ skyrmions which seems to have several bound states. Of these, the state corresponding to 3 pairs is the most bound, the other two (2 triples and 6 individual ones) are at best only local minima. For 7 skyrmions the lowest energy configuration corresponds to 2 pairs and a triple, which can be thought of a bound state of 4 and 3. The other state, corresponding to a triple sandwiched between two pairs has a higher energy and is at most a local minimum. In fig 2. we present energy densities of these two configurations.
For higher $`n`$ the situation becomes even more complicated; the lowest energy states always involve sets of pairs, but there seem to be other states involving triples and even singles.
One question one can ask is whether the $`n=2`$ state involves two skyrmions “on top of each other” or slightly displaced. When started from the original configuration of two skyrmions “on top of each other” slightly perturbed the system evolves a little, but this evolution is comparable to the original perturbation. Thus we believe the two skyrmions are slightly displaced, but this displacement is almost infinitesimal.
### 4.4 $`V=1\varphi _3^2`$
The model with this potential was investigated in great detail by Weidig . As both $`\varphi _3=\pm 1`$ correspond to the vacuum there is no energy argument which stops the skyrmions from “lying on top of each other”. This is in fact what Weidig saw in his simulations; the states found by the shooting method (and so corresponding to many skyrmions “on top of each other”) are the minima of the energy. Thus the energy densities of lowest energy multiskyrmion configurations have a ring-like structure.
### 4.5 $`V=\varphi _3^2(1\varphi _3^2)`$
The extra factor $`\varphi _3^2`$ in this potential causes it to have an additional zero and so the vacuum structure is even richer. However, like for the $`1\varphi _3^2`$ potential, there are no energetic arguments which would disfavour the skyrmions lying “on top of each other”. Indeed, this is what the simulations have shown; and we have checked this by starting with the hedgehog ansatz and $`n`$ solitons on a circle. The simulations have shown that the hedghog ansatz configurations are the minimal-energy solutions. The shapes of energy density have radial symmetry and resemble the configurations of Weidig. The difference is in the energy levels (in this potential the energy per skyrmion decreases monotonically with $`n`$). Our results are presented in table 2. In this (and the next table) we have added a factor $`\frac{1}{2}`$ to the sigma term and chose the $`\theta `$ coefficients to be $`\theta _S=0.2`$ and $`\theta _V=0.05`$. The values of the energies are divided by $`4\pi `$.
### 4.6 $`V=\varphi _3^2(1\varphi _3)`$
Like in the previous case the extra factor $`\varphi _3^2`$ changes the nature of the solutions. In a way the effects due to this potential are closer to those of the “old baby Skyrme model”, in that $`V(\varphi _3=1)0`$.
We have found that for this potential only the hedghog ansatz with $`n=1`$ has minimal-energy and that for $`1<n8`$ the skyrmions “on top of each other”, when perturbed, evolve into configurations corresponding to $`n`$-skyrmions that lie on regular polygons. When starting from a circular set-up for this potential we have calculated the total energy of $`n`$ 1-skyrmions as a function of the radius of the circle to get the initial state with minimum energy. This way the initial state settles down to the stable state sooner. We have confirmed this by performing also some simulations with different radi. We have also added to the initial configuration a non-symmetrical exponentially decaying perturbation to check whether the resultant multi-soliton states are stable, and whether, after dissipation, we reproduce the previous solutions.
In fig. 3 we present picture of the energy density of the lowest energy field configurations involving 6 skyrmions.
Our results show that this time (when we compare this case with $`V=1\varphi _3`$) the bindings are more comparable; in fact three skyrmions are more bound than two, and so the pattern is very different. It is interesting to note that the most bound system involves 6 skyrmions.
Our results on the binding energies etc are presented in table 3.
## 5 Comparison of potentials between 2 skyrmions
We have also looked at the “potentials” between two skyrmions in our models. To do this we had to decide how to construct field configurations involving two skyrmions. We decided to do this as follows: We have taken $`f(r)`$ for $`n=1`$ and computed the field $`W`$ from (30). Then we combined two such fields with $`\delta \chi =\pi `$ and varied the distance between them computing the energy as a function of the distance. We examined the following ways of combining 2 skyrmions:
$$W=W_1+W_2$$
(21)
$$W=W_1W_2$$
(22)
$$W=W_1+W_2W_1W_2$$
(23)
$$W=W_1W_2+W_1W_2.$$
(24)
As could be expected, we have found that for very small $`r`$ none of them gives reliable results. However, when two skyrmions are well separated the combination (23) gives us the lower energy than (21) and so approximates the 2 skyrmion field more accurately. The figures presented below show our numerical results for the following potentials
$$V=(1\varphi _3)$$
(25)
$$V=(1\varphi _3^2)$$
(26)
$$V=\varphi _3^2(1\varphi _3^2)$$
(27)
$$V=\varphi _3^2(1\varphi _3)$$
(28)
As expected when the two skyrmions are far from each other the energy approaches twice the energy of one skyrmion. In our plots fig 5. we have indicated the energy of the hedgehog field with $`n=2`$. Clearly, two skyrmions of this field are at $`r=0`$ but for better visualisation we have indicated it by a line. The low $`r`$ part of each plot (with points indicated by $`\mathrm{}`$) overestimates the energy and cannot be trusted (at these points the skyrmions deform each other and are not given by ansaetze (21-24). However, as the real energies are lower we see that our results confirm again that for the potentials (26) and (27) the hedgehog fields are minimal-energy solutions but for (25) and (28) the minimal-energy solutions are different. The difference is more obvious in the (28) case as the plot indicates that the minimum of the energy is at $`r=3.75`$. For (25) it may seem that the hedgehog field with $`n=2`$ is the minimal energy state but when we have evolved it the energy density expanded a bit thus showing that two baby skyrmions separate from each other a little. At the same time the energy has decreased and so we can conclude that for this potential the minimum of energy field configuration also differs from the hedgehog field. Incidentally the $`\theta ^{}s`$ used to calculate the potentials (25) and (26) were chosen so that all four curves are compareable. Thus they are different from values used in table 1 and 4.
## 6 Further Comments and Conclusions
The extra coefficient $`\varphi _3^2`$ causes the potential to have a further minimum (at $`\varphi _3=0`$). The question then arises of how this change effects the properties of the solutions. To perform such a comparison of multi-skyrmions of potentials (27) and (26) we have chosen the parameter $`\theta _S`$’s in both models to be such that the one skyrmion solutions of these models have the same energy.
Table 4 shows the energy per skyrmion for the new baby skyrmion model and for the model with the potential (27) respectively. The results indicate that a further minimum in the potential increases the binding of the skyrmions (ie the multiskyrmions have lower energies). The same seems to be true for the other two potentials.
We have also tried to understand the difference between the potentials (28) and (25) (in the (25) the minimum of the potential is very close to $`r=0`$ while for (28) it is clearly much further out). A possible suggestion is that, in the latter case, the skyrmions move to a distance between them at which $`\varphi _3=0`$. Unfortunately this idea is not supported by the results of our numerical simulations. We have looked at the plots of the energy density and of the $`\varphi _3`$ at different distances for 2 baby skyrmions and have found $`\varphi _30`$. When we repeated this study for larger $`\theta _V`$ (to have a more effective potential) the value of $`\varphi _3`$ decreased but was still nonzero. So the behaviour of skyrmions is more complicated and it depends on the properties of the potential in a more global way.
We have also calculated the ionisation energies and various break-up modes of multiskyrmions for potentials (27) and (28). Comparing potentials (25) and (28) we see some further similarity (in each case the system of 6 skyrmions seems to be the most bound) and in each case, eg, the state of 8 skyrmions needs very little energy to break up into (2+6). This is more pronounced in the case of the potential (28) for which the binding is stronger.
In this paper we have looked at baby Skyrme models with more general potentials. We have found that as the skyrmions involve mappings between $`S^2S^2`$ their properties depend crucially on whether the potential vanishes at the positions of the skyrmions (ie at $`\varphi _3=1`$) (we have assumed that at spatial infinity $`\varphi _3=1`$). When the potential vanishes the skyrmions “lie on top of each other”, when it does not - they separate and form interesting lattice like patterns. The shape of these patterns depends on the details of the potential. The same holds for the binding energies of skyrmions in all models. When the potential is more complicated (ie it has further zero) the skyrmions are more bound and in their patterns are more spread out - however, the actual positions and distortions depend on the details of the potential.
When the skyrmions are spread out the system of skyrmions has many local minima with some, larger or smaller, potential barriers between them. Thus for instance, for $`V=1\varphi _3`$, a system of 6 skyrmions has at least 3 local minima, and depending on the initial configuration the system can land in any of them. This is not unexpected and it suggests that similar models in, physically more relevant, 3 spatial dimensions may also have many local minima. Thus, the problem of finding multiskyrmion solutions of models in higher dimensions is clearly very complicated.
ACKNOWLEDEMENTS
We wish to thank Bernard Piette and Tom Weidig for useful comments and discussions. One of us (PE) thanks the University of Mashhad for a grant that made her visit to Durham possible and CPT, the University of Durham for its hospitality.
## 7 Appendix
Here we make some remarks about our numerical procedures
#### Profile functions
We have used the shooting method to determine $`f(r)`$ and have integrated (3) by a fourth-order Runge-Kutta method for any $`n`$. To avoid a singularity at $`r=0`$ we have considered $`rϵ`$ with $`ϵ`$ small. We have used the formulae (17) and (18) with 10000 lattice points and with the spacing $`dr=0.003`$.
#### Hedgehog static solutions
In most of our simulations we have used a $`201\times 201`$ lattice with lattice spacing $`\delta x=\delta y=0.3`$. For each $`(x_i,y_j)`$ $`f(x_i,y_j)`$ was determined by a linear interpolation using the values determined by the shooting method. Given $`f`$ we calculated $`\stackrel{}{\varphi }`$ using the hedgehog field expression (11)
#### A linear superposition for static solutions with $`n>1`$
In an alternative formulation we have used a single complex field $`W`$, which is related to $`\stackrel{}{\varphi }`$ by
$$\varphi _1=\frac{W+W^{}}{1+WW^{}}\varphi _2=i\frac{WW^{}}{1+WW^{}}\varphi _3=\frac{1WW^{}}{1+WW^{}}$$
(29)
where $``$ denotes the complex conjugation. Hence the complex field $`W`$ when expressed in terms of the profile function $`f(r)`$ and $`\theta `$ takes the form
$$W=\mathrm{tan}\left(\frac{f(r)}{2}\right)e^{in\theta }.$$
(30)
Static initial field configurations with $`Q=n`$ were formed by a linear superposition. When the baby skyrmions are far from each other (but not too close to the boundary) we can construct a configuration with charge $`n`$ from $`W`$ with $`n=1`$ by a linear superposition
$$W(x,y)=\underset{\alpha }{}W_\alpha (xx_\alpha ,yy_\alpha ),$$
(31)
where ($`x_\alpha `$,$`y_\alpha `$) is the location of the centre of the $`\alpha `$th skyrmion and $`W_\alpha `$ is the field (30) with $`n=1`$ and appropriate $`f(r)`$. We have used this method of constructing initial configurations in some of our simulations. Namely, we have placed $`n`$ baby skyrmions at equal distances from the origin with relative phase shifts $`\delta \chi =\frac{2\pi }{n}`$ between them (for maximal attraction) and then used (29) to get $`\stackrel{}{\varphi }`$.
#### Time evolution of the static solutions
We integrated the equation of motion for each component of $`\stackrel{}{\varphi }`$ independently. In this manner, during the simulations the field $`\stackrel{}{\varphi }`$ would gradually move away from the unit sphere $`𝒮^2`$. To correct this, every few iterations, we kept on rescaling the field as follows:
$$\varphi _a\frac{\varphi _a}{\sqrt{\stackrel{}{\varphi }\stackrel{}{\varphi }}}$$
(32)
and
$$_t\varphi _a_t\varphi _a\frac{_t\stackrel{}{\varphi }\stackrel{}{\varphi }}{\stackrel{}{\varphi }\stackrel{}{\varphi }}\varphi _a.$$
(33)
Another problem we have had to face involved using a finite lattice. Thus we had to make sure that the boundary effects did not alter our results. To be certain of this we varied lattice spacing and the number of lattice points.
When the fields were not at the minimum of the energy we allowed them to flow to this minimum - reducing the energy by a damping term.
$$_{tt}\varphi _a=K_{ab}^1_b(\stackrel{}{\varphi },_t\stackrel{}{\varphi },_i\stackrel{}{\varphi })\gamma _t\varphi _a,$$
(34)
where $`\gamma `$ is the damping coefficient. We set $`\gamma `$ to 0.1. This term takes the energy out of system.
We have used a fourth-order Runge-Kutta method to simulate the evolution of the field in time working in double precision. The time step was $`\delta t=0.1`$ throughout. We have performed our simulation using the $`\stackrel{}{\varphi }`$ formulation with the derivatives replaced by finite differences<sup>4</sup><sup>4</sup>4we have used the 9-point laplacian(as explained in , see also ). |
warning/0001/hep-lat0001021.html | ar5iv | text | # One-Loop Self Energy and Renormalization of the Speed of Light for some Anisotropic Improved Quark Actions
## I Introduction
Examples of successful employment of anisotropic lattices in lattice QCD simulations have been increasing lately. They include extensive studies of the glueball spectrum, investigations of heavy hybrid states and calculations of quarkonium fine structure . In most cases one is dealing with large states requiring large spatial volumes and also signals that can only be extracted from high statistics data. Working with highly improved actions on coarse lattices helps with the large volume and statistics problems, however, a coarse temporal lattice spacing means that correlation functions fall off very rapidly. This last problem can be circumvented by going to an anisotropic lattice which allows for a much finer temporal grid. The correlation functions can be sampled much more frequently in a given physical time region where the signal is still good.
Another potential use of anisotropic lattices would be in simulations of matrix elements in hadronic states with large momenta. These typically occur in semileptonic decays of heavy hadrons. Once one goes beyond spectrum calculations to matrix elements, one is faced with the matching problem between operators in the coarse highly improved lattice theory and continuum QCD. In this article we take a first step in accumulating necessary renormalization factors based on perturbation theory. We carry out the one-loop self energy calculation in several anisotropic improved quark actions. This gives us renormalization of the rest mass $`M_1`$, the wave function renormalization $`Z_2`$ and the “speed of light” renormalization, $`C_0`$, a quantity which will be defined more precisely in the following sections. We treat both massive and massless quarks. Perturbation theory can be used not only in operator matchings, but also to fix parameters in the lattice actions. $`C_0`$ is an example of one such parameter.
In the next section we introduce the gauge and quark actions considered in this article. Section 3 describes the general formalism that we employ for the self energy calculation with massive fermions. We follow closely the work of the Fermilab group which we could straightforwardly extend to anisotropic actions. Section 4 discusses specific one-loop contributions for mass, wavefunction and speed of light renormalizations. Our results are tabulated in section 5 for various choices of actions, fermion masses and degree of anisotropy. Some calculational details are left for appendices, where we describe Feynman rules and IR subtractions in our calculations.
## II Gauge and Quark Actions
We work with two classes of gauge actions denoted $`𝒮_G^I`$ and $`𝒮_G^{II}`$ , with
$`𝒮_G^I`$ $`=`$ $`\beta {\displaystyle \underset{x,s>s^{}}{}}{\displaystyle \frac{1}{\chi }}\left\{c_0^G{\displaystyle \frac{P_{ss^{}}}{u_s^4}}+c_1^G{\displaystyle \frac{R_{ss^{}}}{u_s^6}}+c_1^G{\displaystyle \frac{R_{s^{}s}}{u_s^6}}\right\}`$ (2)
$`\beta {\displaystyle \underset{x,s}{}}\chi \left\{c_0^G{\displaystyle \frac{P_{st}}{u_s^2u_t^2}}+c_1^G{\displaystyle \frac{R_{st}}{u_s^4u_t^2}}+c_1^G{\displaystyle \frac{R_{ts}}{u_t^4u_s^2}}\right\}`$
and
$`𝒮_G^{II}`$ $`=`$ $`\beta {\displaystyle \underset{x,s>s^{}}{}}{\displaystyle \frac{1}{\chi }}\left\{{\displaystyle \frac{5}{3}}{\displaystyle \frac{P_{ss^{}}}{u_s^4}}{\displaystyle \frac{1}{12}}{\displaystyle \frac{R_{ss^{}}}{u_s^6}}{\displaystyle \frac{1}{12}}{\displaystyle \frac{R_{s^{}s}}{u_s^6}}\right\}`$ (4)
$`\beta {\displaystyle \underset{x,s}{}}\chi \left\{{\displaystyle \frac{4}{3}}{\displaystyle \frac{P_{st}}{u_s^2u_t^2}}{\displaystyle \frac{1}{12}}{\displaystyle \frac{R_{st}}{u_s^4u_t^2}}\right\}.`$
The $`x`$ sum is over lattice sites and the variable $`s`$ runs over spatial directions. $`\beta 2N_c/g^2`$, $`\chi `$ is the anisotropy parameter
$$\chi =a_s/a_t$$
(5)
and
$`P_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{N_c}}Real\left(Tr\{U_\mu (x)U_\nu (x+a_\mu )U_\mu ^{}(x+a_\nu )U_\nu ^{}(x)\}\right),`$ (6)
$`R_{\mu \nu }`$ $`=`$ $`{\displaystyle \frac{1}{N_c}}Real\left(Tr\{U_\mu (x)U_\mu (x+a_\mu )U_\nu (x+2a_\mu )U_\mu ^{}(x+a_\mu +a_\nu )U_\mu ^{}(x+a_\nu )U_\nu ^{}(x)\}\right).`$ (7)
$`u_s`$ and $`u_t`$ are the tadpole improvement parameters $`u_0`$ for spatial and temporal link variables respectively .
The parameters $`c_0^G`$ and $`c_1^G`$ in action $`𝒮_G^I`$ are constrained to satisfy $`c_0^G+8c_1^G=1`$. The Symanzik improved gauge action, in which $`O(a^2)`$ errors are removed, corresponds to $`c_0^G=5/3`$ and $`c_1^G=1/12`$ , whereas $`c_0^G=3.648`$ and $`c_1^G=0.331`$, for $`\chi =1`$, leads to one of the RG improved Iwasaki actions . In $`𝒮_G^{II}`$ parameters have been fixed to the Symanzik values. We will be working mainly with Symanzik improved actions and present RG improved results only for a few cases. We note that the action $`𝒮_G^{II}`$ is intrinsically asymmetric even for the isotropic limit $`\chi =1`$.
The most highly improved quark action that we have analysed is the $`D234`$ action .
$`𝒮_{D234}^I`$ $`=`$ $`a_s^3a_t{\displaystyle \underset{x}{}}\overline{\mathrm{\Psi }}_c\{\gamma _t{\displaystyle \frac{1}{a_t}}(_t{\displaystyle \frac{1}{6}}C_{3t}_t^{(3)})+{\displaystyle \frac{C_0}{a_s}}\stackrel{}{\gamma }(\stackrel{}{}{\displaystyle \frac{1}{6}}C_3\stackrel{}{}^{(3)})+m_0`$ (10)
$`{\displaystyle \frac{ra_s}{2}}\left[{\displaystyle \frac{1}{a_t^2}}(_t^{(2)}{\displaystyle \frac{1}{12}}C_{4t}_t^{(4)})+{\displaystyle \frac{1}{a_s^2}}{\displaystyle \underset{j}{}}(_j^{(2)}{\displaystyle \frac{1}{12}}C_4_j^{(4)})\right]`$
$`ra_s{\displaystyle \frac{C_F}{4}}{\displaystyle \frac{i\sigma _{\mu \nu }\stackrel{~}{F}^{\mu \nu }}{a_\mu a_\nu }}\}\mathrm{\Psi }_c`$
$`=`$ $`{\displaystyle \underset{x}{}}\overline{\mathrm{\Psi }}_L\{\gamma _t(_t{\displaystyle \frac{1}{6}}C_{3t}_t^{(3)})+{\displaystyle \frac{C_0}{\chi }}\stackrel{}{\gamma }(\stackrel{}{}{\displaystyle \frac{1}{6}}C_3\stackrel{}{}^{(3)})+a_tm_0`$ (13)
$`{\displaystyle \frac{r}{2}}\left[\chi (_t^{(2)}{\displaystyle \frac{1}{12}}C_{4t}_t^{(4)})+{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}(_j^{(2)}{\displaystyle \frac{1}{12}}C_4_j^{(4)})\right]`$
$`r{\displaystyle \frac{C_F}{4}}i\sigma _{\mu \nu }\stackrel{~}{F}^{\mu \nu }{\displaystyle \frac{a_sa_t}{a_\mu a_\nu }}\}\mathrm{\Psi }_L.`$
The quark fields $`\mathrm{\Psi }_c`$ and the dimensionless lattice fields $`\mathrm{\Psi }_L`$ are related through
$$\mathrm{\Psi }_L=a_s^{3/2}\mathrm{\Psi }_c.$$
(14)
The dimensionless derivatives $`^{(n)}`$ and field strength tensors $`\stackrel{~}{F}^{\mu \nu }`$ are tadpole improved and defined in the Appendix. We use the convention $`\sigma _{\mu \nu }=\frac{1}{2}[\gamma _\mu ,\gamma _\nu ]`$ and set $`r=1`$ in all our calculations. At tree-level the coefficients $`C_0`$, $`C_3`$, $`C_{3t}`$, $`C_4`$, $`C_{4t}`$ and $`C_F`$ are equal to one. The quark action is then tree-level accurate through $`O(a_s^3)`$ and $`O(a_t^3)`$. $`C_0`$ is what we call the “speed of light”. This parameter is adjusted, in general either perturbatively or nonperturbatively, to ensure correct dispersion relations for particles. In anticipation of working on anisotropic lattices with $`a_t`$ much finer than $`a_s`$, one can drop the higher order improvement terms in the temporal derivatives by setting $`C_{3t}=C_{4t}=0`$, without loosing accuracy. We call this action $`𝒮_{D234}^{II}`$.
$`𝒮_{D234}^{II}`$ $`=`$ $`{\displaystyle \underset{x}{}}\overline{\mathrm{\Psi }}_L\{\gamma _t_t+{\displaystyle \frac{C_0}{\chi }}\stackrel{}{\gamma }(\stackrel{}{}{\displaystyle \frac{1}{6}}C_3\stackrel{}{}^{(3)})+a_tm_0`$ (17)
$`{\displaystyle \frac{r}{2}}\left[\chi _t^{(2)}+{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}(_j^{(2)}{\displaystyle \frac{1}{12}}C_4_j^{(4)})\right]`$
$`r{\displaystyle \frac{C_F}{4}}i\sigma _{\mu \nu }\stackrel{~}{F}^{\mu \nu }{\displaystyle \frac{a_sa_t}{a_\mu a_\nu }}\}\mathrm{\Psi }_L`$
The familiar $`O(a)`$ accurate clover quark action corresponds to setting $`C_3=C_4=0`$ in the above and using a less improved field strength tensor $`F^{\mu \nu }`$ (also defined in the Appendix) rather than $`\stackrel{~}{F}^{\mu \nu }`$.
$`𝒮_{clover}`$ $`=`$ $`{\displaystyle \underset{x}{}}\overline{\mathrm{\Psi }}_L\{\gamma _t_t+{\displaystyle \frac{C_0}{\chi }}\stackrel{}{\gamma }\stackrel{}{}+a_tm_0`$ (19)
$`{\displaystyle \frac{r}{2}}[\chi _t^{(2)}+{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}_j^{(2)}]r{\displaystyle \frac{C_F}{4}}i\sigma _{\mu \nu }F^{\mu \nu }{\displaystyle \frac{a_sa_t}{a_\mu a_\nu }}\}\mathrm{\Psi }_L`$
We have carried out one-loop self energy calculations for several combinations of the above gauge and quark actions, for both massless and massive quarks. We list the specific actions considered in Table I. For actions $`𝒮^A`$ and $`𝒮^A^{}`$ massless results have already appeared in . We agree with their results and we include these cases here for completeness. With action $`𝒮^C`$ we treat only the massless case, since our formalism for massive quarks, following , requires that the only time derivatives be in the $`_t`$ and $`_t^{(2)}`$ terms. Both $`𝒮_{D234}^{II}`$ and $`𝒮_{clover}`$ satisfy this condition, but $`𝒮_{D234}^I`$ does not.
## III General Formalism for Self Energy Calculations
In this section we summarize the formalism for self energy calculations, along the lines of reference . Perturbative calculations for massive Wilson quarks are also described in reference . We concentrate on the massive case, since massless lattice perturbation theory has been in the literature for decades.
For massive fermions we use either $`𝒮_{D234}^{II}`$ or $`𝒮_{clover}`$. The fermion self energy $`\mathrm{\Sigma }(p)`$ is defined in terms of the momentum space propagators $`\overline{G}(p)`$ and $`\overline{G}_0(p)`$ for the full and free theories respectively, as
$$\overline{G}^1(p)=\overline{G}_0^1(p)\mathrm{\Sigma }(p).$$
(20)
Carrying out the Fourier transform in $`p_0`$ one defines
$`G(t,\stackrel{}{p})`$ $`=`$ $`{\displaystyle _{\pi /a_t}^{\pi /a_t}}{\displaystyle \frac{dp_0}{2\pi }}e^{ip_0t}\overline{G}(p_0,\stackrel{}{p})`$ (21)
$``$ $`𝒵_2(\stackrel{}{p})e^{E(\stackrel{}{p})t}\mathrm{\Gamma }_{proj}+\mathrm{}.`$ (22)
$`\mathrm{\Gamma }_{proj}`$ is a projection operator in Dirac space. The ellipses refer to lattice artifacts and additional multi-particle states that could be created by the lattice fermion field operator $`\mathrm{\Psi }`$ beyond the single quark state. The rest mass, $`M_1`$, is defined as
$$M_1=E(\stackrel{}{p}=\stackrel{}{0}).$$
(23)
We do not consider the kinetic mass, $`M_2`$ in this article. We will renormalize at the point $`(p_0,\stackrel{}{p})=(iM_1,\stackrel{}{0})`$ and define the wave function renormalization constant
$$Z_2=𝒵_2(\stackrel{}{p}=\stackrel{}{0}).$$
(24)
For a zero spatial momentum quark propagating forward in time one expects ($`t>0`$)
$`G(t,0)`$ $`=`$ $`{\displaystyle _{\pi /a_t}^{\pi /a_t}}{\displaystyle \frac{dp_0}{2\pi }}e^{ip_0t}\overline{G}(p_0,0)`$ (25)
$``$ $`Z_2e^{M_1t}{\displaystyle \frac{1+\gamma _0}{2}}+\mathrm{}.`$ (26)
Our goal in this section is to relate $`Z_2`$ and $`M_1`$ to parameters in the action and to $`\mathrm{\Sigma }(p)`$. In order to orient ourselves, however, it is useful to first consider the free case with $`\mathrm{\Sigma }(p)=0`$.
### A Free Anisotropic Propagator
The free propagator $`\overline{G}_0(p_0,\stackrel{}{p}=0)`$ for both actions $`𝒮_{D234}^{II}`$ and $`𝒮_{clover}`$ becomes (for $`r=1`$)
$`{\displaystyle \frac{1}{a_t}}\overline{G}_0(p_0,\stackrel{}{p}=0)`$ $`=`$ $`{\displaystyle \frac{1}{i\gamma _0\mathrm{sin}(a_tp_0)+a_tm_0+\chi \chi \mathrm{cos}(a_tp_0)}}`$ (27)
$`=`$ $`{\displaystyle \frac{i\gamma _0\mathrm{sin}(a_tp_0)+a_tm_0+\chi \chi \mathrm{cos}(a_tp_0)}{(\mathrm{sin}(a_tp_0))^2+[(a_tm_0+\chi )\chi \mathrm{cos}(a_tp_0)]^2}}.`$ (28)
In terms of the variable
$$ze^{ia_tp_0}=e^{a_tE}$$
(29)
($`p_0=iE`$), one finds two zeros of the denominator corresponding to positive energy solutions.
$$z_1=\frac{(a_tm_0+\chi )\sqrt{(a_tm_0+\chi )^2+1\chi ^2}}{\chi 1}$$
(30)
and
$$\stackrel{~}{z}_1=\frac{(a_tm_0+\chi )\sqrt{(a_tm_0+\chi )^2+1\chi ^2}}{\chi +1}.$$
(31)
The other two zeros, $`z_2`$ and $`\stackrel{~}{z}_2`$ correspond to negative energy solutions, $`z_2=1/z_1`$ and $`\stackrel{~}{z}_2=1/\stackrel{~}{z}_1`$. The integral over $`p_0`$ in (25) can be done as a contour integral around the unit circle in the variable $`z`$. One picks up contributions from both positive energy solutions (for $`t>0`$).
$`poleatz_1`$ $`:`$ $`{\displaystyle \frac{(1+\gamma _0)}{2}}{\displaystyle \frac{e^{M_1^{(0)}t}}{\sqrt{(a_tm_0+\chi )^2+1\chi ^2}}},`$ (32)
$`poleat\stackrel{~}{z}_1`$ $`:`$ $`{\displaystyle \frac{(1\gamma _0)}{2}}{\displaystyle \frac{e^{\stackrel{~}{M}_1^{(0)}t}}{\sqrt{(a_tm_0+\chi )^2+1\chi ^2}}},`$ (33)
with
$$a_tM_1^{(0)}=\mathrm{ln}(z_1),a_t\stackrel{~}{M}_1^{(0)}=\mathrm{ln}(\stackrel{~}{z}_1).$$
(34)
Clearly, $`z_1`$ is the physical positive energy solution. The second solution $`\stackrel{~}{z}_1`$ is a lattice artifact, similar to the time doubler for $`r1`$ in isotropic actions. The solution $`\stackrel{~}{z}_1`$ disappears in the isotropic limit, $`\chi 1`$, where $`a_t\stackrel{~}{M}_1^{(0)}\mathrm{}`$. In the same limit the physical solution $`z_1`$ goes over into the well known result
$$z_1\frac{1}{1+a_tm_0}.$$
(35)
The gap between $`\stackrel{~}{M}_1^{(0)}`$ and $`M_1^{(0)}`$, measured in units of $`1/a_s`$ is
$$a_s(\stackrel{~}{M}_1^{(0)}M_1^{(0)})=\chi \mathrm{ln}\frac{(\chi +1)}{(\chi 1)},$$
(36)
independent of $`m_0`$. This becomes $`\mathrm{}`$ at $`\chi =1`$ and approaches $`2`$ as $`\chi \mathrm{}`$. The size of this gap, $`a_s\delta E2`$, is hence equal to or larger than the amount by which conventional spatial doublers are raised through the Wilson mechanism. We will henceforth ignore $`\stackrel{~}{z}_1`$ and concentrate on the physical pole at $`z=z_1`$. Comparing (32) with (25) one sees that there is nontrivial mass dependent wave function renormalization even at tree-level with
$$Z_2^{(0)}=\frac{1}{\sqrt{(a_tm_0+\chi )^2+1\chi ^2}}=\frac{1}{\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})}.$$
(37)
This has been pointed out several times in the literature .
A useful way to rewrite (30) is
$$a_tm_0+\chi =\chi \mathrm{cosh}(a_tM_1^{(0)})+\mathrm{sinh}(a_tM_1^{(0)}).$$
(38)
### B Mass Renormalization
In the interacting case one has a nontrivial $`\mathrm{\Sigma }(p)`$ which we write as
$$a_t\mathrm{\Sigma }(p)=i\gamma _0B_0(p,m_0)\mathrm{sin}(a_tp_0)+i\frac{1}{\chi }\underset{j}{}[\gamma _jB_j(p,m_0)\mathrm{sin}(a_sp_j)]+C(p,m_0).$$
(39)
The $`\stackrel{}{p}=0`$ propagator becomes
$$\frac{1}{a_t}\overline{G}(p_0,\stackrel{}{p}=0)=\frac{i\gamma _0(1B_0)\mathrm{sin}(a_tp_0)+a_tm_0+\chi C\chi \mathrm{cos}(a_tp_0)}{(1B_0)^2[\mathrm{sin}(a_tp_0)]^2+[(a_tm_0+\chi )C\chi \mathrm{cos}(a_tp_0)]^2},$$
(40)
where, $`B_0=B_0(p_0,m_0)`$ and $`C=C(p_0,m_0)`$ are evaluated at $`\stackrel{}{p}=0`$. If $`p_0=iE`$ is the location of a pole in (40), the following implicit equation must be satisfied.
$$(1B_0(iE,m_0))\mathrm{sinh}(a_tE)=\pm [(a_tm_0+\chi )C(iE,m_0)\chi \mathrm{cosh}(a_tE)].$$
(41)
One can check that the “+” sign leads to the pole $`M_1^{(0)}`$ in the free limit. Hence, the implicit equation for $`M_1`$ is given by
$$\chi \mathrm{cosh}(a_tM_1)+\mathrm{sinh}(a_tM_1)=a_tm_0+\chi +B_0(iM_1,m_0)\mathrm{sinh}(a_tM_1)C(iM_1,m_0).$$
(42)
In a perturbative calculation of $`M_1`$ one expands
$$M_1=M_1^{(0)}+\alpha _sM_1^{(1)}+O(\alpha _s^2).$$
(43)
$`B_0`$ and $`C`$ in (42) start out $`O(\alpha _s)`$, so through one-loop their argument can be replaced by the tree-level $`M_1^{(0)}`$. Expanding the LHS also through $`O(\alpha _s)`$ and taking (38) into account, one finds
$`\alpha _sa_tM_1^{(1)}`$ $`=`$ $`{\displaystyle \frac{B_0(iM_1^{(0)},m_0)\mathrm{sinh}(a_tM_1^{(0)})C(iM_1^{(0)},m_0)}{\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})}}`$ (44)
$`=`$ $`Z_2^{(0)}tr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}a_t\mathrm{\Sigma }(p_0=iM_1^{(0)},\stackrel{}{p}=0)\right\},`$ (45)
where the trace is taken over Dirac space. We note that in the $`M_1^{(0)}=0,m_0=0`$ limit, the $`\gamma _0`$ part of the trace $`tr\{(\gamma _0+1)\mathrm{\Sigma }\}`$ does not contribute and one has
$$\alpha _sa_tM_1^{(1)}(0)=tr\left\{a_t\mathrm{\Sigma }(0)\right\}/4=C(0,0).$$
(46)
In order to have massless quarks remain massless under renormalization, one needs to carry out additive mass renormalization and $`M_1^{(1)}`$ in (44) requires a subtraction. This subtraction must be done without jeopardizing the pole condition (42). There is a standard way to accomplish this. Let $`m_c`$ be the value of the bare quark mass parameter $`m_0`$ for which the physical quark rest mass vanishes ($`M_1=0)`$. Eqn.(42) then tells us that $`m_c`$ is implicitly defined through
$$a_tm_cC(0,m_c)=0.$$
(47)
In equations such as (40) or (42) one always has the combination $`a_tm_0C`$. Using (47) one can add and subtract $`a_tm_c`$ so that
$$a_tm_0Ca_t(m_0m_c)(CC(0,m_c))=a_tm\stackrel{~}{C}.$$
(48)
Previous derivations go through with $`m_0`$ replaced by
$$mm_0m_c$$
(49)
and $`C(iM_1,m_0)`$ by
$$\stackrel{~}{C}(iM_1,m_0)=C(iM_1,m_0)C(0,m_c).$$
(50)
In most lattice simulations, $`m_c`$ and hence also $`m`$ are determined nonperturbatively from the simulations themselves. For $`\stackrel{~}{C}`$, however, one still often uses the one-loop result
$$\stackrel{~}{C}(iM_1^{(0)},m_0)=\stackrel{~}{C}(iM_1^{(0)},m)=C(iM_1^{(0)},m)C(0,0).$$
(51)
$`M_1^{(0)}`$ is now given in terms of $`m`$ rather than $`m_0`$. We will be presenting our results as functions of $`a_sM_1^{(0)}`$, with the understanding that the shift $`m_0m_0m_c`$ has been carried out and that, for instance, $`M_1^{(0)}`$ is given by
$$a_tm+\chi =\chi \mathrm{cosh}(a_tM_1^{(0)})+\mathrm{sinh}(a_tM_1^{(0)})$$
(52)
rather than by (38). In (44) one needs to replace $`C`$ by $`\stackrel{~}{C}`$. Our final formula for the one-loop mass correction, measured in units of $`1/a_s`$ then becomes
$`\alpha _sa_sM_{1,sub}^{(1)}`$ $`=`$ $`\chi {\displaystyle \frac{B_0(iM_1^{(0)},m)\mathrm{sinh}(a_tM_1^{(0)})\stackrel{~}{C}(iM_1^{(0)},m)}{\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})}}`$ (53)
$`=`$ $`Z_2^{(0)}tr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}\left[a_s\mathrm{\Sigma }(p_0=iM_1^{(0)},\stackrel{}{p}=0,m)a_s\mathrm{\Sigma }(0,\stackrel{}{0},0)\right]\right\}.`$ (54)
This expression vanishes automatically for $`M_1^{(0)}=m=0`$. We prefer to measure dimensionful quantities in terms of $`1/a_s`$ rather than $`1/a_t`$. When exploring $`\chi 1`$ it makes more sense to fix $`a_s`$ and let $`a_t`$ be arbitrarily fine, rather than to fix $`a_t`$ and let $`a_s`$ become arbitrarily coarse. In the isotropic limit (53) agrees with formulas in the literature .
### C Wave Function Renormalization
In order to extract a general formula for the wave function renormalization $`Z_2`$ we need to find the residue of $`\overline{G}(p_0,\stackrel{}{p}=0)`$ at the pole $`p_0=iM_1`$. In terms of the variable $`z`$ the Fourier transform in (25) has the form
$$_{|z|=1}\frac{dz}{(2\pi i)z}(z)^{t/a_t}\frac{g(z)}{f(z)},$$
(55)
where the integral is taken over the unit circle. To find the residue we expand the denominator around $`z_1=e^{a_tM_1}`$
$$f(z)=(zz_1)\left(\frac{df}{dz}\right)_{z=z_1}+\mathrm{}.$$
(56)
The contribution from the physical pole to $`G(t,0)`$ is then
$$e^{M_1t}\left(\frac{g(z)}{zf^{}(z)}\right)_{z=z_1}.$$
(57)
One finds for the numerator
$$g(z=z_1)=(\gamma _0+1)(1B_0(iM_1,m))\mathrm{sinh}(a_tM_1)$$
(58)
and for the denominator
$`2(1B_0(iM_1,m))\mathrm{sinh}(a_tM_1)\times `$ (59)
$`\left\{\chi \mathrm{sinh}(a_tM_1)+\mathrm{cosh}(a_tM_1)+\left(i{\displaystyle \frac{d}{d(a_tp_0)}}[iB_0(p_0,m)\mathrm{sin}(a_tp_0)+C(p_0,m)]\right)_{p_0=iM_1}\right\}`$ (60)
using
$$\left(z\frac{df}{dz}\right)_{z=z_1}=i\left(\frac{df}{d(a_tp_0)}\right)_{p_0=iM_1}.$$
(61)
One can now read off $`Z_2`$ and after recognizing the last term in (59) as derivatives acting on different parts of $`a_t\mathrm{\Sigma }(p_0,\stackrel{}{p}=0,m)`$, one obtains
$$Z_2^1=\chi \mathrm{sinh}(a_tM_1)+\mathrm{cosh}(a_tM_1)+itr\left(\frac{(\gamma _0+1)}{4}\frac{d}{dp_0}\mathrm{\Sigma }(p_0,\stackrel{}{p}=0,m)\right)_{p_0=iM_1}.$$
(62)
The one-loop approximation to $`Z_2`$ is obtained by expanding $`M_1`$ once again in $`\alpha _s`$.
$`Z_2^1`$ $`=`$ $`\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})+\alpha _sa_tM_{1,sub}^{(1)}(\chi \mathrm{cosh}(a_tM_1^{(0)})+\mathrm{sinh}(a_tM_1^{(0)}))`$ (64)
$`+itr\left({\displaystyle \frac{(\gamma _0+1)}{4}}{\displaystyle \frac{d}{dp_0}}\mathrm{\Sigma }(p_0,\stackrel{}{p}=0,m)\right)_{p_0=iM_1^{(0)}}`$
$`=`$ $`Z_2^{(0)1}[\mathrm{\hspace{0.33em}1}+{\displaystyle \frac{\alpha _s}{\chi }}a_sM_{1,sub}^{(1)}(\chi \mathrm{cosh}(a_tM_1^{(0)})+\mathrm{sinh}(a_tM_1^{(0)}))Z_2^{(0)}`$ (66)
$`+itr\left({\displaystyle \frac{(\gamma _0+1)}{4}}{\displaystyle \frac{d}{dp_0}}\mathrm{\Sigma }(p_0,\stackrel{}{p}=0,m)\right)_{p_0=iM_1^{(0)}}Z_2^{(0)}]+O(\alpha _s^2).`$
In the last expression we have found it convenient to factor out the tree-level $`Z_2^{(0)1}`$. Equations (62) and (64) go over into the formulas of in the isotropic limit.
### D Speed of Light Renormalization
In order to discuss renormalization of the speed of light one needs to look at the inverse momentum space propagator at small but nonzero spatial momentum.
$`a_t\overline{G}^1(p)`$ $`=`$ $`a_t\overline{G}_0^1(p)a_t\mathrm{\Sigma }(p)`$ (67)
$`=`$ $`i\gamma _0(1B_0)\mathrm{sin}(a_tp_0)+i{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}[\gamma _j(C_0K_jB_j)\mathrm{sin}(a_sp_j)]`$ (69)
$`+a_tm+\chi \chi \mathrm{cos}(a_tp_0)C,`$
with $`K_j=1`$ for $`𝒮_{clover}`$ and $`K_j=(4\mathrm{cos}(a_sp_j))/3`$ for $`𝒮_{D234}^{I,II}`$. One can rewrite $`\overline{G}^1(p)`$ as
$`a_t\overline{G}^1(p)`$ $`=`$ $`(1B_0)\left\{i\gamma _0\mathrm{sin}(a_tp_0)+i{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}[\gamma _j{\displaystyle \frac{(C_0K_jB_j)}{(1B_0)}}\mathrm{sin}(a_sp_j)]\right\}`$ (71)
$`+a_tm+\chi \chi \mathrm{cos}(a_tp_0)C.`$
$`C_0`$ is adjusted so that for small $`a_sp_j`$ the relative coefficient of the $`\gamma _0\mathrm{sin}(a_tp_0)`$ and the $`\gamma _ja_sp_j/\chi `$ terms remains equal to unity. $`(C_0K_jB_j)\mathrm{sin}(a_sp_j)=(C_0B_j)a_sp_j`$ for all quark actions in this limit (of course, $`K_j\mathrm{sin}(a_sp_j)`$ is a better approximation to the continuum $`a_sp_j`$ in the $`D234`$ action than in the clover action), and one has
$$C_0=1+B_j(C_0)B_0(C_0)\mathrm{\hspace{0.33em}1}+B_j(C_0=1)B_0(C_0=1)+O(\alpha _s^2).$$
(72)
Just as with $`Z_2`$ we will define the speed of light renormalization at the zero spatial momentum mass shell point $`p=(iM_1,\stackrel{}{0})`$. From (39) the two terms $`B_j`$ and $`B_0`$ needed for $`C_0`$ at one-loop can be extracted through
$`B_j`$ $`=`$ $`i{\displaystyle \frac{\chi }{4}}tr\left(\gamma _j{\displaystyle \frac{}{(a_sp_j)}}a_t\mathrm{\Sigma }(p)\right)_{p=(iM_1^{(0)},\stackrel{}{0})},`$ (73)
$`B_0`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \frac{tr(\gamma _0a_t\mathrm{\Sigma }(iM_1^{(0)},\stackrel{}{0}))}{\mathrm{sinh}(a_tM_1^{(0)})}}m>0`$ (74)
$`or`$ (75)
$`=`$ $`{\displaystyle \frac{i}{4}}tr\left(\gamma _0{\displaystyle \frac{}{p_0}}\mathrm{\Sigma }(p)\right)_{p=(0,\stackrel{}{0})}m=0.`$ (76)
We note that in the massive case there is nontrivial renormalization of $`C_0`$ even in the isotropic limit $`\chi =1`$, due to our noncovariant mass shell condition, $`p=(iM_1,\stackrel{}{0})`$. Nevertheless, we believe the above definition of the renormalization of the speed of light is a sensible and physical one.
## IV One-Loop Contributions to $`\mathrm{\Sigma }(p)`$
In the previous section one-loop corrections for $`M_1`$, $`Z_2`$ and $`C_0`$ were determined in terms of traces over $`\mathrm{\Sigma }(p)`$ or over derivatives acting on $`\mathrm{\Sigma }(p)`$. In this section we describe the lattice perturbation theory diagrams that contribute to $`\mathrm{\Sigma }(p)`$ at one loop. For all quark actions considered one can write
$$\mathrm{\Sigma }(p)=\mathrm{\Sigma }^{reg}(p)+\mathrm{\Sigma }^{tad}(p)+\mathrm{\Sigma }^{t.i.}(p).$$
(77)
$`\mathrm{\Sigma }^{reg}`$ is the regular “rainbow” diagram, the only diagram that exists in a continuum self energy calculation. $`\mathrm{\Sigma }^{tad}`$ denotes contributions from the lattice artifact tadpole diagram and $`\mathrm{\Sigma }^{t.i.}`$ comes from perturbatively expanding the $`u_s`$’s and $`u_t`$’s entering definitions of tadpole improved derivatives (see Appendix). The main idea behind tadpole-improved perturbation theory is to have $`\mathrm{\Sigma }^{t.i.}`$ cancel the bulk of $`\mathrm{\Sigma }^{tad}`$. We find in many instances, especially with $`𝒮_{clover}`$, that cancellation is complete if one uses the Landau link definition for $`u_s`$ and $`u_t`$ and works in Landau gauge.
### A $`\mathrm{\Sigma }^{reg}(p)`$
In terms of the gauge propagator $`D_{\mu \nu }`$, quark propagator $`\overline{G}_0`$ and the vertex functions $`V_\mu `$, one can write $`a_t\mathrm{\Sigma }^{reg}(p)`$ as the following loop integral over the dimensionless momentum variables $`\pi k_\mu \pi `$.
$`a_t\mathrm{\Sigma }^{reg}(p)`$ (78)
$`=g^2{\displaystyle \frac{4}{3}}{\displaystyle \underset{\mu ,\nu }{}}{\displaystyle \frac{d^4k}{(2\pi )^4}\left\{V_\mu (ap,apk)\frac{\overline{G}_0(apk)}{a_t}V_\nu (apk,ap)\right\}D_{\mu \nu }(k,\alpha _g)}`$ (79)
$`=g^2{\displaystyle \frac{4}{3}}{\displaystyle \underset{\mu ,\nu }{}}{\displaystyle \frac{d^4k}{(2\pi )^4}\left\{V_\mu (ap,apk)[i\gamma K\mathrm{sin}+\mathrm{\Omega }]_{(apk)}V_\nu (apk,ap)\right\}\frac{D_{\mu \nu }(k,\alpha _g)}{(K^2\mathrm{sin}^2+\mathrm{\Omega }^2)_{(apk)}}}`$ (80)
where, $`ap`$ stands for $`(a_tp_0,a_s\stackrel{}{p})`$ and $`K\mathrm{sin}`$ for $`K_0\mathrm{sin}((apk)_0)`$ or for $`K_j\mathrm{sin}((apk)_j)/\chi `$. $`K_\mu (k)`$, $`\mathrm{\Omega }(k)`$, $`V_\mu (k^{},k)`$ and $`D_{\mu \nu }(k,\alpha _g)`$ are detailed in the Appendix. The argument $`\alpha _g`$ in the gauge propagator comes from the gauge fixing term, with $`\alpha _g=1`$ and $`\alpha _g=0`$ corresponding to Feynman and Landau gauges respectively. Our codes have been written for general $`\alpha _g`$ and we have used gauge invariance of $`M_1`$ and $`C_0`$ as one check on our results.
Equation (78) has the familiar form for a self energy integral. The only subtlety is to verify that one is indeed calculating $`\mathrm{\Sigma }^{reg}`$ measured in units of $`1/a_t`$, given the conventions in our Feynman rules. As explained in the Appendix, we choose to work with a dimensionless momentum space gauge propagator $`D_{\mu \nu }`$. It comes from the Fourier transform of the dimensionless correlator $`(a_\mu A_\mu )(a_\nu A_\nu )`$. The relation between $`D_{\mu \nu }`$ and a more conventional propagator of dimension 1/(energy)<sup>2</sup>, denoted $`\stackrel{~}{D}_{\mu \nu }`$, is
$$D_{\mu \nu }=\frac{a_\mu a_\nu }{a_s^3a_t}\stackrel{~}{D}_{\mu \nu }.$$
(82)
Our vertex functions, $`V_\mu `$, are also subtlely different from those in isotropic lattice perturbation theory. They keep tract of the $`1/a_\mu `$ in the derivatives, i.e. of whether one has a $`1/a_s`$ or $`1/a_t`$ there. If $`\stackrel{~}{V}_\mu `$ are vertex functions normalized such that $`\stackrel{~}{V}_\mu i\gamma _\mu `$ for all $`\mu `$ in the continuum limit, then the relation to the $`V_\mu `$ of (78) is given by
$$V_\mu =\frac{a_t}{a_\mu }\stackrel{~}{V}_\mu .$$
(83)
Using (82) and (83) one can write
$`a_t\mathrm{\Sigma }^{reg}(p)`$ $`=`$ $`g^2{\displaystyle \frac{4}{3}}{\displaystyle \underset{\mu ,\nu }{}}{\displaystyle \frac{d^4k}{(2\pi )^4}\left\{\left(\frac{a_t}{a_\mu }\stackrel{~}{V}_\mu \right)\frac{\overline{G}_0}{a_t}\left(\frac{a_t}{a_\nu }\stackrel{~}{V}_\nu \right)\right\}\left(\frac{a_\mu a_\nu }{a_s^3a_t}\stackrel{~}{D}_{\mu \nu }\right)}`$ (84)
$`=`$ $`a_tg^2{\displaystyle \frac{4}{3}}{\displaystyle \underset{\mu ,\nu }{}}{\displaystyle \frac{d^4k}{(2\pi )^4a_s^3a_t}\left\{\stackrel{~}{V}_\mu \overline{G}_0\stackrel{~}{V}_\nu \right\}\stackrel{~}{D}_{\mu \nu }}.`$ (85)
To evaluate (78) we made extensive use of the symbolic manipulation package Mathematica. The integrals themselves were done using the VEGAS program . The various steps involving Mathematica were to 1. calculate the products $`V_\mu [i\gamma K\mathrm{sin}+\mathrm{\Omega }]V_\nu `$; 2. carry out the Dirac traces such as $`tr\{(1+\gamma _0)\mathrm{\Sigma }\}`$ ; 3. take derivatives with respect to external momenta ; 4. put things on the mass shell $`p=(iM_1,\stackrel{}{0})`$ ; and 5. use trigonometric identities to re-express the full integrands in (78) in terms of powers of $`\widehat{k}_\mu 2\mathrm{sin}(k_\mu /2)`$. The last step facilitated speedy evaluation of the integrand by VEGAS.
Both $`M_1`$ and $`C_0`$ are physical quantities. In addition to being gauge invariant they are also IR finite. The wave function renormalization $`Z_2`$, on the other hand, is gauge dependent and also generally logarithmically IR divergent. In any calculation of a physical quantity this IR divergence will eventually be cancelled by vertex corrections and/or matching to continuum operators. In this article we will isolate the gauge dependent IR divergence in $`Z_2`$, verify that it is the same as in the corresponding continuum theory and present results for the remaining IR finite parts. The IR divergence is found in the contribution from $`\mathrm{\Sigma }^{reg}`$ to $`Z_2`$. More specifically it resides in the following term in (64)
$$itr\left(\frac{(\gamma _0+1)}{4}\frac{d}{dp_0}\mathrm{\Sigma }^{reg}(p_0,\stackrel{}{p}=0,m)\right)_{p_0=iM_1^{(0)}}Z_2^{(0)}.$$
(86)
We adopt the method of reference to subtract IR divergent contributions inside integrands and rewrite (86) as
$`{\displaystyle _k}\left\{{\displaystyle \underset{\mu ,\nu }{}}itr\left({\displaystyle \frac{(\gamma _0+1)}{4}}{\displaystyle \frac{d}{dp_0}}\left[V_\mu {\displaystyle \frac{\overline{G}_0}{a_t}}V_\nu \right]_{p_0=iM_1^{(0)}}D_{\mu \nu }Z_2^{(0)}\right)_{sub}(k,m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda )\right\}`$ (87)
$`+F(m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda ),`$ (88)
with
$$F(m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda )=_k_{sub}(k,m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda )$$
(89)
and
$$_kg^2\frac{4}{3}\frac{d^4k}{(2\pi )^4}.$$
(90)
Explicit forms for $`_{sub}(k,m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda )`$ and $`F(m_{\mathrm{𝑒𝑓𝑓}},\mathrm{\Lambda },\lambda )`$ are given in the Appendix. $`\lambda `$ is a gluon mass introduced to regulate IR divengences. $`_{sub}`$ has been constructed to match the same IR divergence as the first term inside the integral in (87). As a result the integral becomes independent of $`\lambda `$. The other condition on $`_{sub}`$ is that the integral (89) be easy to do analytically. The simplest approach is to use a continuum self energy expression for $`_{sub}`$ with an appropriate choice for the mass parameter $`m_{\mathrm{𝑒𝑓𝑓}}`$. The need to adjust $`m_{\mathrm{𝑒𝑓𝑓}}`$ to optimize matching of the IR behaviours in $`_{sub}`$ and the lattice integrand, was emphasized in reference and following that work we find
$$a_tm_{\mathrm{𝑒𝑓𝑓}}=\mathrm{sinh}(a_tM_1^{(0)})\frac{\mathrm{cosh}(a_tM_1^{(0)})+\chi \mathrm{sinh}(a_tM_1^{(0)})}{1+\chi \mathrm{sinh}(a_tM_1^{(0)})}.$$
(91)
The same $`_{sub}`$ and $`m_{\mathrm{𝑒𝑓𝑓}}`$ work for both the clover and D234 quark actions since the IR structure of the two theories agree. Finally, $`\mathrm{\Lambda }\pi `$ in the above expressions is a cutoff imposed on $`_{sub}`$ so that $`_{sub}=0`$ for $`k^2>\mathrm{\Lambda }^2`$. The full expression (87) must be independent of $`\mathrm{\Lambda }`$.
### B $`\mathrm{\Sigma }^{tad}(p)`$
The second contribution to $`\mathrm{\Sigma }(p)`$ is the tadpole contribution $`\mathrm{\Sigma }^{tad}(p)`$ coming from the two-gluon emission vertices listed in the Appendix. For quark action $`𝒮_{D234}^{II}`$ one has
$`a_t\mathrm{\Sigma }^{tad}(p)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left[i\gamma _0\mathrm{sin}(a_tp_0)\chi \mathrm{cos}(a_tp_0)\right]{\displaystyle _k}D_{00}`$ (94)
$`+{\displaystyle \frac{1}{2\chi }}{\displaystyle \frac{1}{3}}{\displaystyle \underset{j}{}}{\displaystyle _k}D_{jj}\{i\gamma _j[(3+C_3)\mathrm{sin}(a_sp_j)2C_3\mathrm{sin}(2a_sp_j)\mathrm{cos}({\displaystyle \frac{k_j}{2}})]`$
$`[(3+C_4)\mathrm{cos}(a_sp_j)C_4\mathrm{cos}(2a_sp_j)\mathrm{cos}^2({\displaystyle \frac{k_j}{2}})]\}.`$
$`\mathrm{\Sigma }^{tad}(p)`$ in the case of $`𝒮_{clover}`$ is obtained by setting $`C_3=C_4=0`$ in the above expression. The appropriate traces and derivatives with respect to external momenta can be carried out immediately and one has
$`tr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}a_t\mathrm{\Sigma }^{tad}\right\}_{p=(iM_1^{(0)},\stackrel{}{0})}`$ (95)
$`=\{\begin{array}{c}\frac{1}{2}[\mathrm{sinh}(a_tM_1^{(0)})+\chi \mathrm{cosh}(a_tM_1^{(0)})]_kD_{00}+\frac{1}{2\chi }_j_kD_{jj}𝒮_{clover}\hfill \\ \\ \frac{1}{2}[\mathrm{sinh}(a_tM_1^{(0)})+\chi \mathrm{cosh}(a_tM_1^{(0)})]_kD_{00}+\frac{1}{6\chi }_j_kD_{jj}[4\mathrm{cos}^2(\frac{k_j}{2})]𝒮_{D234}^{II}\hfill \end{array}`$ (99)
(100)
(101)
$`itr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}{\displaystyle \frac{d}{dp_0}}\mathrm{\Sigma }^{tad}\right\}_{p=(iM_1^{(0)},\stackrel{}{0})}={\displaystyle \frac{1}{2}}[\mathrm{cosh}(a_tM_1^{(0)})+\chi \mathrm{sinh}(a_tM_1^{(0)})]{\displaystyle _k}D_{00}`$ (102)
$`𝒮_{clover}\&𝒮_{D234}^{II}`$ (103)
(104)
(105)
$`B_j^{tad}B_0^{tad}=\{\begin{array}{c}\frac{1}{2}_kD_{jj}\frac{1}{2}_kD_{00}𝒮_{clover}\hfill \\ \\ \frac{2}{3}_kD_{jj}\mathrm{sin}^2(\frac{k_j}{2})\frac{1}{2}_kD_{00}𝒮_{D234}^{II}\hfill \end{array}`$ (109)
where $`B_j^{tad}`$ and $`B_0^{tad}`$ are the contributions from the tadpole diagram to (73) and (74) or (76). All the integrals are IR finite and very easy to carry out numerically. Contributions from $`\mathrm{\Sigma }^{tad}`$ typically dominate over those from $`\mathrm{\Sigma }^{reg}`$ but the bulk if not all of it is cancelled by $`\mathrm{\Sigma }^{t.i.}`$.
### C $`\mathrm{\Sigma }^{t.i.}(p)`$
The lattice covariant derivatives in the quark actions are tadpole-improved. They are listed in the Appendix. In momentum space one has, for instance
$$_\mu i\mathrm{sin}(a_\mu p_\mu )/u_\mu i\mathrm{sin}(a_\mu p_\mu )[1+\alpha _su_\mu ^{(2)}]+O(\alpha _s^2),$$
(110)
where we have perturbatively expanded
$$u_\mu =1\alpha _su_\mu ^{(2)}+O(\alpha _s^2).$$
(111)
Even in the absence of the regular and tadpole one-loop diagrams there are hence $`O(\alpha _s)`$ terms in the quark propagator. We denote the inverse quark propagator with the $`u_\mu `$’s still in place as $`\overline{G}_{0,u0}^1(p)`$, so that $`\overline{G}_0^1(p)\overline{G}_{0,u0=1}^1(p)`$. Through $`O(\alpha _s)`$ eqn.(20) can be written as
$$\overline{G}^1=\overline{G}_{0,u0}^1\mathrm{\Sigma }^{reg}\mathrm{\Sigma }^{tad}\overline{G}_0^1\mathrm{\Sigma }^{reg}\mathrm{\Sigma }^{tad}\mathrm{\Sigma }^{t.i.}$$
(112)
or
$$\mathrm{\Sigma }^{t.i.}=\overline{G}_{0,u0=1}^1\overline{G}_{0,u0}^1.$$
(113)
From the difference in (113) one sees that one link hops bring in factors of $`11/u_\mu \alpha _su_\mu ^{(2)}`$ and two link hops factors of $`11/u_\mu ^22\alpha _su_\mu ^{(2)}`$ etc. Using these rules one finds
$`a_t\mathrm{\Sigma }^{t.i.}(p)`$ $`=`$ $`\alpha _su_t^{(2)}[i\gamma _0\mathrm{sin}(a_tp_0)+\chi \mathrm{cos}(a_tp_0)]+`$ (116)
$`\alpha _su_s^{(2)}{\displaystyle \frac{1}{3\chi }}{\displaystyle \underset{j}{}}\{i\gamma _j[(3+C_3)\mathrm{sin}(a_sp_j)C_3\mathrm{sin}(2a_sp_j)]`$
$`+[(3+C_4)\mathrm{cos}(a_sp_j){\displaystyle \frac{C_4}{2}}\mathrm{cos}(2a_sp_j)]\}.`$
The relevant traces and derivatives become
$`tr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}a_t\mathrm{\Sigma }^{t.i.}\right\}_{p=(iM_1^{(0)},\stackrel{}{0})}`$ (125)
$`=\{\begin{array}{c}[\mathrm{sinh}(a_tM_1^{(0)})+\chi \mathrm{cosh}(a_tM_1^{(0)})]\alpha _su_t^{(2)}\frac{3}{\chi }\alpha _su_s^{(2)}𝒮_{clover}\hfill \\ \\ [\mathrm{sinh}(a_tM_1^{(0)})+\chi \mathrm{cosh}(a_tM_1^{(0)})]\alpha _su_t^{(2)}\frac{1}{\chi }\frac{7}{2}\alpha _su_s^{(2)}𝒮_{D234}^{II}\hfill \end{array}`$
$`itr\left\{{\displaystyle \frac{(\gamma _0+1)}{4}}{\displaystyle \frac{d}{dp_0}}\mathrm{\Sigma }^{t.i.}\right\}_{p=(iM_1^{(0)},\stackrel{}{0})}=[\mathrm{cosh}(a_tM_1^{(0)})+\chi \mathrm{sinh}(a_tM_1^{(0)})]\alpha _su_t^{(2)}`$
$`𝒮_{clover}\&𝒮_{D234}^{II}`$
$`B_j^{t.i.}B_0^{t.i.}=\{\begin{array}{c}\alpha _s(u_t^{(2)}u_s^{(2)})𝒮_{clover}\hfill \\ \\ \alpha _s(u_t^{(2)}\frac{2}{3}u_s^{(2)})𝒮_{D234}^{II}\hfill \end{array}`$ (131)
The Landau mean link definition of $`u_\mu `$ is given by
$$u_\mu \frac{1}{3}TrU_\mu _{\alpha _g=0}1\alpha _su_\mu ^{(2)}=1\frac{1}{2}_kD_{\mu \mu }(\alpha _g=0).$$
(132)
If one evaluates $`\mathrm{\Sigma }^{tad}`$ in Landau gauge then (95) & (125) , (102) & (125) and (109) & (131) cancel for $`𝒮_{clover}`$. (102) & (125) also cancel for $`𝒮_{D234}^{II}`$ and for the other two traces cancellation is almost complete. The difference between contributions from $`\mathrm{\Sigma }^{tad}`$ and $`\mathrm{\Sigma }^{t.i.}`$ would go away if one replaces $`\mathrm{cos}^2(k/2)`$ and $`\mathrm{sin}^2(k/2)`$ by their averages $`1/2`$. Hence, it is easy to see in this calculation how tadpole improving terms in the lattice action eliminates lattice artifact contributions in perturbation theory.
## V Results
In this section we summarize results for the one-loop coefficients, $`a_sM_{1,sub}^{(1)}`$, $`Z_2^{(1)}`$ and $`C_0^{(1)}`$ for mass, wave function and speed of light renormalizations respectively. These follow from equations (53), (64) and (72) - (76) and each has, as explained in the previous section, contributions from regular, tadpole and t.i. diagrams. The numbers in our Tables are coefficients of $`\alpha _s`$. The Landau mean link definition of $`u_\mu `$ is used throughout to implement tadpole improvement.
### A $`a_sM_{1,sub}^{(1)}`$
In Table II we present results for $`a_sM_{1,sub}^{(1)}`$ for action $`𝒮^A`$ for several values of $`a_sM_1^{(0)}`$. We list separately contributions from $`\mathrm{\Sigma }^{reg}`$, $`\mathrm{\Sigma }^{tad}`$ and $`\mathrm{\Sigma }^{t.i.}`$. The fourth column gives the gauge invariant combination (reg + tad) and the sixth column gives $`a_sM_{1,nosub}^{(1)}`$ (reg + tad + t.i.), the full tadpole improved one-loop correction before subtraction. Carrying out the subtraction according to (53), one obtains $`a_sM_{1,sub}^{(1)}`$ which is given in the last column
$$a_sM_{1,sub}^{(1)}=a_sM_{1,nosub}^{(1)}\frac{a_sM_{1,nosub}^{(1)}(0)}{\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})}.$$
(133)
All our calculations have been carried out for two values of the gauge fixing parameter $`\alpha _g`$, 1.0 and 0.0. Table II lists both sets of results and one sees that gauge invariant quantities are independent of $`\alpha _g`$ within numerical integration errors ( which we take to be at the $`\pm 0.003`$ to $`\pm 0.006`$ level depending on the mass). Our results for $`a_sM_1^{(0)}=0`$ agree with those from .
We plot $`a_sM_{1,sub}^{(1)}`$ versus $`a_sM_1^{(0)}`$ in Fig. 1 . One sees that the mass dependence is smooth and that one reaches saturation rapidly already around $`a_sM_1^{(0)}3.05.0`$. We also compare with non-tadpole improved results for which the curve saturates around 1.827 rather than around 1.077. In considering the large mass limit it is useful to note that the factor $`Z_2^{(0)}=1/[\chi \mathrm{sinh}(a_tM_1^{(0)})+\mathrm{cosh}(a_tM_1^{(0)})]`$ appearing in (53) and (64) vanishes exponentially in this limit. The only terms that survive into the static limit are those where $`Z_2^{(0)}`$ is multiplied by an exponentially increasing function of $`a_sM_1^{(0)}`$. It is easy to see, for instance, that the subtraction term in (133) or the spatial tadpoles in (95) become irrelevant in the static limit. Furthermore $`a_sM_{1,sub}^{(1)}`$ becomes identical for $`𝒮_{Wilson}`$, $`𝒮_{clover}`$ and $`𝒮_{D234}^{II}`$ in this limit and the only difference between the current calculations and those of reference resides in the glue action (we have verified that by switching to the unimproved Wilson glue action results of are reproduced).
In Tables III, IV and V we summarize results for $`a_sM_{1,sub}^{(1)}`$ for the other actions listed in Table I. We also list the combination (reg + tad) for each case. If one chooses to implement tadpole improvement differently from what we have done here or decides not to tadpole improve, then $`a_sM_{1,sub}^{(1)}`$ can be calculated straightforwardly from (reg + tad) and the formulas presented in this paper. For action $`𝒮^C`$ we list only massless results for reasons explained in section II. The dispersion relations of this and similar actions are discussed in reference . Our choices for anisotropy values in Tables IV and V, were dictated in part with an eye towards practical numerical simulations. Values such as $`\chi =3.6`$ and $`\chi =5.3`$ were taken from recent work on nonperturbative determinations of the renormalized anisotropy in pure glue theory . For one value $`a_sM_1^{(0)}=1`$ we plot $`a_sM_{1,sub}^{(1)}`$ versus $`\chi `$ in Fig. 2 using the action $`𝒮^D`$. From Tables IV and V and from Fig. 2 one sees that the dependence of $`a_sM_{1,sub}^{(1)}`$ on $`\chi `$ is very mild. Effects of tadpole improvement are significant only for small values of $`\chi `$. This is because due to cancellations in the subtraction of equation (133) only the temporal tadpole and the temporal Landau link term $`u_t^{(2)}`$ contribute to $`a_sM_{1,sub}^{(1)}`$ and both these become small as $`\chi `$ increases.
Finally we mention that the one-loop expression for the critical bare mass $`m_c`$ is given by
$$a_tm_c=\frac{\alpha _s}{\chi }a_sM_{1,nosub}^{(1)}(0)+O(\alpha _s^2).$$
(134)
### B $`Z_2^{(1)}`$
Starting with (64) we define
$$Z_2=Z_2^{(0)}[\mathrm{\hspace{0.17em}1}+\alpha _s(Z_2^{(1)}+Z_2^{(1)IR})+O(\alpha _s^2)]$$
(135)
with
$$Z_2^{(1)IR}=\{\begin{array}{c}\frac{1}{3\pi }[\mathrm{\hspace{0.33em}\hspace{0.33em}1}+(\alpha _g1)]\mathrm{ln}(\lambda ^2)m=0\hfill \\ \\ \frac{1}{3\pi }[2+(\alpha _g1)]\mathrm{ln}(\lambda ^2)m>0\hfill \end{array}$$
(136)
$`\lambda `$ is the gluon mass in units of $`1/a_s`$. It is the coefficient of $`\alpha _s`$ after factoring out $`Z_2^{(0)}`$ that has the same IR $`\mathrm{ln}(\lambda )`$ and $`\mathrm{ln}(am)`$ structure as the continuum wave function renormalization constant. From (64) one also sees that there are two contributions to $`Z_2^{(1)}`$, one coming from the $`d/dp_0`$ derivative term and the second from the expansion of $`M_1`$. Accordingly we write
$$Z_2^{(1)}=Z_{2,dp_0}^{(1)}+Z_{2,M_1}^{(1)}.$$
(137)
In the literature $`Z_{2,M_1}^{(1)}`$ is not always included as part of the definition of $`Z_2^{(1)}`$. $`Z_{2,dp_0}^{(1)}`$ alone with unimproved Wilson glue goes over in the large mass limit to the wave function renormalization of reference . Including $`Z_{2,M_1}^{(1)}`$ leads to the static result of reference which has been used in many subsequent static calculations, for instance in . This latter static value also corresponds to the large mass limit of the one-loop $`Z_2`$ calculated in many versions of NRQCD actions .
Table VI presents results for $`Z_{2,dp_0}^{(1)}`$ and the full $`Z_2^{(1)}`$ for the action $`𝒮^A`$. Again we agree with reference for $`m=0`$. However, one notices that the massive data do not tend towards the massless result as $`a_sM_1^{(0)}`$ decreases. This is because our massive numbers include $`\mathrm{ln}(am)`$ contributions which will eventually diverge, whereas in the massless theory we have set the fermion mass identical to zero from the beginning. This leads to different IR structure for the two theories (see Appendix B for some further discussions). In a matching calculation one will be looking at differences between the lattice and continuum $`Z_2`$. As long as IR divergences are handled in the same manner in the lattice and continuum evaluations, one should not run into any problems and the $`m0`$ limit should be smooth. For instance, using dimensional regularization in the $`\overline{MS}`$ scheme one finds in Feynman gauge the UV finite continuum results
$$Z_2^{(1)cont.}=\{\begin{array}{c}\frac{1}{3\pi }\left[\mathrm{ln}(\frac{\widehat{\lambda }^2}{\mu ^2})+\frac{1}{2}\right]m=0\hfill \\ \\ \frac{1}{3\pi }\left[\mathrm{ln}(\frac{m^2}{\mu ^2})+2\mathrm{ln}(\frac{m^2}{\widehat{\lambda }^2})4\right]m>0\hfill \end{array}$$
(138)
Taking the difference between continuum and lattice wave function renormalization constants, it makes sense to consider the following subtracted $`Z_2`$ factors.
$$Z_{2,diff}^{(1)}=Z_{2,dp_0}^{(1)}\{\begin{array}{c}\frac{1}{3\pi }\frac{1}{2}m=0\hfill \\ \\ \frac{1}{3\pi }[\mathrm{\hspace{0.17em}3}\mathrm{ln}(a_sM_1^{(0)})^24]m>0\hfill \end{array}$$
(139)
Numbers for $`Z_{2,diff}^{(1)}`$ are given in Table VII and one sees that the $`m0`$ behaviour is smooth.
In Tables VIII through XII we present $`Z_{2,dp_0}^{(1)}`$ and $`Z_2^{(1)}`$ for other actions. One does not find any dramatic changes with differing actions and/or anisotropy. The IR subtractions of Appendix B worked well in all actions for $`a_sM_1^{(0)}<5`$. For larger masses VEGAS errors became large especially for $`\chi =1`$. Hence, we only present results up to $`a_sM_1^{(0)}=5`$. For $`\chi >1`$ problems were less severe in general. In Tables VI - XII the numerical integration errors are at the $`\pm 0.02`$ level for $`a_sM_1^{(0)}=5.0`$ and $`\chi <3`$, at the $`\pm 0.006`$ level for $`a_sM_1^{(0)}`$ = 0.01, 0.05 and 0.10 and at the $`\pm 0.004`$ level for all other cases. A more sophisticated method for handling IR divergent integrals appears necessary if accurate results are required for larger masses. Many quantities, however, are close to saturation by the time one reaches $`a_sM_1^{(0)}=5`$.
### C $`C_0^{(1)}`$
For the speed of light renormalization we present the most detailed results for action $`𝒮^B`$ rather than for $`𝒮^A`$ since the former, for $`\chi >1`$, is genuinely anisotropic. In Table XIII we list separately $`regular`$ and $`tadpole`$ diagram contributions, their gauge invariant sum $`(reg+tad)(C_0^{(1)}not.i.)`$ and the fully tadpole improved result $`(C_0^{(1)}witht.i.)`$, all for action $`𝒮^B`$ and at fixed anisotropy $`\chi =4.0`$. $`C_0^{(1)}`$ is independent of $`\alpha _g`$ within numerical integration errors which are the most severe when using (74) for nonzero but small masses. In Figure 3. we plot $`C_0^{(1)}`$ both with and without tadpole improvement versus $`a_sM_1^{(0)}`$. Table XIV summarizes results for several $`\chi `$ values with action $`𝒮^B`$ and in Figure 4. we plot $`C_0^{(1)}`$ versus $`\chi `$ for fixed $`a_sM_1^{(0)}=1.0`$. One sees that tadpole improvement has significant effect and causes $`C_0^{(1)}`$ to switch sign for $`\chi >1`$. Among other things this allows for a smooth $`\chi 1`$ limit.
In Table XV we present results for actions $`𝒮^A`$ and $`𝒮^A^{}`$. In these isotropic actions nontrivial $`C_0`$ comes about because our mass-shell condition $`p=(iM_1,\stackrel{}{0})`$ distinguishes between spatial and temporal directions once $`M_1>0`$. Table XVI summarizes results for action $`𝒮^D`$. Here tadpole improvement does not decrease the magnitude of the one-loop correction, however, for a wide range of mass values it is still true that $`C_0^{(1)}`$ switches sign for $`\chi >1`$ and that the $`\chi 1`$ limit becomes smoother after tadpole improvement.
## VI Summary
We have carried out one-loop perturbative renormalization of the fermion rest mass $`M_1`$, wave function renormalization $`Z_2`$ and the speed of light $`C_0`$ for a range of highly improved actions on isotropic and anisotropic lattices. We find that the dependence of the one-loop coefficients on the anisotropy parameter $`\chi =a_s/a_t`$ and on the tree-level mass parameter $`a_sM_1^{(0)}`$ is mild, especially after tadpole improvement of the actions. Furthermore, none of the coefficients are particularly large. $`M_1`$ and $`C_0`$ exhibit smooth behaviour as one approaches the massless, large mass, $`\chi 1`$ and large $`\chi `$ limits. This also holds for $`Z_2`$ if more physical combinations such as the difference between continuum and lattice wave function renormalizations are considered. The next stage in our program would be to extend the present calculations to vertex corrections and to matchings between continuum and lattice currents and other multi-fermion operators.
## VII Acknowledgments
This research is supported by grants from the US Department of Energy, DE-FG02-91ER40690, and from the Graduiertenkolleg. The authors thank Peter Lepage for many useful conversations. S.G. gratefully acknowledges a grant given by the Max Kade Foundation. J.S. thanks Cornell University for its hospitality while the project was being initiated and the theoretical physics group at the University of Glasgow where part of the work was carried out. Support from a UK PPARC Visiting Fellowship PPA/V/S/1997/00666 is gratefully acknowledged.
## A Definitions and Feynman rules
In this Appendix we summarize definitions for various terms in the lattice actions and present Feynman rules for gauge and quark propagators and for vertex functions.
Covariant Derivatives Acting on Quark Fields
$`_\mu \mathrm{\Psi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{u_\mu }}[U_\mu (x)\mathrm{\Psi }(x+a_\mu )U_\mu ^{}(xa_\mu )\mathrm{\Psi }(xa_\mu )]`$ (A1)
$`_\mu ^{(2)}\mathrm{\Psi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{u_\mu }}[U_\mu (x)\mathrm{\Psi }(x+a_\mu )+U_\mu ^{}(xa_\mu )\mathrm{\Psi }(xa_\mu )]\mathrm{\hspace{0.33em}2}\mathrm{\Psi }(x)`$ (A2)
$`_\mu ^{(3)}\mathrm{\Psi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{u_\mu ^2}}[U_\mu (x)U_\mu (x+a_\mu )\mathrm{\Psi }(x+2a_\mu )U_\mu ^{}(xa_\mu )U_\mu ^{}(x2a_\mu )\mathrm{\Psi }(x2a_\mu )]`$ (A4)
$`{\displaystyle \frac{1}{u_\mu }}[U_\mu (x)\mathrm{\Psi }(x+a_\mu )U_\mu ^{}(xa_\mu )\mathrm{\Psi }(xa_\mu )]`$
$`_\mu ^{(4)}\mathrm{\Psi }(x)`$ $`=`$ $`{\displaystyle \frac{1}{u_\mu ^2}}[U_\mu (x)U_\mu (x+a_\mu )\mathrm{\Psi }(x+2a_\mu )+U_\mu ^{}(xa_\mu )U_\mu ^{}(x2a_\mu )\mathrm{\Psi }(x2a_\mu )]`$ (A6)
$`\mathrm{\hspace{0.17em}4}{\displaystyle \frac{1}{u_\mu }}[U_\mu (x)\mathrm{\Psi }(x+a_\mu )+U_\mu ^{}(xa_\mu )\mathrm{\Psi }(xa_\mu )]+\mathrm{\hspace{0.33em}6}\mathrm{\Psi }(x)`$
Field Strength Tensors
For the unimproved $`F_{\mu \nu }`$ of the clover action we use
$`F_{\mu \nu }(x)`$ $`=`$ $`{\displaystyle \frac{1}{2i}}\left(\mathrm{\Omega }_{\mu \nu }(x)\mathrm{\Omega }_{\mu \nu }^{}(x)\right),`$ (A7)
$`\mathrm{\Omega }_{\mu \nu }(x)`$ $`=`$ $`{\displaystyle \frac{1}{4u_\mu ^2u_\nu ^2}}{\displaystyle \underset{\{(\alpha ,\beta )\}_{\mu \nu }}{}}U_\alpha (x)U_\beta (x+a_\alpha )U_\alpha (x+a_\alpha +a_\beta )U_\beta (x+a_\beta ),`$ (A8)
with $`\{(\alpha ,\beta )\}_{\mu \nu }=\{(\mu ,\nu ),(\nu ,\mu ),(\mu ,\nu ),(\nu ,\mu )\}`$ for $`\mu \nu `$ and $`U_\mu (x+a_\mu )U_\mu ^{}(x)`$. The $`O(a^2)`$ improved field strength tensor of the D234 actions is
$`\stackrel{~}{F}_{\mu \nu }(x)={\displaystyle \frac{5}{3}}F_{\mu \nu }(x)`$ (A9)
$`{\displaystyle \frac{1}{6}}[{\displaystyle \frac{1}{u_\mu ^2}}(U_\mu (x)F_{\mu \nu }(x+a_\mu )U_\mu ^{}(x)+U_\mu ^{}(xa_\mu )F_{\mu \nu }(xa_\mu )U_\mu (xa_\mu ))(\mu \nu )]`$ (A10)
$`+{\displaystyle \frac{1}{6}}({\displaystyle \frac{1}{u_\mu ^2}}+{\displaystyle \frac{1}{u_\nu ^2}}2)F_{\mu \nu }(x).`$ (A11)
The last term ensures that factors of $`1/u_\mu `$ are correctly removed from those contributions to $`UF_{\mu \nu }U^{}`$ and $`U^{}F_{\mu \nu }U`$ that end up being four link objects rather than six link ones. In a one-loop calculation, however, one can set $`u_\mu =1`$ everywhere in the definition of the field strength tensor and this correction term is irrelevant.
Both the above covariant derivatives and the field strength tensor are dimensionless. Factors of $`1/a_t`$ and $`1/a_s`$ are inserted explicitly where necessary such as in (10).
Gauge Propagator
The isotropic Symanzik improved gauge action has been discussed quite extensively in the literature . Here we summarize formulas for the anisotropic generalization. We start from the gauge actions $`𝒮_G^I`$ or $`𝒮_G^{II}`$ and add to it a gauge fixing term
$`𝒮_{gf}`$ $`=`$ $`{\displaystyle \frac{1}{2\alpha _g}}a_s^3a_t{\displaystyle \underset{x}{}}\left[{\displaystyle \frac{1}{a_t}}_tA_t+{\displaystyle \frac{1}{a_s}}{\displaystyle \underset{j}{}}_jA_j\right]^2`$ (A12)
$`=`$ $`{\displaystyle \frac{1}{2\alpha _g}}{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{x}{}}\left[\chi ^2_t(a_tA_t)+{\displaystyle \underset{j}{}}_j(a_sA_j)\right]^2,`$ (A13)
with $`_\mu A_\mu (x)A_\mu (x+a_\mu /2)A_\mu (xa_\mu /2)`$. Equation (A13) expresses $`𝒮_{gf}`$ in terms of dimensionless gauge fields $`a_\mu A_\mu `$. It is convenient to do so, especially since $`𝒮_G^{I,II}`$ are already in dimensionless form with factors of $`\chi `$ and $`1/\chi `$ properly put in place. If $`\overline{A}_\mu (k)`$ is the Fourier transform of $`(a_\mu A_\mu )`$, the quadratic terms in the gauge action become
$$𝒮_G^{(0)I,II}+𝒮_{gf}=\frac{1}{2}\underset{\mu \nu }{}_\pi ^\pi \frac{d^4k}{(2\pi )^4}\left(\overline{A}_\mu (k)M_{\mu \nu }(k)\overline{A}_\nu (k)\right),$$
(A14)
where
$`M_{00}`$ $`=`$ $`\chi \left[{\displaystyle \frac{\chi ^2}{\alpha _g}}\widehat{k}_0^2+{\displaystyle \underset{j}{}}\widehat{k}_j^2q_{0j}\right]`$ (A15)
$`M_{jj}`$ $`=`$ $`{\displaystyle \frac{1}{\chi }}\left[{\displaystyle \frac{1}{\alpha _g}}\widehat{k}_j^2+\chi ^2\widehat{k}_0^2q_{0j}+{\displaystyle \underset{lj}{}}\widehat{k}_l^2q_{lj}\right]`$ (A16)
$`M_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{\chi }}\left[{\displaystyle \frac{1}{\alpha _g}}\widehat{k}_i\widehat{k}_j\widehat{k}_i\widehat{k}_jq_{ij}\right]`$ (A17)
$`M_{0j}=M_{j0}`$ $`=`$ $`\chi \left[{\displaystyle \frac{1}{\alpha _g}}\widehat{k}_0\widehat{k}_j\widehat{k}_0\widehat{k}_jq_{0j}\right]`$ (A18)
and
$$\widehat{k}_\mu 2\mathrm{sin}(\frac{k_\mu }{2}).$$
(A19)
The $`q_{\mu \nu }`$ need to be specified only for $`\mu \nu `$ and one has
$`q_{\mu \nu }`$ $`=`$ $`1c_1^G(\widehat{k}_\mu ^2+\widehat{k}_\nu ^2)\mu \nu 𝒮_G^I`$ (A20)
$`q_{ij}`$ $`=`$ $`1+{\displaystyle \frac{1}{12}}(\widehat{k}_i^2+\widehat{k}_j^2)ij𝒮_G^{II}`$ (A22)
$`q_{0j}`$ $`=`$ $`1+{\displaystyle \frac{1}{12}}\widehat{k}_j^2𝒮_G^{II}`$ (A23)
We have inverted the $`4\times 4`$ matrix $`M_{\mu \nu }`$ using Mathematica keeping $`q_{\mu \nu }`$ general. For both gauge actions, $`𝒮_G^I`$ and $`𝒮_G^{II}`$ the free gauge propagator has the structure
$$D_{\mu \nu }(k)=M_{\mu \nu }^1=\frac{1}{(\widehat{k}^2)^2}\left[\alpha _g\widehat{k}_\mu \widehat{k}_\nu \chi +\frac{f_N^{\mu \nu }(\widehat{k}_\rho ,q_{\rho \sigma },\chi )}{f_D(\widehat{k}_\rho ,q_{\rho \sigma },\chi )}\right],$$
(A24)
with
$$\widehat{k}^2=\chi ^2\widehat{k}_0^2+\underset{j}{}\widehat{k}_j^2.$$
(A25)
The term proportional to the gauge fixing parameter $`\alpha _g`$ has the familiar form
$$\alpha _g\frac{\widehat{k}_\mu \widehat{k}_\nu }{(\widehat{k}^2)^2}\chi =\alpha _g\frac{a_\mu a_\nu }{a_s^3a_t}\frac{\widehat{k}_\mu \widehat{k}_\nu /(a_\mu a_\nu )}{[(\widehat{k}_0/a_t)^2+_j(\widehat{k}_j/a_s)^2]^2}$$
(A26)
with the conversion factor $`a_\mu a_\nu /a_s^3a_t`$ mentioned in (82). This factor results because we are looking at the propagator for dimensionless gauge fields $`a_\mu A_\mu `$ and because we carried out a dimensionless Fourier transform. The second term in (A24) is much more complicated. If one writes
$`f_N^{00}(\widehat{k}_\rho ,q_{\rho \sigma },\chi )`$ $`=`$ $`{\displaystyle \frac{1}{\chi }}\stackrel{~}{f}_N^{00}`$ (A27)
$`f_N^{jj}(\widehat{k}_\rho ,q_{\rho \sigma },\chi )`$ $`=`$ $`\chi \stackrel{~}{f}_N^{jj}`$ (A28)
$`f_N^{ij}(\widehat{k}_\rho ,q_{\rho \sigma },\chi )`$ $`=`$ $`\chi \widehat{k}_i\widehat{k}_j\stackrel{~}{f}_N^{ij}`$ (A29)
$`f_N^{0j}(\widehat{k}_\rho ,q_{\rho \sigma },\chi )`$ $`=`$ $`\chi \widehat{k}_0\widehat{k}_j\stackrel{~}{f}_N^{0j}`$ (A30)
one can show that $`f_D`$ and all the $`\stackrel{~}{f}_N^{\mu \nu }`$ are functions only of $`(\chi \widehat{k}_0)^2,\widehat{k}_j^2,q_{\rho \sigma }`$ with no other $`\chi `$ dependence or odd powers of $`\widehat{k}_\rho `$. We have not shown color indices in the above expressions. The gluon propagator is diagonal in color.
Quark Propagator
The inverse free quark propagator for $`𝒮_{D234}^I`$ is given by
$$a_t\overline{G}_0^1(k)=i\gamma _0K_0(k_0)\mathrm{sin}(k_0)+i\frac{C_0}{\chi }\underset{j}{}\gamma _jK_j(k_j)\mathrm{sin}(k_j)+\mathrm{\Omega }(k_0,\stackrel{}{k})$$
(A31)
with
$`K_0`$ $`=`$ $`1+{\displaystyle \frac{C_{3t}}{3}}{\displaystyle \frac{C_{3t}}{3}}\mathrm{cos}(k_0)`$ (A32)
$`K_j`$ $`=`$ $`1+{\displaystyle \frac{C_3}{3}}{\displaystyle \frac{C_3}{3}}\mathrm{cos}(k_j)`$ (A33)
and
$`\mathrm{\Omega }`$ $`=`$ $`\chi \left[2(1+{\displaystyle \frac{C_{4t}}{3}})\mathrm{sin}^2({\displaystyle \frac{k_0}{2}}){\displaystyle \frac{C_{4t}}{6}}\mathrm{sin}^2(k_0)\right]`$ (A35)
$`{\displaystyle \frac{1}{\chi }}{\displaystyle \underset{j}{}}\left[2(1+{\displaystyle \frac{C_4}{3}})\mathrm{sin}^2({\displaystyle \frac{k_j}{2}}){\displaystyle \frac{C_4}{6}}\mathrm{sin}^2(k_j)\right]+a_tm.`$
Propagators for the other quark actions can be obtained by setting the appropriate $`C_{i(t)}`$ equal to zero. Quark propagators are diagonal in color.
Vertex Functions
In deriving the one- and two-gluon emission vertices we have used the method described in . We list again results only for $`𝒮_{D234}^I`$. Those for other quark actions follow trivially. The general form for a single gluon emission vertex is
$$V_\mu (k^{},k)i\gamma _\mu W_\mu W_\mu ^{}\underset{\nu }{}\sigma _{\nu \mu }W_{\nu \mu }^{\prime \prime }$$
(A36)
where $`\mu `$ is the polarization of the emitted gluon, $`k^{}`$ the momentum of the outgoing quark and $`k`$ the momentum of the incoming quark. We suppress the color factor $`T_{bc}^a`$ which should multiply each of the above terms. Using the variables
$$k_\mu ^\pm \frac{1}{2}(k^{}\pm k)_\mu $$
(A37)
one has
$`W_0`$ $`=`$ $`(1+{\displaystyle \frac{C_{3t}}{3}})\mathrm{cos}(k_0^+){\displaystyle \frac{C_{3t}}{3}}\mathrm{cos}(2k_0^+)\mathrm{cos}(k_0^{})`$ (A38)
$`W_j`$ $`=`$ $`{\displaystyle \frac{C_0}{\chi }}\left[(1+{\displaystyle \frac{C_3}{3}})\mathrm{cos}(k_j^+){\displaystyle \frac{C_3}{3}}\mathrm{cos}(2k_j^+)\mathrm{cos}(k_j^{})\right]`$ (A39)
$`W_0^{}`$ $`=`$ $`\chi \left[(1+{\displaystyle \frac{C_{4t}}{3}})\mathrm{sin}(k_0^+){\displaystyle \frac{C_{4t}}{6}}\mathrm{sin}(2k_0^+)\mathrm{cos}(k_0^{})\right]`$ (A40)
$`W_j^{}`$ $`=`$ $`{\displaystyle \frac{1}{\chi }}\left[(1+{\displaystyle \frac{C_4}{3}})\mathrm{sin}(k_j^+){\displaystyle \frac{C_4}{6}}\mathrm{sin}(2k_j^+)\mathrm{cos}(k_j^{})\right]`$ (A41)
and
$`W_{j0}^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}(2k_j^{})\mathrm{cos}(k_0^{}){\displaystyle \frac{1}{3}}[5\mathrm{cos}(2k_j^{})\mathrm{cos}(2k_0^{})]`$ (A42)
$`W_{0j}^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}(2k_0^{})\mathrm{cos}(k_j^{}){\displaystyle \frac{1}{3}}[5\mathrm{cos}(2k_j^{})\mathrm{cos}(2k_0^{})]`$ (A43)
$`W_{ij}^{\prime \prime }`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \frac{1}{\chi }}\mathrm{sin}(2k_i^{})\mathrm{cos}(k_j^{}){\displaystyle \frac{1}{3}}[5\mathrm{cos}(2k_i^{})\mathrm{cos}(2k_j^{})].`$ (A44)
For the clover action the factor $`\frac{1}{3}[5\mathrm{cos}(2k_\mu ^{})\mathrm{cos}(2k_\nu ^{})]`$ in $`W_{\mu \nu }^{\prime \prime }`$ should be replaced by $`1`$.
For the two-gluon emission vertex we do not present the most general result, but restrict ourselves to those terms necessary for the tadpole diagram $`\mathrm{\Sigma }^{tad}`$. For instance, the $`\sigma _{\mu \nu }F_{\mu \nu }`$ term does not contribute to the tadpole diagram. We also omit terms that vanish upon symmetrizing between the two gluons. If $`V_{\mu _1\mu _2}^{(2)}(k^{},k,q_1,q_2)`$ stands for the emission vertex for gluons of momentum $`q_i`$ and polarization $`\mu _i`$, with $`k_\mu =k_\mu ^{}+q_{1,\mu }+q_{2,\mu }`$, one has
$`V_{00}^{(2)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}\gamma _0\left[(1+{\displaystyle \frac{C_{3t}}{3}})\mathrm{sin}(k_0^+){\displaystyle \frac{2}{3}}C_{3t}\mathrm{sin}(2k_0^+)\mathrm{cos}({\displaystyle \frac{q_{1,0}}{2}})\mathrm{cos}({\displaystyle \frac{q_{2,0}}{2}})\right]`$ (A46)
$`{\displaystyle \frac{\chi }{2}}\left[(1+{\displaystyle \frac{C_{4t}}{3}})\mathrm{cos}(k_0^+){\displaystyle \frac{1}{3}}C_{4t}\mathrm{cos}(2k_0^+)\mathrm{cos}({\displaystyle \frac{q_{1,0}}{2}})\mathrm{cos}({\displaystyle \frac{q_{2,0}}{2}})\right]`$
$`V_{jj}^{(2)}`$ $`=`$ $`{\displaystyle \frac{i}{2}}{\displaystyle \frac{C_0}{\chi }}\gamma _j\left[(1+{\displaystyle \frac{C_3}{3}})\mathrm{sin}(k_j^+){\displaystyle \frac{2}{3}}C_3\mathrm{sin}(2k_j^+)\mathrm{cos}({\displaystyle \frac{q_{1,j}}{2}})\mathrm{cos}({\displaystyle \frac{q_{2,j}}{2}})\right]`$ (A49)
$`{\displaystyle \frac{1}{2\chi }}\left[(1+{\displaystyle \frac{C_4}{3}})\mathrm{cos}(k_j^+){\displaystyle \frac{1}{3}}C_4\mathrm{cos}(2k_j^+)\mathrm{cos}({\displaystyle \frac{q_{1,j}}{2}})\mathrm{cos}({\displaystyle \frac{q_{2,j}}{2}})\right].`$
The color factor for these vertex functions is $`(T^{a_1}T^{a_2})_{bc}`$.
## B IR Subtractions
In this Appendix we list the IR subtraction, $`_{sub}`$ of equation (87), necessary to control numerical integration of IR divergent integrals. A gluon mass, $`\lambda /a_s`$, is introduced into $`D_{\mu \nu }`$ by replacing the first factor in (A24) by
$$\frac{1}{(\widehat{k}^2)^2}\frac{1}{\widehat{k}^2}\frac{1}{\widehat{k}^2+\lambda ^2}.$$
(B1)
The lattice wave function renormalization $`Z_2`$ must reproduce the same IR divergence structure as in continuum QCD. For $`Z_2^1`$ at one-loop the IR divergence is
$`{\displaystyle \frac{\alpha _s}{3\pi }}[1(\alpha _g1)]\mathrm{ln}(\lambda ^2)m=0`$ (B2)
$`{\displaystyle \frac{\alpha _s}{3\pi }}[\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}2}(\alpha _g1)]\mathrm{ln}(\lambda ^2)m>0`$ (B3)
We note that by the $`m=0`$ theory we mean one in which the quark mass has been set to zero before taking the limit $`\lambda 0`$. This is the usual practice in much of the literature on massless lattice perturbation theory. Alternatively one could take the limit $`am0`$ and $`\lambda 0`$ keeping $`am\lambda `$. Since we want to compare with some of the massless literature with improved glue actions (e.g. ) we adopt the first definition in this article. In our massive calculations we do not go to extremely small masses and have not attempted to isolate $`\mathrm{ln}(am)`$ contributions.
For our VEGAS integrations it was convenient to separate the $`d/dp_0`$ derivative in (87) into two parts
$$\frac{d}{dp_0}\left[V_\mu \frac{\overline{G}_0}{a_t}V_\nu \right]\frac{d}{dp_0}\left[\frac{VGV_{num}}{VGV_{den}}\right]=\frac{VGV_{num}^{}}{VGV_{den}}\frac{VGV_{num}}{(VGV_{den})^2}VGV_{den}^{}.$$
(B4)
Corresponding to the two parts with derivatives acting on the numerator or denominator, respectively, we introduce two separate subtraction terms $`_{sub}^{num}`$ and $`_{sub}^{den}`$. These are obtained by calculating the self energy diagram in continuum Euclidean perturbation theory with an appropriate mass $`m_{\mathrm{𝑒𝑓𝑓}}`$ and the mass-shell condition $`p=(im_{\mathrm{𝑒𝑓𝑓}},\stackrel{}{0})`$. The effective mass follows from expanding the lattice integrand in (87) about small $`k`$ and comparing with the continuum calculation . It is given in (91). After converting to the dimensionless integration variables $`k_\mu `$ of (87) one has for the part with the derivative acting on the denominator
$`_{sub}^{den}=\theta (\mathrm{\Lambda }^2k^2)`$ (B5)
$`\times \left\{{\displaystyle \frac{4\chi (\chi ^2k_0^2+b^2/4)((k^2)^2b^2\chi ^2k_0^2)}{(k^2+\lambda ^2)((k^2)^2+b^2\chi ^2k_0^2)^2}}+(\alpha _g1)\chi {\displaystyle \frac{\chi ^2k_0^2(b^2+2k^2)}{k^2(k^2+\lambda ^2)((k^2)^2+b^2\chi ^2k_0^2)}}\right\},`$ (B6)
with
$$k^2=\chi ^2k_0^2+\underset{j}{}k_j^2b=2a_sm_{\mathrm{𝑒𝑓𝑓}}.$$
(B8)
The $`\theta `$-function imposes a cutoff on $`_{sub}^{den}`$ so that it vanishes identically for $`k^2>\mathrm{\Lambda }^2`$, where $`\mathrm{\Lambda }`$ is some number $`0<\mathrm{\Lambda }\pi `$. The subtraction term can be integrated analytically to give
$`F^{den}(m_{\mathrm{𝑒𝑓𝑓}}>0,\mathrm{\Lambda },\lambda )={\displaystyle _k}_{sub}^{den}(k,m_{\mathrm{𝑒𝑓𝑓}}>0,\mathrm{\Lambda },\lambda )`$ (B9)
$`={\displaystyle \frac{\alpha _s}{3\pi }}\{[2\mathrm{ln}({\displaystyle \frac{\mathrm{\Lambda }^2}{\lambda ^2}})+\mathrm{\hspace{0.17em}2}\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }+\sqrt{b^2+\mathrm{\Lambda }^2}}{b}}\right)+{\displaystyle \frac{4\mathrm{\Lambda }^2}{b^4}}(b^2+3\mathrm{\Lambda }^2)+{\displaystyle \frac{\sqrt{b^2+\mathrm{\Lambda }^2}}{b^4}}2\mathrm{\Lambda }(b^26\mathrm{\Lambda }^2)]`$ (B10)
$`+(\alpha _g1)[\mathrm{ln}({\displaystyle \frac{\mathrm{\Lambda }^2}{\lambda ^2}})+{\displaystyle \frac{2\mathrm{\Lambda }^2}{b^4}}(\mathrm{\Lambda }^2+2b^2){\displaystyle \frac{\mathrm{\Lambda }(2\mathrm{\Lambda }^2+3b^2)}{b^4}}\sqrt{b^2+\mathrm{\Lambda }^2}\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }+\sqrt{b^2+\mathrm{\Lambda }^2}}{b}}\right)]\}`$ (B11)
and
$$F^{den}(m_{\mathrm{𝑒𝑓𝑓}}0,\mathrm{\Lambda },\lambda )=\frac{\alpha _s}{3\pi }\left[1+\frac{\alpha _g1}{2}\right]\mathrm{ln}(\frac{\mathrm{\Lambda }^2}{\lambda ^2}).$$
(B13)
The contribution in (B4) from the derivative acting on the numerator is
$$_{sub}^{num}=\theta (\mathrm{\Lambda }^2k^2)\left\{\frac{2\chi k^2}{((k^2)^2+b^2\chi ^2k_0^2)(k^2+\lambda ^2)}+(\alpha _g1)\chi \frac{k^22\chi ^2k_0^2}{((k^2)^2+b^2\chi ^2k_0^2)(k^2+\lambda ^2)}\right\},$$
(B14)
which leads to
$`F^{num}(m_{\mathrm{𝑒𝑓𝑓}}>0,\mathrm{\Lambda },\lambda )={\displaystyle _k}_{sub}^{num}(k,m_{\mathrm{𝑒𝑓𝑓}}>0,\mathrm{\Lambda },\lambda )`$ (B15)
$`=`$ $`{\displaystyle \frac{\alpha _s}{3\pi }}\{[4\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }+\sqrt{b^2+\mathrm{\Lambda }^2}}{b}}\right){\displaystyle \frac{4\mathrm{\Lambda }^2}{b^2}}+{\displaystyle \frac{4\mathrm{\Lambda }}{b^2}}\sqrt{b^2+\mathrm{\Lambda }^2}]`$ (B17)
$`+(\alpha _g1)[\mathrm{ln}\left({\displaystyle \frac{\mathrm{\Lambda }+\sqrt{b^2+\mathrm{\Lambda }^2}}{b}}\right){\displaystyle \frac{2\mathrm{\Lambda }^2}{b^4}}(\mathrm{\Lambda }^2+2b^2)+{\displaystyle \frac{\mathrm{\Lambda }(2\mathrm{\Lambda }^2+3b^2)}{b^4}}\sqrt{b^2+\mathrm{\Lambda }^2}]\}`$
and
$$F^{num}(m_{\mathrm{𝑒𝑓𝑓}}0,\mathrm{\Lambda },\lambda )=\frac{\alpha _s}{3\pi }\left[\mathrm{\hspace{0.17em}2}+\frac{\alpha _g1}{2}\right]\mathrm{ln}(\frac{\mathrm{\Lambda }^2}{\lambda ^2}).$$
(B18)
$`F^{den}+F^{num}`$ reproduces the IR divergent logarithms of (B2). |
warning/0001/physics0001038.html | ar5iv | text | # 1 Introduction
## 1 Introduction
Through the years, a number of authors have attempted to avoid the problems inherent in the point-particle model by focussing upon finite soliton-like structures. Fields interacting non-linearly provide the binding without invoking any phenomenological elements. Einstein and Rosen pointed out many years ago that particles should be contained within a field theory and not exist as independent entities. Rosen made considerable progress in implementing such a program in a gauge-invariant manner by minimally coupling a scalar field to the Maxwell field. However, the soliton solutions yielded negative masses. Later, neutral quantized particle states of positive mass were found and a more complicated model invoking up to three scalar fields coupled to the Maxwell field was shown to be capable of modeling the known massive leptons. However, the particles were spinless and the view then was that a subsequent quantization of the theory would induce spin.
In 1991, one of the present authors suggested an alternative route to elementary particle modelling, namely as solitons of Dirac–Maxwell theory. Since Dirac–Maxwell theory had been so successful in describing electron spin and magnetic moment, predicting the existence of the positron and refining the energy levels in interacting systems such as hydrogen, it seemed reasonable that this might successfully extend to a self-interacting soliton structure to model the elementary particles themselves. Spin would already exist in such a model via the spinor structure of the wave function. Shortly thereafter, such solitons were found and their properties studied. A few years later, Lisi independently discovered some of the results in. Recently, there has been a revival of interest in this field and in particular, the issue of gravitational coupling in the Dirac–Maxwell system has been considered. However, there was the misconception that gravitation was a necessary ingredient for the creation of the soliton.
In this paper, we develop the essential results in and and discuss the role of gravitation in soliton structure. The experimental inputs are the respective masses of the electron, muon and tau, their charge and as a constraint, the upper limit to their size which is $`10^{16}`$ cm. The plan of the paper is as follows: in sec. 2, we set out the essential coupled Dirac–Maxwell equations to be solved. The structure of the Dirac wave function in spherical coordinates is given and particularized to the case of electric field dominance. The equation is separated in sec. 3 and we contrast the standard treatment in which a potential function is imposed such as in the case of hydrogen and the present case of the soliton where the derivation of the potential is part of the problem. The formal structure of the potential in terms of the Green’s function is given. It is shown that there do exist spherically symmetric potentials for appropriate choices of quantum numbers.
In sec. 4, the spherically symmetric energy-momentum tensor is derived. The relationship between the parameters in the Dirac equation and the physically measured quantities is discussed and the expression for the spatial spread of the soliton is given. The various constraints including singularity avoidance lead to the required boundary conditions for the problem.
In sec. 5, the results are presented. New variables of convenience for numerical integration are introduced. The parameters leading to twenty ground state solitons are listed. It is found that there is a critical range which leads to solitons within the experimentally observed upper limit to the size of the electron. Excited states are presented and the mass ratios are found.
In the final section 6, the essential achievements as well as the limitations of the results are discussed. It is stressed that the solitons have been found without the requirement of significant gravitational interaction and it is conjectured that gravity will be significant for Dirac–Maxwell solitons when $`e/m1`$ in units for which $`G=c=1`$. In cgs units, this is $`2.58\times 10^4\text{esu}\text{gm}^1`$. By contrast, the $`e/m`$ ratio for the electron is $`2.04\times 10^{21}`$ or $`5.27\times 10^{17}\text{esu}\text{gm}^1`$ in cgs units.
## 2 Derivation of the Equations
The field equations are obtained from the Lagrangian of quantum electrodynamics
$$L=i\mathrm{}c\overline{\psi }\gamma ^\mu _\mu \psi mc^2\overline{\psi }\psi \frac{1}{16\pi }F^{\mu \nu }F_{\mu \nu }e\overline{\psi }\gamma ^\mu \psi A_\mu $$
(1)
where $`\psi =(\psi _1,\psi _2,\psi _3,\psi _4)^\text{T}`$ is the Dirac spinor, $`\overline{\psi }=\psi ^{}\gamma ^0=(\psi _1^{},\psi _2^{},\psi _3^{},\psi _4^{})`$, $`A^\mu =(\phi ,𝑨)`$ is the electromagnetic four-vector potential and $`F^{\mu \nu }=^\mu A^\nu ^\nu A^\mu `$ is the Maxwell tensor. The $`\gamma ^\mu `$ are $`4\times 4`$ Hermitian anticommuting matrices of the unit square
$$\gamma ^0=\left(\begin{array}{cc}I& 0\\ 0& I\end{array}\right),\gamma ^k=\left(\begin{array}{cc}0& \sigma _k\\ \sigma _k& 0\end{array}\right),k=1,2,3$$
where $`I`$ is the unit $`2\times 2`$ matrix and the $`\sigma _k`$ are the Pauli matrices
$$\sigma _1=\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right),\sigma _2=\left(\begin{array}{cc}0& i\\ i& 0\end{array}\right),\sigma _3=\left(\begin{array}{cc}1& 0\\ 0& 1\end{array}\right).$$
Variation with respect to $`A_\mu `$ and $`\overline{\psi }`$ respectively, yield the field equations
$$F_{,\nu }^{\mu \nu }=4\pi \overline{\psi }\gamma ^\mu \psi $$
(2)
$$i\mathrm{}c\gamma ^\mu _\mu \psi mc^2\psi e\gamma ^\mu \psi A_\mu =0.$$
(3)
If $`\psi `$ is chosen to be an energy eigenstate with energy $`E`$ and one chooses a static charge distribution with a four-vector potential of the form
$$A^\mu =(\varphi (r,\theta ,\phi ),𝑨^k(r,\theta ,\phi )),k=1,2,3$$
then the equations (2)-(3) are reduced to
$`\left[i\mathrm{}c𝜶\mathbf{}\mathbf{}+\alpha _4mc^2e𝜶\mathbf{}𝑨+e\varphi E\right]\psi =0`$ (4)
$`^2\varphi `$ $`=`$ $`4\pi e\psi ^{}\psi `$ (5)
$`\mathbf{}\times \left(\mathbf{}\times 𝑨\right)`$ $`=`$ $`4\pi e\psi ^{}𝜶\psi `$ (6)
where $`\alpha ^k=\gamma ^0\gamma ^k`$.
In spherical coordinates, $`(x,y,z)=(r\mathrm{sin}\theta \mathrm{cos}\phi ,r\mathrm{sin}\theta \mathrm{sin}\phi ,r\mathrm{cos}\theta )`$, the Dirac wave function has the structure
$$\psi (r,\theta ,\phi )_{\begin{array}{c}\\ [j=l+1/2]\end{array}}=\left(\begin{array}{c}\sqrt{\frac{lm}{2l+1}}gY_l^m\\ \sqrt{\frac{l+m+1}{2l+1}}gY_l^{m+1}\\ i\sqrt{\frac{l+m}{2l1}}fY_{l1}^m\\ i\sqrt{\frac{lm1}{2l1}}fY_{l1}^{m+1}\end{array}\right),\psi (r,\theta ,\phi )_{\begin{array}{c}\\ [j=l1/2]\end{array}}=\left(\begin{array}{c}\sqrt{\frac{l+m+1}{2l+1}}gY_l^m\\ \sqrt{\frac{lm}{2l+1}}gY_l^{m+1}\\ i\sqrt{\frac{lm+1}{2l+3}}fY_{l+1}^m\\ i\sqrt{\frac{l+m+2}{2l+3}}fY_{l+1}^{m+1}\end{array}\right)$$
(7)
where $`f=f(r)`$, $`g=g(r)`$ and the $`\{Y_l^m(\theta ,\phi )\}_{l,m}`$ is the set of orthonormal spherical harmonics defined for $`l=0,1,\mathrm{}`$, $`m=l,l+1,\mathrm{},l`$ and
$$Y_l^m(\theta ,\phi )=\sqrt{\frac{2l+1}{4\pi }\frac{(lm)!}{(l+m)!}}P_l^m(\mathrm{cos}\theta )e^{im\phi }.$$
(8)
In addition, $`m`$ is an integer such that $`jm+1/2j`$; $`(m+1/2)\mathrm{}`$ is the $`z`$-component of the total angular momentum.
Consider the spinor with $`j=1/2`$, $`l=0`$ and $`m=0`$ which implies from the above representation (7)
$`4\pi \psi ^{}\psi `$ $`=`$ $`f(r)^2+g(r)^2`$
$`4\pi \psi ^{}𝜶\psi `$ $`=`$ $`2f(r)g(r)\mathrm{sin}\theta (\mathrm{sin}\phi ,\mathrm{cos}\phi ,0)^\text{T}.`$
Resolving equations (5) and (6) into spherical coordinates gives
$`^2\varphi `$ $`=`$ $`e\left(f(r)^2+g(r)^2\right)`$
$`\mathbf{}\times \left(\mathbf{}\times 𝑨\right)|_{\widehat{r}}`$ $`=`$ $`0`$
$`\mathbf{}\times \left(\mathbf{}\times 𝑨\right)|_{\widehat{\theta }}`$ $`=`$ $`0`$
$`\mathbf{}\times \left(\mathbf{}\times 𝑨\right)|_{\widehat{\phi }}`$ $`=`$ $`2ef(r)g(r)\mathrm{sin}\theta .`$
Therefore, a four-vector potential of the form
$$A^\mu =(\varphi (r),A_\phi (r,\theta )\mathrm{sin}\phi ,A_\phi (r,\theta )\mathrm{cos}\phi ,0)$$
should be chosen where the components satisfy
$$\frac{d^2\varphi }{dr^2}+\frac{2}{r}\frac{d\varphi }{dr}=e\left(f(r)^2+g(r)^2\right)$$
(9)
$$\frac{^2A_\phi }{r^2}+\frac{2}{r}\frac{A_\phi }{r}+\frac{\mathrm{cot}\theta }{r^2}\frac{A_\phi }{\theta }+\frac{1}{r^2}\frac{^2A_\phi }{\theta ^2}\frac{A_\phi }{r^2\mathrm{sin}^2\theta }=2ef(r)g(r)\mathrm{sin}\theta .$$
(10)
Since the right hand side of equation (10) is nonzero, the theory can only be exact if $`A_\phi `$ is nonzero. However at this point we will impose the assumption of electric field dominance and hence the dominance of $`\varphi `$ over $`𝑨`$ or $`f(r)`$ dominance over $`g(r)`$.
For the validity of the approximation $`𝑨=\text{0}`$, one radial component of the spinor must dominate over the other so that
$$fgf^2+g^2.$$
It will be demonstrated that such objects do exist within the non-linear field. With this approximation the equations to solve reduce to a Dirac equation coupled to a Poisson equation:
$$\left[i\mathrm{}𝜶\mathbf{}\mathbf{}+\alpha _4mc^2+V(r)\right]\psi =E\psi $$
(11)
$$^2V=4\pi e^2\psi ^{}\psi .$$
(12)
With these facts in mind, we now turn to the separation of the stationary Dirac equation (11) with respect to a general central potential and the derivation of the form of $`\psi ^{}\psi `$ for a general set of quantum numbers.
## 3 Separation of the Equation
The separation procedure follows that given in Bethe and Salpeter. First one introduces quantum numbers $`l`$ and $`j`$; $`l`$ is the orbital angular momentum quantum number as well as being an integer $`0`$; $`j`$ is the total angular momentum quantum number and can assume just the two values $`l+1/2`$ and $`l1/2`$, (but only $`+1/2`$ for $`l=0`$). The forms assumed by the four components of $`\psi `$ are given explicitly in (7).
The explicit form of the Dirac equation (11) for the four components of the wave function is:
$`{\displaystyle \frac{\psi _3}{z}}+{\displaystyle \frac{\psi _4}{x}}i{\displaystyle \frac{\psi _4}{y}}{\displaystyle \frac{i}{\mathrm{}c}}\left[EV(r)mc^2\right]\psi _1`$ $`=`$ $`0`$ (13)
$`{\displaystyle \frac{\psi _4}{z}}{\displaystyle \frac{\psi _3}{x}}i{\displaystyle \frac{\psi _3}{y}}+{\displaystyle \frac{i}{\mathrm{}c}}\left[EV(r)mc^2\right]\psi _2`$ $`=`$ $`0`$ (14)
$`{\displaystyle \frac{\psi _1}{z}}+{\displaystyle \frac{\psi _2}{x}}i{\displaystyle \frac{\psi _2}{y}}{\displaystyle \frac{i}{\mathrm{}c}}\left[EV(r)+mc^2\right]\psi _3`$ $`=`$ $`0`$ (15)
$`{\displaystyle \frac{\psi _2}{z}}{\displaystyle \frac{\psi _1}{x}}i{\displaystyle \frac{\psi _1}{y}}+{\displaystyle \frac{i}{\mathrm{}c}}\left[EV(r)+mc^2\right]\psi _4`$ $`=`$ $`0.`$ (16)
Therefore, by inserting the assumed wave functions (7) into (13)-(16) and using identities similar to (A.6) we find that the following two coupled equations between $`f`$ and $`g`$ hold:
$`{\displaystyle \frac{1}{\mathrm{}c}}\left[EV(r)+mc^2\right]f(r)\left[{\displaystyle \frac{dg}{dr}}+{\displaystyle \frac{1+\kappa }{r}}g(r)\right]`$ $`=`$ $`0`$ (17)
$`{\displaystyle \frac{1}{\mathrm{}c}}\left[EV(r)mc^2\right]g(r)+\left[{\displaystyle \frac{df}{dr}}+{\displaystyle \frac{1\kappa }{r}}f(r)\right]`$ $`=`$ $`0`$ (18)
where the new quantum number $`\kappa `$ is defined as
$$\kappa =\{\begin{array}{ccc}l1\text{for}\hfill & j=l+1/2\hfill & (l=0,1,\mathrm{})\hfill \\ l\text{for}\hfill & j=l1/2\hfill & (l=1,2,\mathrm{}).\hfill \end{array}$$
(19)
These equations are valid for all spherically symmetric potentials $`V(𝒓)=V(r)`$ and together they replace expression (11).
At this point, the standard procedure is to specify an external spherically symmetric potential, an example of which is the electrostatic potential energy of the proton-electron interaction. That is, simply
$$V(r)=\frac{Ze^2}{r},$$
which is the fundamental solution of Laplace’s equation
$$^2V=4\pi Ze^2\delta ^3(𝒓)$$
(20)
where $`\delta ^3(𝒓)`$ is a three dimensional Dirac delta function centered at the origin. This is consistent with the far range<sup>1</sup><sup>1</sup>1By far range, we mean those distances much larger than the Bohr radius $`r\mathrm{}^2/me^2`$. behaviour that we expect to find for the self-field of the fermion since, when we compare (20) with (12) we see that the fermion is treated as an object without structure through the equality,
$$\psi ^{}\psi =\delta ^3(𝒓).$$
There is one additional problem that must be explored, namely how to couple relation (5) to (17)-(18). This will be achieved in three parts. First, we find the Green’s function for the equation (5). Second, we find an analytic form for the probability density $`\psi ^{}\psi `$ using (7). Once this equation is known, we can proceed to the third step which is to find $`V(r)`$ by forming the convolution of the Green’s function of step one, with the probability density of step two.
The potential $`V(r)`$ satisfies the Poisson equation (12) and by assuming that the solution is sufficiently regular, this can be converted to an integral equation
$$V(𝒓)=4\pi e^2G(𝒓,𝒓^{})\psi ^{}(𝒓^{})\psi (𝒓^{})𝑑𝒓^{}$$
(21)
where $`G(𝒓,𝒓^{})`$ is the Green’s function of the Laplacian operator in three dimensions
$$G(𝒓,𝒓^{})=\frac{1}{4\pi }\frac{1}{\left|𝒓𝒓^{}\right|}=\underset{l=0}{\overset{\mathrm{}}{}}\frac{1}{2l+1}\frac{r_<^l}{r_>^{l+1}}\underset{m=l}{\overset{l}{}}Y_l^m(\theta ,\phi )Y_l^m(\theta ^{},\phi ^{}).$$
(22)
With the Green’s function determined, we can turn our attention to the probability density. This is accomplished by using a pair of identities for the associated Legendre functions<sup>2</sup><sup>2</sup>2Both Eqs. (23)-(24) follow directly from Eqs. (8.5.1) and (8.5.3) of Abramowitz & Stegun.
$`(1\mu ^2)\left(P_l^{m+1}\right)^2`$ $`=`$ $`\left[(lm)\mu P_l^m(l+m)P_{l1}^m\right]^2,`$ (23)
$`(1\mu ^2)\left(P_{l1}^{m+1}\right)^2`$ $`=`$ $`\left[(l+m)\mu P_{l1}^m(lm)P_l^m\right]^2`$ (24)
together with the definition of the spherical harmonics (8). The resulting expression for the charge density of the Dirac particle is given by
$$\psi ^{}\psi =\frac{f^2+g^2}{2l+1}\left[(lm)\left|Y_l^{m+1}\right|^2+(l+m+1)\left|Y_l^m\right|^2\right]$$
(25)
when $`j=l+1/2`$ and
$$\psi ^{}\psi =\frac{f^2+g^2}{2l+1}\left[(l+m+1)\left|Y_l^{m+1}\right|^2+(lm)\left|Y_l^m\right|^2\right]$$
(26)
when $`j=l1/2`$.
Therefore by using (21), (22) and (25), one obtains the expression
$`V(𝒓)`$ $`=`$ $`4\pi e^2{\displaystyle G(𝒓,𝒓^{})\psi ^{}(𝒓^{})\psi (𝒓^{})𝑑𝒓^{}}`$
$`=`$ $`{\displaystyle \frac{4\pi e^2}{2l^{}+1}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}\frac{r_<^l}{r_>^{l+1}}\underset{m=l}{\overset{l}{}}\frac{1}{2l+1}Y_l^m(\theta ,\phi )Y_l^m(\theta ^{},\phi ^{})\left[f(r^{})^2+g(r^{})^2\right]}`$
$`\times \left[(l^{}m^{})\left|Y_l^{}^{m^{}+1}(\theta ^{},\phi ^{})\right|^2+(l^{}+m^{}+1)\left|Y_l^{}^m^{}(\theta ^{},\phi ^{})\right|^2\right]r^2dr^{}d(\mathrm{cos}\theta ^{})d\phi ^{}`$
for the case $`j=l+1/2`$. Similarly, with the use of (26), it can be shown that the potential $`V(𝒓)`$ takes the form
$`V(𝒓)`$ $`=`$ $`{\displaystyle \frac{4\pi e^2}{2l^{}+1}}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}\frac{r_<^l}{r_>^{l+1}}\underset{m=l}{\overset{l}{}}\frac{1}{2l+1}Y_l^m(\theta ,\phi )Y_l^m(\theta ^{},\phi ^{})\left[f(r^{})^2+g(r^{})^2\right]}`$
$`\times \left[(l^{}+m^{}+1)\left|Y_l^{}^{m^{}+1}(\theta ^{},\phi ^{})\right|^2+(l^{}m^{})\left|Y_l^{}^m^{}(\theta ^{},\phi ^{})\right|^2\right]r^2dr^{}d(\mathrm{cos}\theta ^{})d\phi ^{}`$
for the case $`j=l1/2`$. It is to be noted that the primed indices $`(l^{},m^{})`$ correspond to the angular momentum of the particle, while the unprimed indices run over the complete set of permissible angular momentum quantum numbers. By performing the angular integration of the above formulae, one can immediately conclude that both of the above integrals vanish except when $`m=0`$ and $`l=0,2,\mathrm{},2l^{}`$. This implies that
$`V(𝒓)`$ $`=`$ $`4\pi e^2{\displaystyle \frac{2(l^{}j^{})}{2l^{}+1}}{\displaystyle \underset{n=0}{\overset{l^{}}{}}}{\displaystyle \frac{Y_{2n}^0(\theta ,\phi )}{4n+1}}{\displaystyle _{r^{}=0}^{\mathrm{}}}{\displaystyle \frac{r_<^{2n}}{r_>^{2n+1}}}\left[f(r^{})^2+g(r^{})^2\right]r^2𝑑r^{}`$ (27)
$`\times \left[(\kappa ^{}+m^{}+1)l^{},m^{}+1|Y_{2n}^0|l^{},m^{}+1+(\kappa ^{}m^{})l^{},m^{}|Y_{2n}^0|l^{},m^{}\right]`$
where the cases $`j=l\pm 1/2`$ have been combined by the application of the definition of $`\kappa ^{}`$. Expression (27) replaces the equation (12). When written in this form, it is clearly seen that the potential $`V(𝒓)`$ is not in general spherically symmetric. Table 1 lists the potential (27) for $`l^{}=0,1`$ and illustrates the fact that there exists spherically symmetric states with $`l^{}0`$.
A localized solution of this model must satisfy the field equations (17) and (18) for $`f`$ and $`g`$ and a given energy $`E`$ where the potential is given by the expression (27). Moreover, it is required that the total probability
$$\psi |\psi =\underset{i=1}{\overset{4}{}}\psi _i|\psi _i=_0^{\mathrm{}}\left(f^2+g^2\right)r^2𝑑r<\mathrm{}.$$
Since the equations which describe the spatial evolution of the wave function (17)-(18) were derived under the assumption that the potential, $`V(r)`$, is spherically symmetric, they are not valid for an extended Dirac particle in an arbitrary state of angular momentum. We have shown that there do exist certain choices of $`l`$ and $`m`$ where the probability density is spherically symmetric and it is these cases in which our primary interest lies.
We can conclude that with the spinor representation given by (7), there are essentially three differential equations to be solved simultaneously. Equations (17)-(18) specify the spatial evolution of the wave function and equation (27) reflects the spatial extent of the self-field of the particle. A strategy for solving these intrinsically non-linear equations, as well as a few of their interesting properties, will be explored in the following sections.
## 4 Boundary Conditions
From the previous section we have found that the equations to be satisfied for a self-interacting fermion are equations (17)-(18) and
$$^2V=4\pi e^2\frac{2(lj)}{2l+1}\left(f^2+g^2\right)\left[(\kappa +m+1)\left|Y_l^{m+1}(\theta ,\phi )\right|^2+(\kappa m)\left|Y_l^m(\theta ,\phi )\right|^2\right]$$
(28)
where we have combined the $`j=l\pm 1/2`$ cases by using the definition of $`\kappa `$. Since we have assumed that the potential $`V`$ in equations (17)-(18) is spherically symmetric, this necessarily restricts the choice of $`l`$ and $`m`$. Assume from this point on that $`l`$ and $`m`$ are chosen to satisfy this criterion. Consequently, equation (28) becomes
$$^2V=e^2\left(f^2+g^2\right).$$
(29)
Since the soliton asires as a coupling between Dirac and Maxwell fields, the energy $`E`$ that appears in the Dirac equation is not the total energy of the particle. The total energy can be obtained by calculating the $`T_0^0`$ component of the energy-momentum tensor. For our field, the Lagrangian is given by equation (1) where $`A^\mu `$ is the vector potential of the electromagnetic field. One generates the symmetric energy-momentum tensor directly from the Lagrangian in the form
$$T^{\mu \nu }=\frac{L}{g_{\mu \nu }}\frac{g^{\mu \nu }}{2}L.$$
(30)
Applying (30) to (1) yields
$`T^{\mu \nu }`$ $`=`$ $`\left[{\displaystyle \frac{i\mathrm{}c}{2}}\left(\overline{\psi }\gamma ^\mu ^\nu \psi +\overline{\psi }\gamma ^\nu _\mu \psi \right){\displaystyle \frac{1}{4\pi }}F^{\alpha \mu }F^{\beta \nu }g_{\alpha \beta }{\displaystyle \frac{e}{2}}\left(\overline{\psi }\gamma ^\mu \psi A^\nu +\overline{\psi }\gamma ^\nu \psi A^\mu \right)\right]`$
$`{\displaystyle \frac{g^{\mu \nu }}{2}}\left[i\mathrm{}c\overline{\psi }\gamma ^\alpha g_{\alpha \beta }^\beta \psi mc^2\overline{\psi }\psi {\displaystyle \frac{1}{8\pi }}F^{\alpha \beta }F_{\alpha \beta }e\overline{\psi }\gamma ^\alpha \psi A^\beta g_{\alpha \beta }\right].`$
Further simplification gives
$$T_0^0=E\psi ^{}\psi +\frac{1}{8\pi }\left(\frac{d\varphi }{dr}\right)^2.$$
This yields an expression for the total energy, $`E_{\text{tot}}`$, of
$`E_{\text{tot}}`$ $`=`$ $`{\displaystyle T_0^0𝑑V_{ol}}`$ (31)
$`=`$ $`E{\displaystyle _0^{\mathrm{}}}\left(f^2+g^2\right)r^2𝑑r+{\displaystyle \frac{1}{2}}{\displaystyle \left(\frac{d\varphi }{dr}\right)^2r^2𝑑r}`$
where $`dV_{ol}`$ is an infinitesimal volume element. This total energy should be associated with the observed mass of the particle as $`E_{\text{tot}}=mc^2`$. There is still sufficient freedom remaining to set $`lim_r\mathrm{}V(r)=0`$ because the spinor is invariant under the transformation
$$VV+\beta ;EE+\beta $$
for any real-valued $`\beta `$.
The mass $`m`$ and the charge $`e`$ that appear in the Dirac equation are not necessarily the experimentally measured quantities just as the charge that appears at a vertex of a Feynman graph is not the experimentally measured charge of the particle. Because of this, we will replace the $`m`$ in (17)-(18) by the symbol $`\mu `$. In addition, the $`e`$ in (29) will be replaced by an $`ϵ`$. The symbols $`m`$ and $`e`$ will be reserved for the physically observed quantities. With these substitutions, we convert to a set of variables whereby equations (17), (18) and (31) are independent of any physical constants. The particular transformation chosen is
$$f=\eta F,g=\eta G,r=\frac{\mathrm{}x}{\mu c},E=\lambda \mu c^2,V=\mu c^2U$$
where $`\eta ^2=\mu ^3c^4/ϵ^2\mathrm{}^2`$. These redefined variables have the following dimensions in terms of length ($`L`$):
$$[x]=L^0;[\lambda ]=L^0;[F]=L^{3/2};[G]=L^{3/2};[U]=L^0.$$
This yields the transformed equations:
$`\left[\lambda U(x)+1\right]F(x)\left[{\displaystyle \frac{dG}{dx}}+{\displaystyle \frac{1+\kappa }{x}}G(x)\right]`$ $`=`$ $`0`$ (32)
$`\left[\lambda U(x)1\right]G(x)+\left[{\displaystyle \frac{dF}{dx}}+{\displaystyle \frac{1\kappa }{x}}F(x)\right]`$ $`=`$ $`0`$ (33)
$`^2U+\left(F^2+G^2\right)`$ $`=`$ $`0`$ (34)
where $`^2`$ is now the Laplacian with respect to the $`x`$ coordinate. The mass of the soliton comes from the transformed version of the total energy expression (31),
$$mc^2=\frac{\mathrm{}\mu c^3}{ϵ^2}\left[\lambda _0^{\mathrm{}}\left(F^2+G^2\right)x^2𝑑x+\frac{1}{2}_0^{\mathrm{}}\left(\frac{dU}{dx}\right)^2x^2𝑑x\right]$$
(35)
and the total charge is given as the integral of the charge density
$$e=ϵ\rho 𝑑V_{ol}=\frac{\mathrm{}c}{ϵ}_0^{\mathrm{}}\left(F^2+G^2\right)x^2𝑑x.$$
(36)
We will show that if the charge $`ϵ`$ is replaced by $`e`$, that the $`f`$ component of the spinor greatly dominates the $`g`$ component. By substituting $`ϵ=e`$ and choosing a value for $`m`$, the value of $`\mu `$ can be determined numerically once the spatial extent of the soliton is known. In this case, the expectation value of the radius of the particle becomes
$$r=\frac{\mathrm{}c}{\mu c^2}x=\frac{e^2}{\mu c^2}\frac{\mathrm{}c}{e^2}\frac{{\displaystyle ^2Ux^3dx}}{{\displaystyle ^2Ux^2dx}}=\frac{e^2}{\mu c^2}\frac{{\displaystyle ^2Ux^3dx}}{\left[{\displaystyle ^2Ux^2dx}\right]^2}.$$
To stay within the current experimental bounds of the mean charge radius, this value must be less than $`r_{exp}`$ which is $`10^{18}m`$ in the case of an electron. Hence,
$$_0^{\mathrm{}}^2Ux^3dx\frac{r_{exp}}{r_e}\frac{\mu }{m_e}\left[_0^{\mathrm{}}^2Ux^2dx\right]^2$$
where $`r_e`$ is the classical electron radius $`r_e=e^2/m_ec^2`$.
Since we know that $`U`$ has zero slope at $`x=0`$ and that it must behave as $`N/x`$ for large argument ($`N`$ is the amount of enclosed charge), we assume as a first approximation, that $`U`$ can be represented as the electrostatic potential produced by a sphere of radius $`R_0`$ with uniform charge density. Therefore,
$$U(r)=\{\begin{array}{cc}\frac{N}{R_0}\left[\frac{3}{2}\frac{1}{2}\frac{r^2}{R_0^2}\right]\text{for}\hfill & r<R_0\hfill \\ \frac{N}{r}\text{for}\hfill & rR_0.\hfill \end{array}$$
(37)
With this representation, one finds that
$$r=\frac{9}{4}\frac{r_em_eR_0}{\mu N}$$
which means that since $`0<r<r_{exp}`$, we can conclude that
$$0<\frac{R_0}{\mu N}<\frac{4}{9}\frac{r_{exp}}{r_em_e}3.088\times 10^4c^2/\text{MeV}$$
in the case of the electron.
Let $`R_0`$ be defined as the effective range of the non-Coulombic behaviour of the potential energy so that for $`x>R_0`$, $`U(x)N/x`$. Since $`U`$ is a solution to a Poisson equation with a negative definite charge density, $`U(0)`$ must be larger than $`U(R_0)`$. This can be quickly verified by considering the opposite. If $`U(0)<U(R_0)`$ then there exists some $`r(0,R_0)`$ such that $`U^{}(r)=0`$. Therefore integrating (34) from $`0`$ to $`r`$, one obtains
$$r^2U^{}(r)=0=_0^r\left(F^2+G^2\right)x^2𝑑x$$
which is clearly a contradiction.
To determine the initial values of $`F`$ and $`G`$, one simply eliminates either $`F`$ or $`G`$ from (32-33), say $`F`$, which leads to a second order equation for the other, namely,
$$G^{\prime \prime }+PG^{}+QG=0$$
where both $`P`$ and $`Q`$ are functions of $`U`$, $`U^{}`$, $`\kappa `$, $`\lambda `$ and $`x`$. To avoid a singularity in the potential $`U(x)`$, it must be both bounded and have zero slope in a neighbourhood of the origin. Moreover, both $`F`$ and $`G`$ are bounded in this same neighbourhood. From this, it is easy to verify that
$$F(0)=\{\begin{array}{cc}0\hfill & \kappa =1\hfill \\ \text{arbitrary}\hfill & \kappa =+1\hfill \\ 0\hfill & \text{other}\kappa ,\hfill \end{array}$$
(38)
$$G(0)=\{\begin{array}{cc}\text{arbitrary}\hfill & \kappa =1\hfill \\ 0\hfill & \kappa =+1\hfill \\ 0\hfill & \text{other}\kappa .\hfill \end{array}$$
(39)
Furthermore, by examining the indicial equation, it can be shown that no fractional powers exist in a power series solution of either $`F`$ or $`G`$ about the origin $`x=0`$.
Summarizing the boundary conditions:
$$U(x)<U(0)<\mathrm{},x[0,\mathrm{}),$$
$$U^{}(0)=0,$$
together with the conditions (38), (39). For the case $`\kappa =1`$, the initial values of $`U`$, $`G`$ and the energy $`\lambda `$ are determined by the requirement that the wave function $`\psi `$, and hence both $`F`$ and $`G`$, vanish exponentially as $`x\mathrm{}`$.
## 5 Results
In the search for numerical solutions it was specified that $`\kappa =1`$ and $`\lambda =1`$ giving the set of differential equations
$`{\displaystyle \frac{dG}{dx}}`$ $`=`$ $`\left[2U(x)\right]F(x)`$
$`{\displaystyle \frac{dF}{dx}}`$ $`=`$ $`{\displaystyle \frac{2}{x}}F(x)+U(x)G(x)`$
$`^2U`$ $`=`$ $`F^2G^2.`$
To find a soliton, the values of $`F(0)`$, $`G(0)`$ are specified and a search is made for the value of $`U(0)`$ whereby $`lim_x\mathrm{}xF(x)=0`$ and $`lim_x\mathrm{}xG(x)=0`$. Only values of $`G(0)>0`$ are considered because the equations are symmetric under the transformation $`GG`$, $`FF`$, $`UU`$. In a neighbourhood of a ground state soliton, the radial probability density is numerically seen to have a single well-defined minimum value for $`x>0`$. The choice of $`\kappa =1`$ gives the initial condition $`F(0)=0`$.
The choice of $`\lambda =1`$ is simply a numerical convenience. Outside the neighbourhood of a soliton it is expected that the potential will behave as $`U(x)A+B/x`$ for large $`x`$. The value of $`\lambda `$ should have been chosen so that the asymptotic behaviour of the potential $`U(x)`$ is purely Coulombic in nature. By defining a shifted potential $`\stackrel{~}{U}(x)=U(x)lim_x\mathrm{}U(x)`$, this value of $`\lambda `$ must satisfy $`1U(x)=\lambda \stackrel{~}{U}(x)`$. Therefore after a soliton is found the value of $`\lambda `$ is given as $`\lambda =1lim_x\mathrm{}U(x)`$. In addition, the starting value of $`\stackrel{~}{U}(x)`$ is given by $`\stackrel{~}{U}(0)=U(0)+\lambda 1`$.
Using the redefined value of $`\lambda `$ the observed charge and mass of the particle are compared to the values used in the Lagrangian by using the expressions (36) and (35) respectively. By defining
$$\begin{array}{c}𝒫=_0^{\mathrm{}}\left(F^2+G^2\right)x^2𝑑x,𝒳=_0^{\mathrm{}}\left(F^2+G^2\right)x^3𝑑x,\\ =\lambda 𝒫+\frac{1}{2}_0^{\mathrm{}}\left(\frac{dU}{dx}\right)^2x^2𝑑x,\end{array}$$
the charge ratio $`ϵ/e`$ is given as
$$\frac{ϵ}{e}=\frac{\mathrm{}c}{e^2}𝒫=\frac{𝒫}{\alpha }$$
where $`\alpha `$ is the fine structure constant. The mass ratio $`\mu /m=𝒫^2/\alpha `$ and the expectation value for the radius of the soliton is
$$r=\frac{(f^2+g^2)r^3𝑑r}{(f^2+g^2)r^2𝑑r}=\frac{\mathrm{}}{\mu c}\frac{(F^2+G^2)x^3𝑑x}{(F^2+G^2)x^2𝑑x}=r_e\left(\frac{m_e}{m}\right)\frac{𝒳}{𝒫^3}.$$
Both of the quantities $`𝒫`$ and $`𝒳`$ are positive. However, depending upon the value of $`\lambda `$, $``$ could be positive, negative or even zero if the electromagnetic and “bare mass” terms in the energy exactly cancel. A negative value for $``$ will give an unphysical negative value for the observed radius $`r`$. Because of this ambiguity, both the value of $`r`$ and the particle width $`\mathrm{\Delta }r=\sqrt{r^2r^2}`$ are presented. Tables 2 and 3 respectively list the numerical parameters and the observed properties of a number of ground state particles found where $`m`$ was taken to be the observed mass of the electron $`m_e`$.
Figure 1 illustrates the radial behaviour of $`F`$ and $`G`$ for the case $`ϵ=e(i=1)`$. It is to be noted that for $`x>0`$, $`F`$ is much larger than $`G`$ and as a consequence, $`FGF^2+G^2`$. In fact, $`G`$ is so small that it resembles a straight line along the $`x`$ axis. This supports the argument that the four-vector potential can be reasonably approximated with only a radial $`A^0`$ component.
The characteristics of a typical soliton with $`ϵe`$ is illustrated with the choice $`ϵ/e=454.8(i=19)`$. In this case the potential plays a much more dominant role in holding the particle together than in the case $`ϵ=e`$. However, since in this case the approximation of $`FGF^2+G^2`$ is violated, one would have to solve the full model (equations (9)-(10)) to properly analyse this situation. This would be a far more complicated problem. Figure 2 illustrates the radial components of this spinor and it shows that the magnitude of $`G`$ is now comparable to the magnitude of $`F`$. Table 3 also shows that the choice of $`ϵ=389.0e`$, $`\mu =2.360\times 10^{12}m_e`$ ($`i=15`$) yields a soliton with an expectation value for the radius of $`5.05\times 10^{23}m`$. This size is well within the present experimentally determined upper limit for the electron radius of $`10^{18}m`$.
These equations also exhibit excited states. The $`n^{\text{th}}`$ excited state of our field is characterized through the functions $`F_n(x)`$, $`G_n(x)`$ and $`U_n(x)`$ for which the $`G_n`$ component crosses the abscissa $`n+1`$ times while the $`F_n`$ component crosses it $`n`$ times. Once the ground state solution is found, the value of $`\mu `$ can be determined through equation (36). The corresponding $`n^{\text{th}}`$ excited state is that excited state with the same observed charge ratio, $`ϵ/e`$, as the ground state. Therefore, in this interpretation of the theory, the ratio of the mass of the $`n^{\text{th}}`$ excited state to the ground state is given by the expression
$$\frac{m_n}{m_0}=\frac{\mu /m_0}{\mu /m_n}=\frac{\lambda _n{\displaystyle _0^{\mathrm{}}}\left(F_n^2+G_n^2\right)x^2𝑑x+{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{dU_n}{dx}}\right)^2x^2𝑑x}{\lambda {\displaystyle _0^{\mathrm{}}}\left(F^2+G^2\right)x^2𝑑x+{\displaystyle \frac{1}{2}}{\displaystyle _0^{\mathrm{}}}\left({\displaystyle \frac{dU}{dx}}\right)^2x^2𝑑x}.$$
Figure 3 shows radial probability density of the first three states for the case $`G(0)=1`$. Each of these solitons has a different value of $`ϵ/e`$.
Figure 4 illustrates the behaviour of the mass ratio, $`\mu /m`$, as a function of the charge ratio $`ϵ/e`$ for the ground state and the first two excited states. For each class of particles there is a charge ratio where the electromagnetic and bare mass components of the energy balance making $`=0`$. At this value of $`ϵ/e`$, the mass ratio $`\mu /m\mathrm{}`$. At charge ratios less than this critical value the mass ratio is negative whereas charge ratios above this critical value result in a positive value of $`\mu /m`$. There is numerical evidence that each class of particles has an upper bound for the charge ratio. Above this maximum charge ratio we were unable to find any solutions such that $`lim_x\mathrm{}xF(x)0`$ or $`lim_x\mathrm{}xG(x)0`$. This necessarily restricts the definition of the mass ratio defined above. Figure 4 also illustrates the fact that at moderate charge ratios, the electromagnetic field does not contain an appreciable amount of the particle energy resulting in the behaviour $`|\mu /m|ϵ/e`$.
The mass ratios of the first and second excited states with respect to the ground state solutions are shown in figure 5. This ratio is only defined up to a maximum value of $`ϵ/e`$ since beyond $`ϵ/e550`$, a ground state fails to exist. For excited states, this maximum admissible charge ratio increases. This implies that for a fixed value of $`ϵ/e`$ there may not exist a ground state solution, but there will be arbitrarily many excited states. As is readily apparent from figure 5, the only appreciable mass splitting occurs for large charge ratios. However, it is precisely for large charge ratios where our approximation that $`FGF^2+G^2`$ breaks down.
## 6 Concluding Remarks
We have seen that spherically symmetric Dirac-Maxwell solitons can be constructed and with a charge and mass to model the electron successfully. However, it should be noted that the higher energy excited states of this form did not yield the large mass separations of the muon and tau relative to the electron in this model. The search thus far has been restricted to spherical solitons. It is conceivable that a relaxation of this restriction or some other change in conditions would increase the mass splitting. In any event, we have shown that Dirac-Maxwell solitons exist and are capable of modelling an electron where the charge-to-mass ratio is the observed $`10^{21}`$ in units in which $`G=c=1`$. Furthermore, we have found a charge-to-mass ratio that simultaneously yields the observed charge and mass of the electron as well as exhibiting a degreee of compactification that is well within the current experimental upper limit. Finster et al. have considered Einstein–Dirac–Maxwell (EDM) solitons and concluded that it is the interaction with gravitation which is responsible for the existence of bound states. However, we see here that bound states exist with negligible gravitational interaction. While the $`e/m`$ ratio at which significant gravitational coupling sets in is yet to be determined for EDM solitons, it is our conjecture that this will be so at the same level that was found earlier in the case of minimally coupled scalar interaction, namely for $`e/m1`$. The known fundamental charged particles of nature, on the other hand have enormous $`e/m`$ ratios.
REFERENCES
## Appendix A Derivatives of $`f(r)Y_l^m(\theta ,\phi )`$
In the Dirac wave equation, all of the derivatives are with respect to Cartesian coordinates. We can change to a spherical polar representation via the transformation
$$\begin{array}{c}x=r\mathrm{sin}\theta \mathrm{cos}\phi \hfill \\ y=r\mathrm{sin}\theta \mathrm{sin}\phi \hfill \\ z=r\mathrm{cos}\theta .\hfill \end{array}$$
By applying the chain rule, it is trivial to show that this changes the first order partial derivatives via
$`{\displaystyle \frac{}{x}}`$ $`=`$ $`\mathrm{sin}\theta \mathrm{cos}\phi {\displaystyle \frac{}{r}}+{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\phi }{r}}{\displaystyle \frac{}{\theta }}{\displaystyle \frac{\mathrm{sin}\phi }{r\mathrm{sin}\theta }}{\displaystyle \frac{}{\phi }}`$ (A.1)
$`{\displaystyle \frac{}{y}}`$ $`=`$ $`\mathrm{sin}\theta \mathrm{sin}\phi {\displaystyle \frac{}{r}}+{\displaystyle \frac{\mathrm{cos}\theta \mathrm{sin}\phi }{r}}{\displaystyle \frac{}{\theta }}+{\displaystyle \frac{\mathrm{cos}\phi }{r\mathrm{sin}\theta }}{\displaystyle \frac{}{\phi }}`$ (A.2)
$`{\displaystyle \frac{}{z}}`$ $`=`$ $`\mathrm{cos}\theta {\displaystyle \frac{}{r}}{\displaystyle \frac{\mathrm{sin}\theta }{r}}{\displaystyle \frac{}{\theta }}.`$ (A.3)
If the functions $`\psi _j(j=1,\mathrm{},4)`$ from expression (7) are substituted into (13)-(16), and if one uses the formulas given in Bethe and Salpeter for the derivatives of a function of the form $`f(r)Y_l^m(\theta ,\phi )`$ with respect to $`x`$, $`y`$, and $`z`$, one finds a coupled pair of first order ordinary equations for $`f(r)`$ and $`g(r)`$.
For example, in order to calculate
$$\frac{}{z}\left[f(r)Y_l^m(\theta ,\phi )\right],$$
we first require the identities
$$\mathrm{cos}\theta P_l^m(\mathrm{cos}\theta )=\frac{1}{2l+1}\left[(lm+1)P_{l+1}^m(\mathrm{cos}\theta )+(l+m)P_{l1}^m(\mathrm{cos}\theta )\right]$$
(A.4)
$$\mathrm{sin}\theta \frac{d}{d\theta }P_l^m(\mathrm{cos}\theta )=\frac{1}{2l+1}\left[l(lm+1)P_{l+1}^m(\mathrm{cos}\theta )(l+1)(l+m)P_{l1}^m(\mathrm{cos}\theta )\right],$$
(A.5)
which can both be verified through the use of Rodrigues’ formula
$$P_l^m(\mu )=\frac{(1)^m}{2^ll!}(1\mu ^2)^{m/2}\frac{d^{l+m}}{d\mu ^{l+m}}(\mu ^21)^l.$$
Writing $`Y_l^m(\theta ,\phi )`$ as a function of $`P_l^m`$ by using (8) gives the relationship
$$\frac{}{z}\left[f(r)Y_l^m(\theta ,\phi )\right]=\sqrt{\frac{2l+1}{4\pi }\frac{(lm)!}{(l+m)!}}e^{im\phi }\left[\mathrm{cos}\theta P_l^m(\mathrm{cos}\theta )\frac{df}{dr}\mathrm{sin}\theta \frac{d}{d\theta }P_l^m(\mathrm{cos}\theta )\frac{f}{r}\right].$$
By substituting (A.4-A.5) in the above, collecting terms, and applying the definition of $`Y_l^m(\theta ,\phi )`$ once again, one obtains the simplification
$`{\displaystyle \frac{}{z}}\left[f(r)Y_l^m(\theta ,\phi )\right]`$ $`=`$ $`\sqrt{{\displaystyle \frac{(lm+1)(l+m+1)}{(2l+1)(2l+3)}}}Y_{l+1}^m(\theta ,\phi )\left[{\displaystyle \frac{df}{r}}{\displaystyle \frac{l}{r}}f\right]`$ (A.6)
$`+\sqrt{{\displaystyle \frac{(lm)(l+m)}{(2l1)(2l+1)}}}Y_{l1}^m(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}+{\displaystyle \frac{l+1}{r}}f\right].`$
Similar relationships for $`\frac{}{x}\pm i\frac{}{y}`$ can be found in Bethe and Salpeter<sup>3</sup><sup>3</sup>3See formula (A.38) and (A.39) respectively in Bethe and Salpeter., but there is a very elegant way to derive these operators by applying the Wigner–Eckart theorem.
First, we evaluate the matrix element $`l0|_0|l0`$ of the gradient operator, which is an example of a vector operator. Specifically,
$$_0=\frac{}{z},_\pm =\frac{1}{\sqrt{2}}\left(\frac{}{x}\pm i\frac{}{y}\right).$$
Since
$$_0f(r)Y_l^0=\frac{l+1}{\sqrt{(2l+1)(2l+3)}}Y_{l+1}^0\left[\frac{df}{dr}\frac{l}{r}f\right]+\frac{l}{\sqrt{(2l1)(2l+1)}}Y_{l1}^0\left[\frac{df}{dr}+\frac{l+1}{r}f\right]$$
for the special case of (A.6) where $`m=0`$, we have
$$l^{}0|_0|l0=\frac{l+1}{\sqrt{(2l+1)(2l+3)}}[\frac{df}{dr}\frac{l}{r}f]\delta _{l+1}^l^{}+\frac{l}{\sqrt{(2l1)(2l+1)}}[\frac{df}{dr}+\frac{l+1}{r}f]\delta _{l1}^l^{}.$$
Now, we are at a point where we can use the Wigner–Eckart theorem. By inspection, the general matrix element is given by
$`l^{}m^{}|_\mu |lm`$ $`=`$ $`(1)^{l^{}m^{}}\left(\begin{array}{ccc}l^{}& 1& l\\ m^{}& \mu & m\end{array}\right)l^{}l`$
$`=`$ $`(1)^m^{}{\displaystyle \frac{\left(\begin{array}{ccc}l^{}& 1& l\\ m^{}& \mu & m\end{array}\right)}{\left(\begin{array}{ccc}l^{}& 1& l\\ 0& 0& 0\end{array}\right)}}l^{}0|_0|l0.`$
After evaluating the $`3j`$ symbols, one can quickly verify the following equations.
$`{\displaystyle \frac{}{z}}[(f(r)Y_l^m(\theta ,\phi )]`$ $`=`$ $`\sqrt{{\displaystyle \frac{(lm+1)(l+m+1)}{(2l+1)(2l+3)}}}Y_{l+1}^m(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}{\displaystyle \frac{l}{r}}f\right]`$ (A.9)
$`+`$ $`\sqrt{{\displaystyle \frac{(lm)(l+m)}{(2l1)(2l+1)}}}Y_{l1}^m(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}+{\displaystyle \frac{l+1}{r}}f\right]`$
$`\left[{\displaystyle \frac{}{x}}+i{\displaystyle \frac{}{y}}\right]\left[f(r)Y_l^m(\theta ,\phi )\right]`$ $`=`$ $`\sqrt{{\displaystyle \frac{(l+m+1)(l+m+2)}{(2l+1)(2l+3)}}}Y_{l+1}^{m+1}(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}{\displaystyle \frac{l}{r}}f\right]`$ (A.10)
$``$ $`\sqrt{{\displaystyle \frac{(lm1)(lm)}{(2l1)(2l+1)}}}Y_{l1}^{m+1}(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}+{\displaystyle \frac{l+1}{r}}f\right]`$
$`\left[{\displaystyle \frac{}{x}}i{\displaystyle \frac{}{y}}\right]\left[f(r)Y_l^m(\theta ,\phi )\right]`$ $`=`$ $`\sqrt{{\displaystyle \frac{(lm+1)(lm+2)}{(2l+1)(2l+3)}}}Y_{l+1}^{m1}(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}{\displaystyle \frac{l}{r}}f\right]`$ (A.11)
$`+`$ $`\sqrt{{\displaystyle \frac{(l+m1)(l+m)}{(2l1)(2l+1)}}}Y_{l1}^{m1}(\theta ,\phi )\left[{\displaystyle \frac{df}{dr}}+{\displaystyle \frac{l+1}{r}}f\right].`$
Linear combinations of (A.10) and (A.11) yield the derivatives with respect to $`x`$ and $`y`$. |
warning/0001/cond-mat0001281.html | ar5iv | text | # Center-of-Mass Properties of the Exciton in Quantum Wells
## I Introduction
Excitons dominate the optical properties of low-dimensional semiconductor heterostructures such as quantum wells (QW) and quantum wires. The relative motion of the constituent particles and their center-of-mass (COM) motion determine different characteristics of the optical spectra and exciton kinetics.
The exciton relative motion in QW is well studied and understood. The confinement of the carriers along one or two spatial directions into regions comparable to or smaller than the bulk exciton size enhances the effect of the electron-hole Coulomb interaction. This results in larger binding energies and oscillator strengths and in an increased stability compared to bulk excitons. Therefore, excitons are observed even at room temperature in these structures. The effect of the reduced dimensionality is as a rule much larger on the exciton groundstate than on the excited states.
Details of the excitonic optical spectra of QW related to the COM motion like, e.g., inhomogeneous broadening and Stokes shift between photoluminescence (PL) and absorption are frequently used for structure characterization. These features are influenced by exciton localization and diffusion in the presence of interface or alloy disorder . Optical spectra and their temporal evolution are determined by the exciton formation processes and the subsequent energy and spin relaxation dynamics. Spatially resolved spectroscopy techniques like micro PL and near-field scanning optical microscopy allow direct observation of exciton COM quantization in local potential minima . All these phenomena are intimately related to the exciton COM properties whereby different energy and COM momentum regions of the exciton dispersion are probed in different processes.
In many heterostructure systems of interest like, e.g., GaAs/AlAs, InGaAs/InP and ZnCdSe/ZnSeS, the exciton can be described in the effective-mass approximation (Wannier exciton) due to its small bulk binding energy (e.g., 4 meV for GaAs). In this approximation, the degeneracy of the valence bands at the center of the Brillouin zone for materials of cubic or zinc-blende symmetry was first taken into account by Dresselhaus . He also pointed out the absence of a well-defined COM transformation due to this degeneracy. Altarelli and Lipari calculated the exciton COM dispersion for direct- and indirect-gap bulk semiconductors. They demonstrated that the ambiguity in the choice of the COM transformation can be used to achieve formal simplicity or optimal numerical convergence. For bulk GaAs, where the heavy- to light-hole mass ratio is large, the exciton dispersions are found to be strongly anharmonic and show avoided crossings between different branches.
In semiconductor QW, the broken translational symmetry in the growth direction leads to the splitting of heavy and light holes at the $`\mathrm{\Gamma }`$-point, and subsequent formation of heavy- and light-hole excitons. Due to the large hole-to-electron mass ratio, the influence of the valence-band mixing on the COM motion is greater than on the relative motion. The exciton dispersions are, thus, strongly non-parabolic. Direct consequences of the exciton dispersion anharmonicity in QW like slow indirect excitonic transitions due to camel-back shaped dispersions have been experimentally observed .
The multiband exciton (i.e., with the full coupling of heavy- and light-hole bands taken into account) has been theoretically and numerically thoroughly investigated at vanishing COM momentum $`Q`$ . The numerical effort for such calculations remains reasonable due to the high symmetry of this point. In contrast, only a few publications on multiband calculations of exciton COM dispersions in QW exist () since these are very demanding. Methods for improving the numerical accuracy and reducing the effort of such calculations are clearly necessary. Particularly useful would be an easy-to-use approach that gives the main features of the exciton dispersion with at least moderate accuracy.
The main focus of the present work lies on the exciton groundstate dispersion and its properties. A secondary goal is to study the feasibility of numerical methods for calculating the exciton groundstate dispersion in more complicated structures like V-groove quantum wires . Excitons in GaAs/Al<sub>0.3</sub>Ga<sub>0.7</sub>As symmetric QW are considered (section II). Two different methods for the dispersion calculation are used: (i) the extension for $`Q0`$ of the well-known expansion of the exciton in the product space of electron and hole subband states (section II A), and (ii) a finite-differences scheme in real space (section II B) with a groundstate-adapted Coulomb discretization . Method (i) gives high-quality numerical results but is not feasible, e.g., for quantum wires of complicated geometry. With method (ii) the groundstate exciton dispersion in a V-groove quantum wire is tractable. Its convergence properties are checked here against the results of (i). Results for V-groove quantum wires will be presented elsewhere.
Improving own previous results , we address in detail the problem of the choice of the COM transformation and introduce an optimized, groundstate adapted COM transformation that greatly enhances the numerical accuracy and stability of our results (section III).
Quasi as a by-product, a semi-analytical expression for the average exciton groundstate mass suitable for exciton localization problems is derived (section IV). This expression is of great practical importance since it gives for not too wide QW reliable mass values. The only necessary ingredients are the lowest subband dispersion and a good estimate for the groundstate Bohr radius.
Finally, we discuss the results of the exciton dispersion calculations in momentum- (section V A) and real-space (section V B). The results of our semi-analytical expression for the average exciton groundstate mass are compared to various other mass expressions as well as with the numerical dispersions separately (section V C).
## II Theoretical model
We consider the well-studied system of direct Wannier excitons in a single symmetric GaAs/Al<sub>0.3</sub>Ga<sub>0.7</sub>As type I quantum well grown in $`100`$ direction. Many aspects of the presented results can be effortlessly extended to Wannier excitons in other, more general heterostructures.
In the envelope function approximation, the Wannier exciton is described by the Hamilton operator
$$\begin{array}{ccc}H& =& H_e(\stackrel{}{r}_e)+H_h(\stackrel{}{r}_h)+V_{Coul}(\stackrel{}{r}_e\stackrel{}{r}_h).\hfill \end{array}$$
(1)
$`H_{e,h}`$ describe the material-dependent bandstructure of the respective particles in the vicinity of the $`\mathrm{\Gamma }`$-point, and $`V_{Coul}`$ stands for the attractive Coulomb interaction. We choose the coordinate system as usual with the $`z`$-axis in growth direction $`100`$; $`\stackrel{}{r}_e=(x_e,y_e,z_e),\stackrel{}{r}_h=(x_h,y_h,z_h)`$ denote the space coordinates of electron and hole, respectively. For the materials involved, the conduction band is to a good approximation parabolic; anharmonicities in the conduction band arise mainly through the interaction with the light and split-off valence bands which is small due to the relatively large band gaps. The valence band is adequately described by the Luttinger Hamiltonian in the axial approximation , which takes into account explicitly the coupling of the heavy- and light-hole bands ($`\mathrm{\Gamma }_8^v`$) but suppresses warping. The coupling to the split-off ($`\mathrm{\Gamma }_7^v`$) band can be safely neglected for subband states with energies up to approximately 50 meV from the band edge because of the relatively large energy separation. We neglect the effect of the different dielectric constants (no image charge effects) , and also all effects that lead to a small spin-splitting like lack of inversion symmetry of the bulk material or the interfaces as well as the exchange part of the Coulomb interaction . The electron spin is irrelevant and will be given a fixed value of $`+1/2`$ in the present work. The quantization axis of the electron spin and of the hole angular momentum $`J`$ is taken along the growth direction, and we use for the valence band edge states the same convention as in Ref. .
The Hamilton operator $`H`$, Eq. (1), acts within these approximations on a four-component envelope function in the product basis of the conduction and valence band edge states $`\left\{|\frac{3}{2}m_J_v|\frac{1}{2}+\frac{1}{2}_c\right\}`$, where the hole-spin projection attains values of $`m_J=+\frac{3}{2},+\frac{1}{2},\frac{1}{2},\frac{3}{2}`$:
$$\begin{array}{ccc}\hfill H_e& =& \left(\frac{\mathrm{}^2}{2m_e}(_{x_e}^2+_{y_e}^2+_{z_e}^2)+V_c\right)I,\hfill \end{array}$$
(2)
$$\begin{array}{ccc}\hfill H_h& =& \frac{\mathrm{}^2}{2m_0}\left(\begin{array}{cccc}𝒫+𝒬& & & 0\\ ^{}& 𝒫𝒬& 0& \\ ^{}& 0& 𝒫𝒬& \\ 0& ^{}& ^{}& 𝒫+𝒬\end{array}\right)+V_vI\hfill \end{array}$$
(3)
with
$`\begin{array}{ccccccc}𝒫\hfill & =& \gamma _1(_{x_h}^2+_{y_h}^2+_{z_h}^2)\hfill & & \hfill 𝒬& =& \gamma _2\left(_{x_h}^2+_{y_h}^22_{z_h}^2\right)\hfill \\ \hfill & =& i\mathrm{\hspace{0.17em}2}\sqrt{3}\gamma _3\left(_{x_h}i_{y_h}\right)_{z_h}\hfill & & \hfill & =& \sqrt{3}\frac{\gamma _2+\gamma _3}{2}(_{x_h}^2_{y_h}^2i\mathrm{\hspace{0.17em}2}_{x_hy_h}^2)\hfill \end{array}`$
and
$$\begin{array}{ccc}V_{Coul}(\stackrel{}{r}_e\stackrel{}{r}_h)& =& \frac{e^2}{ϵ}\frac{1}{\stackrel{}{r}_e\stackrel{}{r}_h^{}}I.\hfill \end{array}$$
(4)
$`I`$ is the $`4\times 4`$ unity matrix. The material parameters $`\gamma _1(z_h),\gamma _2(z_h),\gamma _3(z_h)`$ as well as the offsets $`V_v(z_h),V_c(z_e)`$ are piecewise constant functions of $`z_e,z_h`$. To ensure that the kinetic operators remain Hermitian in the presence of interfaces, we use the symmetric substitutions
$$\gamma _i(_i\gamma +\gamma _i)/2,\gamma _{ij}^2(_i\gamma _j+_j\gamma _i)/2,i=x,y,z.$$
(5)
The in-plane COM momentum $`\stackrel{}{Q}=i\mathrm{}(\stackrel{}{}_e_{}+\stackrel{}{}_h_{})`$ is a constant of motion because the interaction term (4) depends only on the relative distance of the two particles. Reflection with respect to the central $`xy`$ plane, $`\sigma _{xy}`$, is also a symmetry element for symmetric QW. Consequently, the exciton can be characterized by the parity $`P=\pm 1`$. Then, the wavefunction factorizes into
$$\mathrm{\Psi }^{\stackrel{}{Q};Pa}(\stackrel{}{r}_e,\stackrel{}{r}_h)=\frac{e^{i\stackrel{}{Q}\stackrel{}{R}}}{2\pi }\underset{m_J}{}\mathrm{\Psi }_{m_J}^{\stackrel{}{Q};Pa}(z_e,z_h,\stackrel{}{\rho })\frac{3}{2}m_J_v\frac{1}{2}\frac{1}{2}_c,$$
(6)
where $`\stackrel{}{\rho }=\stackrel{}{r}_e_{}\stackrel{}{r}_h_{}`$ is the in-plane particle distance, $`\stackrel{}{R}`$ the COM space coordinate canonically conjugate to $`\stackrel{}{Q}`$, and $`a`$ stands for the remaining quantum numbers related to the relative motion of the exciton.
The COM space coordinate $`\stackrel{}{R}`$ in (6) is not unambiguously defined because of the anharmonic dispersions of the constituent particles . The COM transformation must be linear in order to preserve the canonical commutation relations of space and momentum operators, and it has in general the form
$$\stackrel{}{R}=𝜷\stackrel{}{r}_e_{}+\left(\mathrm{𝟏}𝜷\right)\stackrel{}{r}_h_{},\stackrel{}{k}=i\mathrm{}\left(\left(\mathrm{𝟏}𝜷\right)\stackrel{}{}_e_{}𝜷\stackrel{}{}_h_{}\right).$$
(7)
In the parabolic case, the free parameter $`𝜷`$ is taken as the scalar
$$\beta _{parab}=\frac{m_e}{m_e+m_h}$$
(8)
in order that relative and COM motion completely decouple. For bulk excitons, $`𝜷`$ has been considered in the literature as a scalar, a tensor in real space , or even a spinor . We will return later on to the problem of an appropriate choice for the COM coordinate $`\stackrel{}{R}`$.
Taking into account the electron spin degeneracy, each exciton state is at least fourfold degenerate. It can be shown in a similar way as has been done for the hole subband states in Ref. that the operator $`R_\pi T`$, with the rotation by $`\pi `$ about the $`z`$ axis, $`R_\pi `$, and time-reversal, $`T`$, transforms between the degenerate states of different parity and opposite electron spin. If one combines this operator with the Pauli matrix $`\sigma _y^e`$, which flips only the electron spin, we have apart from an overall phase:
$`\mathrm{\Psi }^{\stackrel{}{Q};Pa}(\stackrel{}{r}_e,\stackrel{}{r}_h)`$ $`=`$ $`\left(\sigma _y^eR_\pi T\right)\mathrm{\Psi }^{\stackrel{}{Q};Pa}(\stackrel{}{r}_e,\stackrel{}{r}_h)`$ (9)
$`=`$ $`{\displaystyle \frac{e^{i\stackrel{}{Q}\stackrel{}{R}}}{2\pi }}{\displaystyle \underset{m_J}{}}\mathrm{\Psi }_{m_J}^{\stackrel{}{Q};Pa^{}}(z_e,z_h,\stackrel{}{\rho })\frac{3}{2}m_J_v\frac{1}{2}+\frac{1}{2}_c.`$ (10)
Comparing (6) with (9), we find:
$$\mathrm{\Psi }_{m_J}^{\stackrel{}{Q};Pa}(z_e,z_h,\stackrel{}{\rho })=\mathrm{\Psi }_{m_J}^{\stackrel{}{Q};Pa^{}}(z_e,z_h,\stackrel{}{\rho }).$$
(11)
That is, the (degenerate) state of reversed parity is obtained by inverting the order of the spin components of the exciton envelope, complex conjugating, and changing the sign of the in-plane relative coordinate. Thus changing the multiband exciton parity with fixed electron spin in symmetric QW corresponds to flipping the hole spin in the single-band exciton case. In the axial approximation and for $`Q=0`$, the different angular momentum components decouple , and changing the sign of $`\stackrel{}{\rho }`$ in Eq. (11) just changes the sign of two spin components leaving the other two unchanged; this holds no longer at $`Q0`$.
We have solved the eigenvalue problem (1) in two ways which will be discussed in turn: (i) in $`\stackrel{}{k}`$-space, expanding Eq. (6) in the product space of the electron and hole subband states , and (ii) in real space, using a finite-differences scheme. The first method gives very accurate results and is used to reveal the main features of the exciton dispersion. The second method is only suitable for the groundstate dispersion but promises to be feasible for more general structures. It is validated by comparing its results with the ones from the first method.
### A Solution in $`\stackrel{}{k}`$-space
As a first step, we calculate the single-particle subband states and their dispersions
$`H_e|n_e\stackrel{}{k}_e;\pm \frac{1}{2}`$ $`=`$ $`_{n_e}(k_e)|n_e\stackrel{}{k}_e;\pm \frac{1}{2}`$ (12)
$`H_h|n_h\stackrel{}{k}_h;p_h`$ $`=`$ $`_{n_h}(k_h)|n_h\stackrel{}{k}_h;p_h`$ (13)
using a transfer-matrix method as in . The respective solutions in the axial approximation are of the form
$`|n_e\stackrel{}{k}_e;\pm \frac{1}{2}`$ $`=`$ $`{\displaystyle \frac{e^{i\stackrel{}{k}_e\stackrel{}{r}_e_{}}}{2\pi }}e^{i\left(\pm {\scriptscriptstyle \frac{1}{2}}\right)\theta _e}\xi _{n_e}(z_e)|\frac{1}{2}\pm \frac{1}{2}_c,`$ (14)
$`|n_h\stackrel{}{k}_h;p_h=\pm 1`$ $`=`$ $`{\displaystyle \frac{e^{i\stackrel{}{k}_h\stackrel{}{r}_h_{}}}{2\pi }}{\displaystyle \underset{m_J}{}}e^{im_J\theta _h}\xi _{n_h,p_h,k_h}^{m_J}(z_h)|\frac{3}{2}m_J_v.`$ (15)
In Eqs. (13,14) $`n_{e,h}`$ denote the subband indices, $`\stackrel{}{k}_{e,h}=(k_{e,h},\theta _{e,h})`$ the respective in-plane wavevectors in polar coordinates and $`p_h`$ the hole parity under $`\sigma _{xy}`$ .
In a second step, the exciton wavefunction for a given COM momentum $`\stackrel{}{Q}`$ is expanded into
$`\mathrm{\Psi }^{\stackrel{}{Q};Pa}(\stackrel{}{r}_e,\stackrel{}{r}_h)`$ $`=`$ $`{\displaystyle \underset{n_en_h}{}}{\displaystyle 𝑑\stackrel{}{k}\phi _{n_en_h}^{\stackrel{}{Q};a}(\stackrel{}{k})|n_e\stackrel{}{k}_e;+\frac{1}{2}_c|n_h\stackrel{}{k}_h;p_h_v},`$ (16)
with subband states of the two particles combined in such a way that the resulting exciton state has the required parity $`P`$ and total momentum $`\stackrel{}{Q}`$:
$$\stackrel{}{k}_e=\stackrel{}{k}+𝜷\stackrel{}{Q},\stackrel{}{k}_h=\stackrel{}{k}\left(\mathrm{𝟏}𝜷\right)\stackrel{}{Q},P=p_h(1)^{n_e+1}.$$
(17)
The last relation reflects that the conduction subband envelopes are even (odd) for odd (even) subband index. Fixing exciton parity $`P`$ and electron spin eliminates any degeneracy at $`\stackrel{}{Q}0`$.
With the expansion (16) and the relations (17), the exciton Schrödinger equation takes the form
$`\left(_{n_e}(\stackrel{}{k}_e)+_{n_h}(\stackrel{}{k}_h)E_a^X(\stackrel{}{Q})\right)\phi _{n_en_h}^{\stackrel{}{Q};a}(\stackrel{}{k})+{\displaystyle \underset{n_e^{}n_h^{}}{}}{\displaystyle 𝑑\stackrel{}{k}^{}V_{\genfrac{}{}{0pt}{}{n_e^{}n_h^{}}{n_en_h}}^\stackrel{}{Q}(\stackrel{}{k},\stackrel{}{k}^{})\phi _{n_e^{}n_h^{}}^{\stackrel{}{Q};a}(\stackrel{}{k}^{})}=0,`$ (18)
where $`E_a^X(\stackrel{}{Q})`$ denotes the energy dispersion of the exciton state. The interaction
$$V_{\genfrac{}{}{0pt}{}{n_e^{}n_h^{}}{n_en_h}}^\stackrel{}{Q}(\stackrel{}{k},\stackrel{}{k}^{})=\frac{1}{2\pi }\frac{e^2}{ϵ}\frac{1}{|\stackrel{}{k}\stackrel{}{k}^{}|}_{\genfrac{}{}{0pt}{}{n_e^{}n_h^{}}{n_en_h}}^\stackrel{}{Q}(\stackrel{}{k},\stackrel{}{k}^{})$$
(19)
is the in-plane 2D Fourier transform of the 3D Coulomb potential modified due to the confinement in $`z`$-direction. The latter is expressed through form factors \[with Eq. (17)\]
$`_{\genfrac{}{}{0pt}{}{n_e^{}n_h^{}}{n_en_h}}^\stackrel{}{Q}(\stackrel{}{k},\stackrel{}{k}^{})`$ $`=`$ $`{\displaystyle \underset{m_J}{}}{\displaystyle 𝑑z_e𝑑z_he^{|\stackrel{}{k}\stackrel{}{k}^{}||z_ez_h|}\xi _{n_e}^{}(z_e)\xi _{n_h,\stackrel{}{k}_h}^{m_J^{}}(z_h)\xi _{n_e^{}}(z_e)\xi _{n_h^{},\stackrel{}{k}_h^{}}^{m_J}(z_h)}.`$ (20)
The above integrals are calculated analytically since the subband states obtained with the transfer matrix method are combinations of exponential and trigonometric functions.
The integrable singularity of the Coulomb potential (19) at $`\stackrel{}{k}=\stackrel{}{k}^{}`$ is taken care of by adding and subtracting in Eq. (18) the analytically integrable term
$$C(\stackrel{}{k},\stackrel{}{k}^{})=\frac{e^2}{ϵ}\frac{1}{2\pi }\left(\frac{1}{|\stackrel{}{k}\stackrel{}{k}^{}|}\frac{1}{max(k,k^{}))}\right).$$
(21)
This gives a smooth “corrected” potential (19) of small absolute magnitude.
To take benefit of the axial approximation, the exciton envelope and the form factors are expanded into 2D angular momentum eigenstates $`\mathrm{exp}(i\mathrm{}\theta )`$. The angular momentum $`\mathrm{}`$ is chosen for every subband combination such that $`\mathrm{}=0`$ corresponds to the respective $`s`$-like exciton at $`Q=0`$.
The resulting set of coupled one-dimensional integral equations is solved numerically for various values of the COM momentum $`Q`$. Results will be presented in section V.
### B Solution in real space
We demonstrated in Ref. that calculations of the multiband-exciton groundstate dispersion are also feasible with a finite-differences scheme in real space. This method is conceptually simple: the Schrödinger equation corresponding to Eqs. (1)-(4),(6) leads to a system of four (number of spin components) coupled partial differential equations in the four dimensional space $`(\stackrel{}{\rho },z_e,z_h)`$. The resulting eigenvalue problem involves a large sparse complex Hermitian matrix with a substantial number, 44, of nonzero off-diagonals. In contrast to the $`\stackrel{}{k}`$-space approach, the method can, in principle, be applied effortlessly to very general heterostructures, like quantum wires and quantum dots. The main drawback is the need for huge amounts of computer memory. Indeed, the dimension of the matrix to be diagonalized scales with the fourth power of the number of grid points per spatial dimension. The most dense grid we used led to a matrix of dimension $`N_D=4\times 31\times 31\times 61\times 61=1.410^7`$ (4 is the number of spin components).
We used the ARPACK package to calculate a few eigenvalues and eigenvectors at the lower end of the spectrum. ARPACK is an efficient implementation of the Implicitly Restarted Arnoldi Method that can be viewed as a synthesis of the Arnoldi/Lanczos process with the Implicitly Shifted QR algorithm. Storage of the (nonzero) matrix elements is not required, only a matrix-vector multiplication utility is needed. Multiple eigenvalues, as they occur in our case, offer no additional problems. However, if one needs all the members of a multiplet the iteration subspace has to be chosen large enough. We find that an iteration subspace of five times the number of the requested eigenvalues (rather than the proposed factor of two ) is usually sufficient. This matter was of no concern for our problem, since using the symmetry considerations, Eq. (11), we can derive from a calculated state also the second one of the doublet. For the largest matrices, we used a factor of three as a compromise between memory demand and CPU time usage.
The matrix resulting from the discretization is highly structured. For minimizing the memory costs and still making full use of the vector registers, we construct the matrix-vector product using auxiliary, much smaller matrices.
Since memory is critical, it is crucial for any real-space approach to optimize the convergence of the relevant quantities with the mesh size. On the one hand, we optimize the COM transformation, as will be discussed in the next section, thereby improving the handling of the kinetic terms. On the other hand, we use a groundstate-adapted discretization of the Coulomb potential which is discussed in more detail in the Appendix. The idea behind this approach is to extract the discretized interaction from a reference system that has the same interaction but a simple kinetic term, and whose groundstate is known analytically. If the groundstate of the reference system is similar enough to the one of the real system, good convergence is expected.
Calculations on a parallel-vector machine of type CRAYJ932 reached performances of 140MFlops/CPU, the peak performance of the limiting BLAS routine being 185MFlops/CPU.
## III Optimized center-of-mass transformation
We return to the ambiguity in the COM transformation (7) which is expressed in the freedom to choose $`𝜷`$. The relevance of $`𝜷`$ for accelerating numerical convergence in dispersion calculations was realized quite early for bulk excitons in , where a scalar $`\beta `$ was optimized in a trial and error procedure. However, there has been no other algorithm to take advantage of this freedom until recently . Before that, there have been just two publications where numerical multiband exciton dispersion in quantum wells were calculated: in Ref. $`\beta =1`$ (in the parabolic case $`m_e=\mathrm{}`$) was taken in order for the form factors (20) to be independent of $`Q`$, and in Ref. no particular choice or handling of $`\beta `$ is mentioned. In analytic expressions, usually the symmetric (in the parabolic case $`m_e=m_h)`$ value $`\beta =1/2`$ is taken .
The effect of the $`\beta `$ choice becomes clear when one evaluates Eq. (7) for two different values $`\beta ,\beta ^{}=\beta +\delta \beta `$ giving $`\stackrel{}{R}^{}=\stackrel{}{R}+(\delta \beta )\stackrel{}{\rho },\stackrel{}{k}^{}=\stackrel{}{k}(\delta \beta )\stackrel{}{Q}`$. Clearly, $`\beta `$ moves artificially part of the plane wave of the COM motion into the relative part of the exciton (6) or, equivalently, it shifts the relative part of the wavefunction in $`\stackrel{}{k}`$-space. A good choice of $`\beta `$, as in the parabolic case (8), keeps the relative part of the exciton in real space as smooth as possible or, equivalently, pins the relative part of the wavefunction in $`\stackrel{}{k}`$-space to the origin. This situation is illustrated in Fig. 1 where we plot the envelope of the $`HH_1C_11s`$ exciton in the single subband approximation using the symmetric value of $`\beta =1/2`$. This value of $`\beta `$ is indeed not optimal, as the large shift demonstrates. Diamonds mark where the position of the origin would be for other values of $`\beta `$.
We introduced in Ref. a quasi-analytical method for determining the optimal choice of the scalar $`\beta `$, which we briefly summarize here. It is motivated by the fact that in the parabolic case the correct COM transformation decouples the relative motion and COM motion completely. A full decoupling is not possible for non-parabolic dispersions. We looked for a choice of $`\beta `$ that decouples “as much as possible”. To quantify this, we inserted in Eq. (1) the general $`\beta `$-dependent COM transformation (7), separated the $`Q`$-dependent terms from the rest
$$\begin{array}{ccc}H\hfill & =& H^{(0)}+H^{(1)}(\beta )Q+H^{(2)}(\beta )Q^2\hfill \end{array},$$
(22)
and viewed these as a perturbation of the $`Q=0`$ exciton. Taking into account the inversion symmetry of the Brillouin zone, the kinetic mass of the groundstate exciton $`g`$ is given in second order perturbation theory by
$$\frac{\mathrm{}^2}{2M_g^X}=gH^{(2)}(\beta )g+\underset{ag}{}\frac{|aH^{(1)}(\beta )g|^2}{E_g^X(0)E_a^X(0)}.$$
(23)
However, the exciton mass must not depend on $`\beta `$. Maximizing the first order contribution in (23) and minimizing this way the strictly positive contribution of the higher states to the groundstate mass leads to the analytical result
$$\begin{array}{ccc}\beta \hfill & =& gH_h^{(2)}|g/gH_e^{(2)}+H_h^{(2)}|g\hfill \end{array}.$$
(24)
$`H_{e,h}^{(2)}`$ are simply the material-dependent coefficients of the $`\beta ^2Q^2`$ terms when inserting (17) into the $`k`$-space representation of the kinetic energies in Eq. (2) and (3), respectively .
The explicit form of Eq. (23) with the contributions from the higher states dropped and $`\beta `$ from Eq. (24) suggests to define COM-related, effective masses $`m_{e,h}^{}`$ for electron and hole
$$1/m_{e,h}^{}=\frac{2}{\mathrm{}^2}g|H_{e,h}^{(2)}|g\text{ satisfying}M_g^X=m_e^{}+m_h^{}.$$
(25)
Numerical results show that the masses obtained from (25) tend to be too small. Nevertheless, the obtained values for $`\beta `$ in were quite reasonable because of the much heavier hole mass. If one actually calculates the contributions of the higher exciton states to $`M_g^X`$ in (23) (which are dropped in (25)), one finds that the only important correction comes from the coupling to the $`LH_1C_11s`$-like state. Taking in Eq. (23) this single correction term into account gives practically the exact curvature of the exciton groundstate dispersion at $`Q=0`$.
The above procedure is not the best for determining the optimal value of $`\beta `$ as the importance of the coupling to higher states demonstrates. It was inspired by the solution of the exciton problem in real space. Let us now look at the form of the exciton wavefunction for $`Q0`$ in the subband expansion (16). The $`Q`$-dependence enters the wavefunction: (i) the need to appropriately combine the subband states to get the right $`Q`$ and (ii) through the need for the envelope to adjust for the anharmonicities in the dispersions. In the perturbation approach described above, we tried to find a COM transformation that keeps for small $`Q`$ values the *entire wavefunction* unchanged as much as possible. However, once the one-particle problem is solved, the $`Q`$-dependence due to the appropriate combination of the subband states (i) is explicitly known. Therefore, a better Ansatz for the wavefunction would be to find a COM transformation that keeps the *envelopes* as much unchanged as possible; that is Eq. (16) with
$$\phi _{n_en_h}^{\stackrel{}{Q};g}(\stackrel{}{k})=\phi _{n_en_h}^{0;g}(\stackrel{}{k}).$$
(26)
The minimization of the energy with respect to $`\beta `$ can be done analytically in the limit $`Q\mathrm{\hspace{0.17em}0}`$. We easily obtain the optimized value
$$\beta _0=\frac{{\displaystyle \underset{n_en_h}{}}{\displaystyle 𝑑\stackrel{}{k}|\phi _{n_en_h}^{0;g}(\stackrel{}{k})|^2(\widehat{Q}\stackrel{}{})^2_{n_h}(\stackrel{}{k})}}{{\displaystyle \underset{n_en_h}{}}{\displaystyle 𝑑\stackrel{}{k}|\phi _{n_en_h}^{0;g}(\stackrel{}{k})|^2(\widehat{Q}\stackrel{}{})^2(_{n_e}(\stackrel{}{k})+_{n_h}(\stackrel{}{k}))}}.$$
(27)
This expression accounts also for the dependence of $`\beta _0`$ on the direction $`\widehat{Q}`$ of the COM momentum in the case of warped valence bands. We took into account the inversion symmetry of the Brillouin zone and assumed that the Coulomb potential (19), (20) can be approximated as a function of the momentum transfer only, $`V(\stackrel{}{k},\stackrel{}{k}^{})V(\stackrel{}{k}\stackrel{}{k}^{})`$, i.e., we neglected any $`Q`$-dependence of the (Coulomb) potential energy of the groundstate. This is expected to be a good approximation since the Coulomb energy depends solely on the charge distribution, which should not be affected significantly by the in-plane motion. Indeed, it has been estimated in Ref. that the error introduced by neglecting in Eq. (20) the $`\stackrel{}{k}`$-dependence of the hole envelopes is about $`5\%`$. Our assumption should lead to even smaller deviations.
Fig. 1 demonstrates the quality of the expression (27). It shows the groundstate envelope at a rather large value of $`Q=0.5`$ nm<sup>-1</sup>, even though Eq. (27) was obtained in the limit $`Q0`$. It is particularly impressive that the respective origin for the optimized choice $`\beta _0`$ from Eq. (27) lies even a bit to the left of the maximum of the envelope. This accounts for the $`p`$-component that deforms the originally radially symmetric $`Q=0`$ envelope; using an angular momentum decomposition of the envelope a minimum number of components would be needed. The slight deformation (these are logarithmic plots) for large values of $`Q`$ is due to the anharmonicity of the one-particle dispersions.
The importance of a suitable choice of the COM transformation for the numerical convergence is illustrated in Fig. 2 for the dispersion of the $`HH_1C_11s`$ exciton of a 5 nm QW. This has been calculated in $`\stackrel{}{k}`$-space for various values of $`\beta `$ with the same basis ($`HH_1C_1,LH_1C_1,\mathrm{}=0,\pm 1,\pm 2`$). The further the used $`\beta `$ lies from the optimal value $`\beta _0`$ ($`\beta _0=0.23`$ in this case) the worse the results are. We did also calculations where for *each* $`Q`$ value an optimal value of $`\beta `$ was obtained by numerical variation. We observed deviations from $`\beta _0`$ less than $`1\%`$ near $`Q=0`$ and not larger than $`10\%`$ at $`Q=0.5`$ nm<sup>-1</sup> even for the widest well (20 nm). At large $`Q`$ the $`\beta =0`$ (in the parabolic case $`m_h=\mathrm{}`$) curve gives slightly better results than $`\beta _0`$ since the $`HH_1`$ subband dispersion gets more flat after the avoided crossing with the $`LH_1`$ subband, but it gives considerably worse results at small $`Q`$.
## IV A simple analytical formula for the average exciton kinetic mass
The analytical variation that led to Eq. (27) gives the groundstate energy up to terms quadratic in $`Q`$. The corresponding groundstate kinetic mass $`M_g^X`$ has again the form
$$M_g^X=m_e^X+m_h^X$$
(28)
with the COM-related effective masses for electron and hole defined as
$$1/m_{e,h}^X=\frac{1}{\mathrm{}^2}\underset{n_en_h}{}𝑑\stackrel{}{k}|\phi _{n_en_h}^{0;g}(\stackrel{}{k})|^2(\widehat{Q}\stackrel{}{})^2_{e,h}(\stackrel{}{k}).$$
(29)
With these masses, the expression for $`\beta _0`$, Eq. (27), has the same form as in the parabolic case (8). Eq. (29) gives the correct results for the free particle case.
We claim that this simple result will be of considerable practical importance. Eq. (29) is physically appealing: it leads to a weighted average of the subband dispersions. Further, it is relatively simple to calculate: It requires only the approximate knowledge of the $`Q=0`$ exciton envelope and of the involved subband dispersions. The numerical calculation of subband dispersions is nowadays an easy task (provided the $`\stackrel{}{k}\stackrel{}{p}`$ parameters are known). Moreover, especially for narrow QW, the envelope of the $`HH_1C_1`$ component of the groundstate exciton is to a very good approximation similar in shape to the groundstate of the 2D exciton,
$$\phi _{1s}^{2D}(\stackrel{}{\rho })=\sqrt{\frac{2}{\pi a_B^2}}e^{a_B\rho },\phi _{1s}^{2D}(\stackrel{}{k})=\sqrt{\frac{2a_B^2}{\pi }}\left(1+\left(a_Bk\right)^2\right)^{3/2}.$$
(30)
The $`LH_1C_1`$ component is quite small, e.g., $`5\%`$ for the 20 nm QW, and can be safely neglected in this context. Therefore, only a good estimate for the effective Bohr radius $`a_B`$ is needed to evaluate (29).
## V Results
We have calculated exciton dispersions both in real and momentum space for GaAs/Al<sub>0.3</sub>Ga<sub>0.7</sub>As $`001`$ QW of various widths. The coupling of heavy and light holes was fully incorporated. The values of the material parameters $`\gamma _1,\gamma _2,\gamma _3,m_e`$ were taken by linear interpolation from the GaAs and AlAs values; the offset ratio was $`V_v/V_c=0.68/0.32`$ and the band gap in meV was taken as $`E_g(x)=1519+1040x+470x^2`$, $`x`$ being the Al content. For the dielectric constant, we adopted $`ϵ=12`$ for both well and barrier material.
### A Subband expansion
The nomenclature is as follows: the exciton in the subband expansion has various $`n_hn_e`$ subband components with the corresponding envelopes $`\phi _{n_en_h}^{\stackrel{}{Q};a}(\stackrel{}{k})`$, Eq. (16). These envelopes have in the axial approximation at $`Q=0`$ a definite angular momentum $`\mathrm{}`$ and will be denoted by $`1s,\mathrm{\hspace{0.17em}2}s,\mathrm{\hspace{0.17em}2}p_\pm ,\mathrm{\hspace{0.17em}3}d_\pm `$ and so on. Each exciton state at finite $`Q`$ will be named according to the main subband component of the corresponding state at $`Q=0`$. That is, speaking of the $`HH_1C_11s`$ exciton means that at $`Q=0`$ its main subband component is the $`HH_1C_1`$ one with an $`1s`$ envelope. Similar to the single-particle hole subband states, which can change their heavy- or light-hole character away from the $`\mathrm{\Gamma }`$-point, the envelope of the main subband component or even the main subband component itself can change with increasing $`Q`$. To denote the main subband component of a state at a given value of $`\stackrel{}{Q}`$, we will speak of the character of the state at this $`\stackrel{}{Q}`$. For example, the $`HH_1C_12p_+`$ exciton has a $`HH_1C_12p_+`$ character at $`Q=0`$ and a $`LH_1C_11s`$ character $`Q0`$.
The exciton dispersions in $`k`$-space are calculated as follows: For each QW, we first calculate the exciton spectrum at $`Q=0`$. Subsequently, a 2D $`1s`$-exciton groundstate function (30) is fitted to the $`HH_1C_1`$ envelope. For the wider QW, also a two-dimensional $`3d`$-exciton function is fitted to the $`LH_1C_1`$ envelope. This fit is used, instead of the numerical envelope, to evaluate the optimized COM transformation (27) because it allows to take advantage of the analytically known derivatives of the fit function. The so calculated value of $`\beta _0`$ is used for the $`Q0`$ calculations.
In Fig. 3, we display the envelopes of the components of the groundstate exciton in the subband expansion (16) for a 15 nm wide QW at $`Q=0`$. The coupling of the higher subbands is rather small, less than $`3\%`$ for the $`LH_1C_1`$ component and even less for the others. The parity and spin selection rules for the Coulomb coupling of the subbands at the $`\mathrm{\Gamma }`$-point in symmetric QW are obeyed, e.g., the admixture of the $`LH_1C_1`$ state vanishes at the $`\mathrm{\Gamma }`$-point since the Coulomb potential is spin-diagonal and the $`HH_1`$ and $`LH_1`$ subband states are pure heavy- and light-hole states, respectively. The $`HH_1C_1`$ envelope is very well approximated by a 2D $`1s`$-exciton function. Deviations are mainly located at the vicinity of the $`HH_1LH_1`$ avoided crossing of the hole-subband dispersions (here, $`k_{ac}0.13`$ nm<sup>-1</sup>). The total in-plane probability distribution follows the form of Eq. (30) even better than the $`HH_1C_1`$ envelope alone; the coupling to the higher subbands allows the exciton to relax further. This, again, supports the notion that the subband mixing has little influence on the charge distribution.
The calculations in $`k`$-space presented in the following take into account only the two lowest hole subbands ($`n_hn_e=HH_1C_1,`$ $`LH_1C_1`$). Inclusion of higher subbands does not enhance the binding energy of the groundstate exciton considerably. For the angular decomposition of the envelope components only the $`s,p_\pm ,d_+`$ ($`\mathrm{}=0,\pm 1,2)`$ components for the $`HH_1C_1`$ and the $`s,p_\pm ,d_{}`$ ($`\mathrm{}=0,\pm 1,2)`$ components for the $`LH_1C_1`$ were considered. Due to the optimized choice of the COM coordinate system, Eq. (27), these few angular momentum components are sufficient to describe the $`HH_1C_11s`$ and $`LH_1C_11s`$ dispersions excellently over the whole range of COM momentum values considered, $`Q0.5`$ nm<sup>-1</sup>: The $`p`$ components account mainly for the deformation of the envelope, and the $`d`$ components take care of the Coulomb coupling to higher states.
In Fig. 4, we show the dispersion calculated in $`k`$-space of the first bound exciton states as well as some of the continuum states in a 15 nm wide QW. Zero of energy is the onset of the $`HH_1C_1`$ continuum at $`Q=0`$. The Bohr radius $`a_B`$ is the one obtained from the same fit of a 2D $`1s`$-exciton to the envelope of the $`HH_1C_1`$ component at $`Q=0`$ which was used to calculate the optimized $`\beta _0`$. The groundstate is $`HH_1C_11s`$. The next state is $`HH_1C_12p_+`$ which turns into $`LH_1C_11s`$ character at the avoided crossing. The $`LH_1C_11s`$ exciton is at $`Q=0`$ the fourth excited state and shows a substantial mixing with the $`HH_1C_13d_+`$ exciton. The $`HH_1C_12p_+`$ exciton has at $`Q=0`$ a slightly lower energy than the $`HH_1C_12p_{}`$ (third state at $`Q=0`$) because it couples to the $`LH_1C_11s`$ outside of the $`\mathrm{\Gamma }`$-point.
As was reported in earlier work , the exciton COM dispersions are highly non-parabolic, much like the hole subband dispersions. This is not surprising, since the dispersion of the conduction subband is parabolic and the hole mass is much larger than the electron mass. Furthermore, it has been claimed in Ref. that the exciton dispersion follows, in a good approximation, the hole subband dispersion. Although this is certainly true in the present case due to the parabolic electron dispersion and the small electron to hole mass ratio, the exciton dispersion is in principle a *two*-particle quantity. In fact, the dispersion of the groundstate exciton follows even more closely, in the studied cases within 1 meV, the electron-hole-pair continuum edge $`_{n_en_h}(\stackrel{}{Q})`$. The latter is defined by
$$_{n_en_h}(\stackrel{}{Q})=\underset{\stackrel{}{k}_e+\stackrel{}{k}_h=\stackrel{}{Q}}{min}\left\{_{n_e}(\stackrel{}{k}_e)+_{n_h}(\stackrel{}{k}_h)\right\},$$
(31)
and represents the minimal kinetic energy of a free electron-hole pair for a given subband combination $`n_en_h`$ and a given $`\stackrel{}{Q}.`$ In the independent-subband approximation, this coincides with the *exciton continuum edge*. Fig. 5(a,b) directly compare for a narrow and a wide QW the exciton groundstate dispersion and the appropriately shifted exciton continuum edge. Also shown is the respective hole dispersion. The latter lies always above the shifted exciton continuum edge since Eq. (31) implies $`_{n_en_h}(\stackrel{}{Q})_{n_e}(0)+_{n_h}(\stackrel{}{Q})`$.
The exciton groundstate dispersion is found to lie systematically below the shifted exciton continuum edge, i.e., the groundstate exciton binding energy becomes larger away from $`Q=0`$. This can be understood based on the fact that for $`Q0`$ the groundstate exciton, Eqs. (16,17), is built from hole subband states around $`(1\beta _0)Q`$. Due to the flatter subband dispersion around this point, the hole mass gets larger and the wavefunction can better adjust to the potential. The $`LH_1C_11s`$ exciton in the 20 nm wide QW, Fig. 5b, is well separated from the spectrum of the $`HH_1C_1`$ exciton. We observe that at small $`Q`$ the $`LH_1C_11s`$ dispersion lies above the respective shifted exciton continuum edge. Indeed, the $`LH_1`$ subband shows a negative mass at the $`\mathrm{\Gamma }`$-point which becomes positive near the avoided crossing. Hence, the respective exciton has to pay with extra kinetic energy in order to achieve a small COM momentum, and its binding energy is decreased. For larger $`Q`$ values it gains again some binding energy.
The enhancement of the groundstate binding energy for $`Q0`$ is particularly large when the exciton is built from hole subband states around the avoided crossings $`(1\beta _0)Q=k_{a.c.}`$. This is demonstrated in Fig. 6 for QW of various widths. Peaks are seen at the respective location of the $`HH_1`$-$`LH_1`$ avoided crossing, marked by arrows. The enhancement of the groundstate binding energy with $`Q`$ is less than $`15\%`$ and is generally larger for wider QW.
In Fig. 7, we plot the exciton dispersion of a 5 nm QW. The exciton dispersions, like the hole subband dispersions, are less anharmonic in the much narrower QW than for the QW of Fig. 4 because of the larger energy separation of the hole subbands. The avoided crossing takes place at rather large $`k`$ ($`k=0.31`$ nm<sup>-1</sup>). The shading indicates the percentage of the contribution of the $`LH_1C_1`$ subband states to the norm of the numerical eigenvectors. The very light shading of the bound exciton states confirms the small admixture of the $`LH_1C_1`$ states in the groundstate. The shading is darker in the vicinity of the avoided crossing. Small-scale intensity variations are numerical artefacts due to the finite $`k`$-space mesh. In the exciton continuum only the shading is shown. The $`LH_1C_11s`$ resonance is clearly seen starting at approximately 18 meV. The resonance is sharper at the beginning and becomes more diffuse at the avoided crossing. It lies slightly above the respective exciton continuum edge.
### B Real-space calculations
For effectively two-dimensional structures with translational symmetry like the considered symmetric QW in axial approximation and, maybe, for some highly idealized quantum wire structures, the real-space method presented in section II B can not compete with the one in $`k`$-space. However, for realistic one-dimensional structures, like V-groove and T-shaped quantum wires, this may be the only feasible approach for calculating the exciton groundstate dispersion. This is due to the high number of confining dimensions for the exciton (in quantum wires four, two for each particle): an expansion in a problem-adapted basis like the product basis of the one-particle eigenstates, for which one expects reasonable convergence, leads to four-dimensional integrals for the Coulomb interaction. An expansion in a basis where the Coulomb potential is simple will probably show a very slow convergence with basis size.
The calculations reported here are mainly to be viewed as tests of the applicability of our real-space approach and its generalization to finite-elements discretization. They are primarily compared with results obtained with the more established $`k`$-space methods. We will therefore discuss the results of the real-space calculations focusing on the convergence properties of the method. Further, wavefunction features are better visualized in real space, in particular, the electron-hole correlation in the growth direction.
For the real-space calculations at finite $`Q`$, we used the optimized $`\beta _0`$ obtained in the respective $`k`$-space calculations. We could have used equally well some other procedure to find the effective Bohr radius $`a_B`$, e.g., a variational one, or we could also have fitted a 2D $`1s`$-exciton function to the in-plane probability distribution of the previously calculated exciton groundstate at $`Q=0`$.
The Coulomb potential was discretized as described in the Appendix. The 3D $`1s`$-exciton in the four dimensional space ($`\stackrel{}{\rho },z_e,z_h`$) with $`m_e^{ref}=0.0665m_0`$, $`m_h^{ref}=0.24m_0`$ was used as reference groundstate. The value for the reference hole mass was taken from Fig. 11, discussed below, as an average value for the range of QW widths considered. This gives a reasonable reference Bohr radius of $`a_B^{ref}=12.2`$ nm; it is nearly the correct value for the in-plane motion or somewhat larger. In the confinement direction the size of the reference wavefunction is larger than the actual one (the exciton is quenched in this direction), too. As discussed in the Appendix, a reference Bohr radius as large as or somewhat larger than the actual one gives good convergence. We did test calculations with $`a_B^{ref}`$ doubled and with $`a_B^{ref}`$ halved and found a qualitatively similar behavior as in Fig. 12 in the Appendix. The integration region was 120 nm wide in the $`\stackrel{}{\rho }`$ directions and 30 nm (60 nm) wide in the $`z`$ directions for the 5 nm (20 nm) QW with a gridpoints distance of about $`a_B/6`$ in the in-plane directions.
The panel in the middle of Fig. 8 shows dispersions of the lowest exciton states (diamonds) from the real-space calculation for a 5 nm wide QW. The exciton dispersions calculated in momentum space are plotted as full curves for comparison. The groundstate binding energy is not yet fully converged: all real-space results were shifted approximately 1 meV to lower energies to match the groundstate energies at $`Q=0`$ for both methods. Numerical tests show that the gridpoint density in the growth direction ($`0.6`$ points/nm) is more critical than in the in-plane direction ($`0.6`$ points/nm). The exciton continuum edge lies also 0.5 meV too high, and the stronger confinement due to the Coulomb interaction lets us expect for the exciton a larger deviation. Nevertheless, the groundstate relative dispersion, $`_g^X(Q)_g^X(0)`$, is converged and reproduces the $`k`$-space results very well. We note that in the case of parabolic one-particle dispersions, this is an exact property of any numerical exciton dispersions. The dispersions of the excited states are not reproduced that well. This is mainly due to their larger spatial extension and smaller energy separation from each other compared to the groundstate.
The panels on the left and on the right in Fig. 8 show logarithmic contour plots of the exciton probability distribution for some characteristic states. The probability distribution is either integrated over $`z_e,z_h`$ and displayed in the $`\stackrel{}{\rho }`$ plane, or integrated over $`\stackrel{}{\rho }`$ and displayed in the $`z_ez_h`$ plane for each spin component separately. The numerically obtained wavefunctions are a linear combination of the two degenerate solutions with opposite parity, Eq. (11). They are disentangled according to parity $`P`$, and only $`P=1`$ states are displayed.
The left panel displaying the groundstate exciton at $`Q=0`$ illustrates its $`HH_1C_11s`$ character: the main spin component is $`m_J=+3/2`$ and has no nodes. The bulk of the exciton is confined in the QW but there is substantial penetration into the barrier, being stronger for the lighter electron (note the logarithmic plot). At $`Q=0.5`$ nm<sup>-1</sup> (right panels), the groundstate exciton has still $`HH_1C_11s`$ character, as is seen in the $`k`$-space calculations, Fig. 3. This does not contradict the strong mixture of heavy- and light-hole bulk states seen in the lower right panel of Fig. 8. Indeed, the exciton is built from hole subband states near $`(1\beta _0)Q`$. This point lies past the $`HH_1LH_1`$ avoided crossing. Hence, the $`HH_1`$ subband states near this point are a strong mixture of heavy- and light-hole bulk bandedge states. The stronger penetration into the barrier of the light-hole component is again related to its smaller mass.
An interesting feature is the larger confinement of the $`m_J=+\frac{3}{2}`$ component at $`Q=0.5`$ nm<sup>-1</sup> compared to $`Q=0`$. This is a consequence of the enhanced exciton binding energy. Altogether, the plots demonstrate that the total charge distribution in not altered much with increasing COM momentum; the somewhat stronger confinement of the heavy hole is at least partly canceled by the larger penetration into the barrier of the light hole. Although here not clearly resolved, the in-plane plots at $`Q=0.5`$ nm<sup>-1</sup> show the slight deformation of the originally radially symmetric wavefunction that was seen in Fig. 1. Again these deformations partly cancel each other in the sum over the spin components, and the in-plane charge distribution remains mainly symmetric. This is more clearly seen for the first excited state at $`Q=0.5`$ nm<sup>-1</sup> which has $`HH_1C_1`$-$`2p_y`$ character (upper right panel in Fig. 8). It is the $`HH_1C_12p_+`$ exciton, which is, again, slightly lower in energy than the $`HH_1C_12p_{}`$ exciton at $`Q=0`$. At large $`Q`$ the character of the $`HH_1C_1`$ envelope changes from $`2p_\pm `$ to $`2p_{y,x}`$.
One does not expect strong electron-hole correlation in the growth direction for the 5 nm QW, which is considerably narrower than the exciton Bohr radius. Here, the confining potentials are on average much stronger than the Coulomb potential and the wavefunction can not relax in this direction. Indeed, in Fig. 9(a) the probability density integrated over $`\stackrel{}{\rho }`$ does not show much correlation: the contour lines are not elongated along the diagonal, $`z_e=z_h`$. However, some correlation exists as the plot of the cut at $`\rho =0`$ in Fig. 9(b) demonstrates. This weak correlation is not included in our $`k`$-space calculations where only one conduction band was used. Its envelope does not depend on $`k`$ and consequently the $`z_e`$ coordinate can be separated, Eq. (16).
In Fig. 10, results are displayed for a 20 nm wide QW. Two set of points are shown for two different mesh sizes. The energies of the groundstate exciton are almost converged for the more dense mesh; the deviation from the results of the $`k`$-space calculations is only 0.2 meV at $`Q=0`$.
The small, but noticeable $`m_J=+1/2`$ component of the groundstate exciton $`HH_1C_11s`$ exciton in the lower left panel corresponds to the substantial admixing of the $`LH_1C_1d_{}`$ exciton seen already for the 15 nm QW in Fig. 3. The admixing is larger for wider QW due to the smaller energy separation of the respective subbands. At $`Q=0.5`$ nm<sup>-1</sup>, it has mainly bulk light-hole character ($`m_J=1/2`$). This, again, does not contradict the $`HH_1C_11s`$ character since the $`HH_1`$ subband has approximately 60% light-hole character beyond the avoided crossing with the $`LH_1`$ subband. The first excited state at $`Q=0.5`$ nm<sup>-1</sup> ($`LH_1C_11s`$ exciton) has $`LH_1C_11s`$ character even though an additional node is seen in the $`z_ez_h`$ probability distribution of the main spin component ($`m_J=+1/2`$). The envelopes of the single-particle subband states at large enough in-plane momentum show more nodes than at the $`\mathrm{\Gamma }`$-point due to the coupling of in-plane and growth directions in the Luttinger Hamiltonian. This is one of the reasons why the expansion Eqs. (16,17) gives very good results with just two subbands, while an expansions in the subband states at the $`\mathrm{\Gamma }`$-point needs more subbands for the same accuracy.
The probability distribution plots illustrate the almost vanishing penetration into the barriers for the wide QW, in contrast to the narrower QW of Fig. 8. All panels show a clear orientation of the contour lines towards the diagonal, $`z_e=z_h`$. This demonstrates the considerable electron-hole correlation in the growth direction for QW wider than one Bohr radius. Still, the stronger correlation has little impact on the energies. Recall that in perturbation theory the first order correction to the wavefunction gives only a second order correction to the energy. This justifies the usual factoring out of the dependence on the growth direction for the much lighter electron.
This, however, is not the case for the much heavier hole. It is a well known fact, that the Coulomb coupling of the hole subbands is considerable. Indeed, for the 20 nm wide QW, neglecting the Coulomb coupling of the $`HH_1C_1`$\- to the $`LH_1C_1`$-excitons leads to an error in the groundstate binding energy larger than $`10\%`$. That is, the correlation of in-plane and confinement directions for the hole is substantial.
For the 20 nm QW, a single calculation at $`Q=0.5`$ nm<sup>-1</sup> was done with $`\beta =1`$ in order to check the relevance of this parameter for the numerical accuracy in the real-space calculations, too. Indeed, the respective groundstate energy lies very far ( 43 meV ! ) above the correct value.
### C Average exciton groundstate mass
In the process of determining the optimal choice of the COM coordinate system, we derived in section III an expression for the kinetic mass of the groundstate exciton, Eqs. (28,29). Two assumptions were essential: First, the Coulomb potential, i.e., the form factors, are a function of the in-plane momentum transfer, and, second, the Ansatz (26) is valid. The numerical results of the previous sections support these assumptions.
Together with the exciton dispersions in Fig. 4, we displayed for the groundstate exciton a parabola with the exciton mass $`M_g^X`$ of section III. This mass is obviously not the one determined by the curvature of the exciton groundstate dispersions at $`Q=0`$. It is rather an average of the curvature of the groundstate dispersion in a region of size $`1/a_B`$. Indeed, this is implied by Eqs. (28,29) and the observation made in section V A that the groundstate dispersion follows closely the respective exciton continuum edge.
This exciton COM momentum region is the one important for exciton localization due to well width fluctuations, interface roughness or alloy fluctuations in QW . Indeed, the exciton averages over smaller scale fluctuations due to its finite size and feels an effective disorder potential that is spatially correlated over the Bohr radius $`a_B`$.
The dependence of various expressions for the groundstate exciton mass on the QW width $`L`$ is displayed in Fig. 11. The conduction band mass was taken material-independent $`m_e^{b,w}=0.0665`$, the shown $`L`$-dependence comes solely from the valence band. Displayed are the masses obtained by: (i) describing the hole in the diagonal Luttinger approximation $`1/m_h=P_w(\gamma _1^w+\gamma _2^w)+P_b(\gamma _1^b+\gamma _2^b)`$, where $`P_{w,b}`$ denote the probability that the hole is in the well and barrier material respectively (long dashed), (ii) taking as hole mass the subband curvature at the $`\mathrm{\Gamma }`$-point that is known analytically (dashed), (iii) using our semi-analytical expression (28,29) for the mass with a 2D $`1s`$-exciton function fitted to the envelope of the $`HH_1C_1`$ component (diamonds), and (iv) the average of the curvature of the numerical exciton dispersion weighted with the same function as (iii) (circles, the line is a guide to the eye).
The top curve displays the values for the Bohr radius that we have used for the calculation of (iii) and (iv). Arrows at either sides of the lower part of Fig. 11 mark the mass in the diagonal approximation in the well and barrier bulk materials.
Fig. 11 demonstrates the failure of the diagonal Luttinger approximation (i) to describe even the $`HH_1`$ subband curvatures at the $`\mathrm{\Gamma }`$-point due to the degeneracy of heavy- and light-hole bands in the unstrained bulk. However, even the correct single-particle subband curvatures at the $`\mathrm{\Gamma }`$-point (ii) fail to describe accurately the curvature of the groundstate exciton dispersion at $`Q=0`$. This is mainly due to the finite extension of the exciton in $`k`$-space that implies an averaging of the subband dispersions over a region of approximately $`1/a_B`$ near the $`\mathrm{\Gamma }`$-point and partly to the Coulomb coupling to higher subbands. Both effects tend to make the groundstate exciton heavier. The mass derived from the curvature of the groundstate exciton at $`Q=0`$ (not shown) lies between curves (ii) and (iii).
The numerically obtained “best” mass values (iv) show a quantitatively and qualitatively different behavior from curves (i) and (ii). For very narrow QW the subbands become flatter at the $`\mathrm{\Gamma }`$-point because of the larger penetration into the barriers where the exciton becomes again heavier, as in models (i), (ii). But, for large $`L`$ the region of the $`HH_1`$-$`LH_1`$ avoided crossing comes to a distance of approximately $`1/a_B`$ to the $`\mathrm{\Gamma }`$-point and the exciton, averaging over the flatter subband dispersion, becomes heavier.
The quality of our semi-analytical result for the average exciton mass (iii) has to be judged according to its deviation from the numerical average (iv). The non-monotonous behavior is clearly seen for the mass (iii) obtained using only the $`HH_1`$ subband dispersion and the fitted Bohr radius. The mass values (iii) are somewhat smaller ($`<`$10%) than the ones of curve (iv). This is due to the enhancement of the binding energy for $`Q0`$ that yields larger average masses (iv) than one would expect based on the one-particle subband dispersions and the $`Q=0`$ groundstate exciton. Our semi-analytical average groundstate exciton mass (iii) reaches a minimum approximately at the QW width where the maximum binding energy is reached.
The small differences between curves (iii) and (iv) demonstrate the quality of our expression (28,29) for the average exciton groundstate mass. For this reasonable and easy-to-use mass expression, only a good estimate for the in-plane Bohr radius and the dispersion of the involved single-particle subbands is needed.
In the paper by Triques and Brum , average exciton effective masses were calculated that are relevant to the formation process of excitons in two different scenarios. These masses are defined by parabolic fits within a relevant energy range of: (a) 5 meV, approximately half the exciton binding energy, in the case that the particles first relax and then form an exciton of kinetic energy lower than the binding energy and (b) 36 meV in the case that the exciton is formed very fast and relaxes initially via optical phonon emission, reaching energies below 36 meV, the energy of the GaAs-LO phonon. These energy ranges translate to $`Q`$ values in general much larger than $`1/a_B`$. Therefore, these average masses should be larger than those of the present work. However, the values published in for the 5 meV mass ($`0.2m_0M_g^X0.3m_0`$ for $`L10`$ nm) are smaller than ours for narrow QW. Our average masses of scenario (a) are not falling below 0.3$`m_0`$ and show a smooth minimum for a QW of width somewhere between 2 nm and 5 nm. This difference can be traced back mainly to the inefficiency for narrow QW of their method involving an expansion in the subband states at the $`\mathrm{\Gamma }`$-point, as already remarked by the authors themselves. Indeed, for narrow QW the few states at the $`\mathrm{\Gamma }`$-point can not provide the needed flexibility to simulate states far away from the $`\mathrm{\Gamma }`$-point. This is substantiated by the lower left panel in Fig. 8 showing the groundstate exciton at $`Q=0.5`$ nm<sup>-1</sup>for the 5 nm QW. The $`z_ez_h`$-plots show a substantial $`m_J=+\frac{1}{2}`$ spin component. In an expansion of the $`z_h`$ dependence in the hole subband states at the $`\mathrm{\Gamma }`$-point this component would need a $`LH_2`$ subband to be described efficiently. However, for the 5 nm QW only the $`HH_1,LH_1,HH_2`$ subbands exist below the top of the barrier.
In Fig. 5(c), we compare the results of Ref. with ours for the 5 nm QW. For this narrow QW, their exciton and $`HH_1`$ subband dispersions are much steeper than ours. The discrepancy is not due to the different parameters as a comparison with the subband dispersion calculated exactly with the transfer matrix method for the parameters of Ref. demonstrates (diamonds).
We remarked already, that different parts of the non-parabolic groundstate exciton dispersion are relevant for different physical process, and we gave exciton localization and cooling of a non-thermal exciton population as examples. We would like to discuss two more possible experimental consequences of anharmonicities using Fig. 4 as illustration. First, an exciton population loosing energy by acoustic phonon emission at low temperatures could experience a bottle-neck effect, i.e., an increased population of the $`k`$-space region near $`0.15`$ nm<sup>-1</sup>. Scattering closer towards the $`\mathrm{\Gamma }`$-point will be suppressed by the decreased density of final states. A second, more directly observable consequence is the temperature dependence of the exciton lifetime. Evaluation of the latter along the lines of Ref. with the groundstate dispersion of Fig. 4 yields a super-linear increase of the exciton lifetime with temperature: The slope increases in the range from 5 to 20 K by a factor 2.2 (not shown).
## VI Concluding remarks
In summary, we performed $`\stackrel{}{k}\stackrel{}{p}`$ multiband exciton dispersion calculations of high accuracy even for very large COM momentum and narrow QW using the well known expansion in subband states. For the high quality of the numerical results, an optimized COM transformation was essential. The optimized COM transformation is related to an average exciton kinetic mass for which a simple semi-analytical expression is derived. This eliminates for many practical purposes the need to actually calculate the groundstate dispersion. The groundstate exciton dispersion is found to follow closely the respective exciton continuum edge.
In addition, we demonstrated that multiband groundstate exciton dispersion calculations are feasible in real space using a finite-differences scheme. Besides the essential optimized COM transformation, a convenient groundstate-adapted discretization of the Coulomb potential enhances numerical accuracy. This method promises to give results for systems where the first one is not practicable. Its generalization for finite-element schemes is straightforward.
## Acknowledgments
This work has been funded by the Deutsche Forschungsgemeinschaft in the frame of SFB 296. We wish also to thank the Rechenzentrum of the Humboldt University and the Konrad-Zuse-Zentrum in Berlin for their support and access to their Cray J932 (project bvph08as).
## Appendix: Discretization of the Coulomb potential
In order to attain a simple discretization for the interaction of some Hamilton operator, Glutsch, Chemla, and Bechstedt proposed to discretize on the same mesh another operator whose groundstate is analytically known with the same interaction but with a simple kinetic term (mass $`m^{ref}`$). We will call this the reference system. For illustration, consider a simple one-dimensional system with known groundstate wavefunction $`g(x)`$ of energy $`E_g^{ref}`$. The Schrödinger equation of the reference system discretized on a mesh $`x_i=i\mathrm{\Delta }_x`$ reads
$$\frac{\mathrm{}^2}{2m^{ref}\mathrm{\Delta }_x^2}\left(g(x_{i+1})2g(x_i)+g(x_{i1})\right)+V(x_i)g(x_i)=E_g^{ref}g(x_i)$$
(32)
and yields the groundstate adapted discretization
$$V(x_i)=E_g^{ref}+\frac{\mathrm{}^2}{2m^{ref}\mathrm{\Delta }_x^2}\frac{g(x_{i+1})2g(x_i)+g(x_{i1})}{g(x_i)},g(x_i)0.$$
(33)
This procedure is very simple, easy to implement, cheap to calculate, and gives for the reference system always the correct groundstate energy, regardless how inappropriate the mesh is. In addition, no special handling for potentials with integrable singularities is needed.
Using this discretization for a reference state similar enough to the one sought, one can expect good convergence with mesh size. In order to check how the dependence of the potential discretization on the reference state influences the results, we performed calculations of the ideal 2D and 3D excitons using discretizations of the Coulomb potential based on reference excitonic groundstates with various Bohr radii. Fig. 12 shows the lowest numerical eigenvalues as a function of gridpoint density. One can see that the correct estimation of the “unknown” groundstate is not so critical: a reference Bohr radius $`a_B`$ within a factor of two from the actual one, $`a_{B0}`$, still gives good results for the groundstate for reasonable mesh densities. Apparently, especially for the 2D case, it is better to choose $`a_B`$ rather larger than smaller in order to get good results also for the excited states.
We also show for the 2D case the results obtained with the potential integrated analytically on each Cartesian mesh element, using
$$\frac{dxdy}{r}=y\mathrm{ln}\left(x+\sqrt{x^2+y^2}\right)+x\mathrm{ln}\left(y+\sqrt{x^2+y^2}\right).$$
(34)
This gives also a very good convergence with the mesh size. The 2D result obtained for the groundstate with the potential integrated analytically only at the origin, and taking everywhere else its value at each gridpoint (’naive’ discretization) is included in Fig. 12. Since both discretizations are the same at the origin the difference does not originate from the divergence at this point. The superiority of the discretizations (33) and (34) is obvious. We therefore expect Eq. (33) to yield good results even with not very dense meshes.
For the 3D case a similar expression to Eq. (34) for the Coulomb potential integrated analytically on a rectangular box can be derived. However, we could not use this result to attain an alternative discretization for the real-space QW calculations because the natural mesh $`(\rho _x,\rho _y,z_e,z_h)`$ is not Cartesian in the relative coordinate $`z_ez_h`$. |
warning/0001/hep-th0001108.html | ar5iv | text | # Aharonov-Bohm Effect in Cyclotron and Synchrotron Radiations
## I Introduction
Aharonov-Bohm (AB) effect plays an important role in quantum theory refining the status of electromagnetic potentials in this theory. First this effect was discussed in relation to a study of interaction between a non-relativistic charged particle and an infinitely long and infinitesimally thin magnetic solenoid field<sup>*</sup><sup>*</sup>*A similar effect was discussed earlier by Ehrenberg and Siday (further AB field). It was discovered that particle wave functions vanish at the solenoid line. In spite of the fact that the magnetic field vanishes out of the solenoid, the phase shift in the wave functions is proportional to the corresponding magnetic flux . A non-trivial particle scattering by the solenoid is interpreted as a possibility for quantum particles to ”feel” potentials of the corresponding electromagnetic field. Indeed, potentials of AB field do not vanish out of the solenoid. A number of theoretical works and convinced experiments was done to clarify and prove the existence of AB effect. The detailed exposition of this activity can be encountered, for example, in . In particular, it was shown that AB scattering is accompanied by an electromagnetic radiation. Pair creation by a photon in the presence of AB field was calculated in . The interaction between electron spin and AB field leads to Dirac wave functions that do not vanish at the solenoid line. Thus, the issue of spin changes slightly the interpretation of AB effect. Theoretical study of AB scattering for spinning particles was presented in many papers, see for example and . AB effect was also discussed in connection with fractional spin and statistics and with cosmic strings . AB effect in anyon scattering was considered in ; radiative corrections to the effect were calculated in . AB scattering within the Chern-Simons theory of scalar particles was studied in . There exist impressive applications of AB effect in solid state physics .
A splitting of Landau levels in a superposition of parallel uniform magnetic field and AB field (further magnetic-solenoid field) gives an example of AB effect for bound states. First, exact solutions of Schrödinger equation in the magnetic-solenoid field (non-relativistic case) were studied in . Then these solutions were used in to discuss AB effect.
It is well-known that a charged particle irradiates moving in a uniform magnetic field. The corresponding radiation is called cyclotron one (CR) in the non-relativistic case; it is called synchrotron radiation (SR) in the relativistic case. In the present article we study how the presence of AB field affects CR and SR. It is clear that classical trajectories do not feel the presence of AB field whenever they do not intersect the solenoid. Thus, from the classical point of view, CR and SR are not affected by the presence of AB field. However, the latter field changes quantum trajectories, thus we expect that in the framework of quantum theory characteristics of CR and SR may be affected by such a field. We calculate spontaneous one-photon radiation of a particle (both spinless and spinning) in the magnetic-solenoid field in the framework of quantum theory. We consider from the beginning quantum relativistic problem in order to analyze consistently both relativistic (SR) and non-relativistic (CR) cases. One ought to mention that conventional CR and SR were studied in detail in numerous works (see and Refs. there). The analysis of the radiation in the magnetic-solenoid field is much more complicated and contains many new aspects and technical details.
The article is organized in the following way: In Sect.II we present exact solutions of Klein-Gordon and Dirac equations in the magnetic-solenoid field and analyze the energy spectrum of particles in such a field. In Sect.III matrix elements of transitions (both for spinless and spinning particles) with one-photon radiation are calculated exactly. In Sect.IV we analyze frequencies of the radiation. Spinless particle radiation is studied in detail in Sect.V. Here an exact expression for the radiation intensity is obtained. The non-relativistic approximation, semiclassical approximation, and weak magnetic field limit are considered. Besides, we reveal some important peculiarities of angular distribution of the radiation. Results that were obtained for spinless particle radiation are generalized to spinning particle case in Sect.VI. Particular emphasis is given to electron transitions that cause low frequency (less that the basic synchrotron frequency) radiation. In the end, we summarize results focusing our attention on manifestations of AB effect in CR and SR.
## II Relativistic particle in magnetic-solenoid field
As was mentioned above, the magnetic-solenoid field is a superposition of a constant uniform magnetic field of strength $`H`$ directed along the axis $`z`$ and a solenoid field (AB field). The latter field is created by an infinitely long and infinitesimally thin solenoid situated along the same axis $`z.`$ The solenoid creates a finite magnetic flux $`\mathrm{\Phi }`$ along the axis $`z`$. The magnetic-solenoid field is given by electromagnetic potentials of the form
$`A_1`$ $`=`$ $`x^2\left({\displaystyle \frac{\mathrm{\Phi }}{2\pi r^2}}+{\displaystyle \frac{H}{2}}\right),A_2=x^1\left({\displaystyle \frac{\mathrm{\Phi }}{2\pi r^2}}+{\displaystyle \frac{H}{2}}\right),A_0=A_3=0,`$ (1)
$`r^2`$ $`=`$ $`(x^1)^2+\left(x^2\right)^2,x^0=ct,𝐱=(x^i),x^1=x,x^2=y,x^3=z.`$ (2)
The potentials (2) define the magnetic field $`𝐇`$of the form
$$𝐇=(0,0,H_z),H_z=H+\mathrm{\Phi }\delta (x^1)\delta (x^2).$$
(3)
It is convenient to present the magnetic flux $`\mathrm{\Phi }`$as
$$\mathrm{\Phi }=(l_0+\mu )\mathrm{\Phi }_0,\mathrm{\Phi }_0=2\pi c\mathrm{}/\left|e\right|,\mathrm{\hspace{0.33em}0}\mu <1,l_0Z.$$
(4)
The integer $`l_0`$ gives a number of quanta $`\mathrm{\Phi }_0`$ in the total flux $`\mathrm{\Phi }`$. The quantity $`\mu `$ will be called the mantissa of the magnetic flux $`\mathrm{\Phi }.`$In the cylindrical coordinates $`r,\phi ,`$
$$x^1=r\mathrm{cos}\phi ,x^2=r\mathrm{sin}\phi ,\rho =\frac{\gamma r^2}{2},\gamma =\frac{|eH|}{c\mathrm{}}>0,$$
(5)
the non zero potentials have the form
$$\frac{|e|}{c\mathrm{}}A_1=\frac{l_0+\mu +\rho }{r}\mathrm{sin}\phi ,\frac{|e|}{c\mathrm{}}A_2=\frac{l_0+\mu +\rho }{r}\mathrm{cos}\phi .$$
(6)
Doing a transformation of relativistic wave functions $`\mathrm{\Psi }(x)=e^{il_0\phi }\stackrel{~}{\mathrm{\Psi }}(x),`$ we can eliminate $`l_0`$ dependence from the potentials. Indeed, electromagnetic potentials enter in relativistic wave equations via operators of momenta $`\widehat{P}_\mu =i\mathrm{}_\mu \frac{e}{c}A_\mu `$. Thus, equations for $`\stackrel{~}{\mathrm{\Psi }}(x)`$ contain momentum operators of the form
$`e^{il_0\phi }P_\mu e^{il_0\phi }=\mathrm{}\left(i_\mu +\overline{A}_\mu \right),\overline{A}_0=\overline{A}_3=0,`$ (7)
$`\overline{A}_1={\displaystyle \frac{\mu +\rho }{r}}\mathrm{sin}\phi ,\overline{A}_2={\displaystyle \frac{\mu +\rho }{r}}\mathrm{cos}\phi .`$ (8)
Therefore, the functions $`\stackrel{~}{\mathrm{\Psi }}(x)`$ depend on the mantissa of the magnetic flux only.
Consider first solutions of Klein-Gordon equation
$$\left(\widehat{P}^2m_0^2c^2\right)\mathrm{\Psi }(x)=0,\widehat{P}_\mu =i\mathrm{}_\mu \frac{e}{c}A_\mu $$
(9)
in the solenoid-magnetic field. The operators $`\widehat{P}_0,\widehat{P}_3,`$ and $`\widehat{L}_z=x^2p_1x^1p_2=i\mathrm{}_\phi `$are integrals of motion in the case under consideration ( $`\widehat{L}`$ is angular momentum operator). We are looking for solutions of (9) that are eigenvectors for these operators,
$$\widehat{P}_0\mathrm{\Psi }=\mathrm{}k_0\mathrm{\Psi },\widehat{P}_3\mathrm{\Psi }=\mathrm{}k_3\mathrm{\Psi },\widehat{L}_z\mathrm{\Psi }=\mathrm{}\left(ll_0\right)\mathrm{\Psi },lZ.$$
(10)
The integer $`l`$is called the azimuthal quantum number. As a consequence of (10) and (9), we have
$$\widehat{P}_r^2\mathrm{\Psi }=\mathrm{}^2k^2\mathrm{\Psi },\widehat{P}_r^2=\widehat{P}_1^2+\widehat{P}_2^2,k_0^2=m^2+k_3^2+k^2,m=\frac{m_0c}{\mathrm{}}.$$
(11)
Solutions of the equations (9), (10) can be written as
$$\mathrm{\Psi }\left(x\right)=e^{i\mathrm{\Gamma }}\psi \left(\rho \right),\mathrm{\Gamma }=k_0x^0+k_3x^3+\left(l_0l\right)\phi ,$$
(12)
where the functions$`\psi \left(\rho \right)`$obey the equation
$$\rho \psi ^{\prime \prime }+\psi ^{}+\left[\overline{n}+\frac{1}{2}\frac{\left(\overline{l}+\rho \right)^2}{4\rho }\right]\psi =0,k^2=2\gamma \left(\overline{n}+\frac{1}{2}\right),\overline{l}=l+\mu .$$
(13)
Solutions of this equation can be expressed via Laguerre functions$`I_{n,m}(x).`$The latter functions are defined (for any complex $`n,m,x,n1,2,3,\mathrm{})`$ by the relation
$$I_{n,m}(x)=\sqrt{\frac{\mathrm{\Gamma }(1+n)}{\mathrm{\Gamma }(1+m)}}\frac{\mathrm{exp}(x/2)}{\mathrm{\Gamma }(1+nm)}x^{\frac{nm}{2}}\mathrm{\Phi }(m,nm+1;x),$$
(14)
where $`\mathrm{\Phi }(\alpha ,\beta ;x)`$is the confluent hypergeometric function (, 9.210). The Laguerre functions$`I_{\alpha +m,m}(x)`$ are quadratically integrable on the interval $`x0`$ whenever $`m=0,1,2,\mathrm{},`$ and $`Re\alpha >1.`$ These functions form a complete and orthonormal set on the interval $`x0`$ whenever $`m=0,1,2,\mathrm{},`$ and $`Im\alpha =0,\alpha >1.`$ Namely,
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}I_{\alpha +n,n}(x)I_{\alpha +n,n}(y)`$ $`=`$ $`\delta \left(xy\right),x,y>0,`$ (15)
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}I_{\alpha +n,n}(x)I_{\alpha +m,m}(x)𝑑x`$ $`=`$ $`\delta _{m,n},\alpha >1,n,m=0,1,2,\mathrm{}.`$ (16)
It is a matter of direct verification (using an equation for the confluent hypergeometric functions) to prove that a general solution of the differential equation
$$4x^2I^{\prime \prime }+4xI^{}\left[x^22x\left(1+s+n\right)+\left(sn\right)^2\right]I=0$$
(17)
has the form $`I=c_1I_{s,n}+c_2I_{n,s}.`$ The functions $`I_{s,n},`$and $`I_{n,s}`$are linearly independent for $`snZ.`$Otherwise we have
$$I_{n,s}=\left(1\right)^{ns}I_{s,n},snZ.$$
(18)
Let $`m`$ be an integer and non-negative; then Laguerre functions are connected to the Laguerre polynomials $`L_m^\alpha (x)`$ by the relation ( , 8.970)
$$I_{\alpha +m,m}(x)=\sqrt{\frac{\mathrm{\Gamma }(1+m)}{\mathrm{\Gamma }(1+\alpha +m)}}\mathrm{exp}(x/2)x^{\frac{\alpha }{2}}L_m^\alpha (x),m=0,1,2,\mathrm{}.$$
(19)
Taking the above information into account, we can see that bounded and quadratically integrable solutions of Eq.(13) are divided in two types, $`\psi _{n,l}^{(j)}(r),j=1,2,`$
$`\psi _{n,l}^{(1)}(r)`$ $`=`$ $`I_{\overline{n},\overline{n}\overline{l}}(\rho ),\overline{n}=n+\mu ,\mathrm{\hspace{0.33em}0}ln,`$ (20)
$`\psi _{n,l}^{(2)}(r)`$ $`=`$ $`I_{\overline{n}\overline{l},\overline{n}}(\rho ),\overline{n}=n,l<0,nZ.`$ (21)
The states of the first type ($`j=1`$) correspond to the energy spectrum of the form
$$k_0^2=m^2+k_3^2+2\gamma (n+\mu +\frac{1}{2}),\mathrm{\hspace{0.33em}0}ln,$$
(22)
and ones of the second type ($`j=2`$) correspond to the following spectrum
$$k_0^2=m^2+k_3^2+2\gamma (n+\frac{1}{2}),l<0.$$
(23)
The integer $`n0`$ is referred to as the principle quantum number. Note that the spectrum (23) of the second type states corresponds exactly to the spectrum of spinless particles in a uniform magnetic filed. The spectrum (22) is deformed by the presence of the solenoind field whenever $`\mu 0.`$ Thus, the solenoind field partially lifts a degeneracy of the magnetic field spectrum with respect to the quantum number $`l`$ whenever $`\mu 0.`$ Namely, in the general case, the particle energy spectrum in the magnetic-solenoid field depends on sign$`l.`$
In accordance with Eq. (21), it is convenient to define an effective quantum number$`\overline{n}`$ by the relation
$$\overline{n}=n+\mu \left(2j\right)=\{\begin{array}{c}n+\mu ,j=1\\ n,j=2,\end{array},n=0,1,2,\mathrm{}.$$
(24)
Then Eqs. (22), (23) can be integrated into a single formula
$$k_0^2=m^2+k_3^2+2\gamma (\overline{n}+\frac{1}{2}),\overline{l}\overline{n}.$$
(25)
We stress that the solutions (21) vanish at $`r=0.`$ That allows us to speak about AB effect in the case under consideration whenever $`\mu 0.`$
Similar to Klein-Gordon equation, the Dirac one (in the magnetic-solenoid field)
$$(\gamma ^\mu \widehat{P}_\mu m_0c)\mathrm{\Psi }(x)=0$$
(26)
admits $`\widehat{P}_0,\widehat{P}_3`$ to be integrals of motion. Besides, $`\widehat{J}_z=\widehat{L}_z+\frac{\mathrm{}}{2}\mathrm{\Sigma }_3(\widehat{J}`$ is the total angular momentum operator and $`𝚺=\mathrm{diag}(\sigma ,\sigma ))`$ is an integral of motion as well. Thus, we are looking for solutions of (26) that are eigenvectors for these integrals of motion,
$$\widehat{P}_0\mathrm{\Psi }=\mathrm{}k_0\mathrm{\Psi },\widehat{P}_3\mathrm{\Psi }=\mathrm{}k_3\mathrm{\Psi },\widehat{J}_z\mathrm{\Psi }=\mathrm{}(ll_0\frac{1}{2})\mathrm{\Psi },.$$
(27)
Solutions of Eqs. (26), (27) can be written in the form
$$\mathrm{\Psi }(x)=N_D\mathrm{exp}(i\mathrm{\Gamma })\left(\begin{array}{c}e^{i\phi }c_1\psi _{n1,l1}^{(j)}\left(\rho \right)\\ ic_2\psi _{n,l}^{(j)}\left(\rho \right)\\ e^{i\phi }c_3\psi _{n1,l1}^{(j)}\left(\rho \right)\\ ic_4\psi _{n,l}^{(j)}\left(\rho \right)\end{array}\right),$$
(28)
where $`N_D`$ is a normalization factor and the constant bispinor $`C=(c_a,a=1,2,3,4)`$ is subjected to the following algebraic system of equations (we use the standard representation for $`\gamma `$-matrices)
$$AC=0,A=\gamma ^0k_0+\gamma ^3k_3\sqrt{2\gamma \overline{n}}\gamma ^1m.$$
(29)
The system (29) has a nontrivial solution whenever
$$detA=\left(k_0^2m^2k_3^22\gamma \overline{n}\right)^2=0.$$
(30)
It follows from (30) that the rank of the matrix $`A`$ equals $`2`$. Thus, a nontrivial general solution of Eqs. (29) contains two arbitrary constants and can be written in the following block form via an arbitrary spinor $`\upsilon `$,
$$C=\left(\begin{array}{c}\left(k_0+m\right)\upsilon \\ \left(\sqrt{2\gamma \overline{n}}\sigma _1k_3\sigma _3\right)\upsilon \end{array}\right),C^+C=2k_0\left(k_0+m\right)\upsilon ^+\upsilon .$$
(31)
As in the spinless particle case, we have here two types of states ($`j=1,2`$). The energy spectrum of spinning particles in the magnetic-solenoid field follows from (30),
$$k_0^2=m^2+k_3^2+2\gamma \overline{n}.$$
(32)
States of the second type (with $`j=2)`$ have one spin orientation only whenever $`n=0.`$ Indeed, in such a case we must set $`c_1=c_3=0.`$ Thus,
$$\upsilon =\left(\begin{array}{c}0\\ 1\end{array}\right),n=0,j=2.$$
(33)
In this case, the wave functions (28) are eigenvectors for the operator $`\mathrm{\Sigma }_3`$ ($`\mathrm{\Sigma }_3\mathrm{\Psi }=\mathrm{\Psi }`$) with the eigenvalue $`1`$ (the electron spin is always opposite to the magnetic field). That fact is well-known in the absence of the solenoid field .
The states of the first type ($`j=1`$) vanish at $`r=0`$ whenever $`l0.`$ For $`l=0,`$ $`\mu 0,`$ these states become singular at $`r=0.`$ However they still can be normalized to a $`\delta `$function. The states of the second type ($`j=2`$) vanish at $`r=0`$ for any $`l.`$
The arbitrary constant spinor $`\upsilon `$ from (31) can be specified by an appropriate choice of spin integrals of motion . In what follows we are going to write $`\upsilon `$ as
$$\upsilon =\frac{1}{2}\left(\begin{array}{c}1+\zeta \\ 1\zeta \end{array}\right),\zeta =\pm 1.$$
(34)
In this case, $`\zeta =+1`$ corresponds to the spin along the magnetic field, and $`\zeta =1`$corresponds to the spin opposite to the magnetic field.
Finally, we briefly review classical motion of a charged particle in the magnetic solenoid field. That is useful for an interpretation of quantum numbers in the problem under consideration. Suppose we consider classical trajectories that do not intersect $`z`$ axis. Such trajectories are not affected by the solenoid field and have the form
$`x^0={\displaystyle \frac{k_0}{m}}\tau ,x^1=R\mathrm{cos}\kappa +x_{\left(0\right)}^1,x^2=R\mathrm{sin}\kappa +x_{\left(0\right)}^2,x^3={\displaystyle \frac{k_3}{m}}\left(\tau \tau _0\right),`$ (35)
$`\kappa =\omega _0\tau +\phi _0,\omega _0={\displaystyle \frac{\gamma }{m}},k_0^2=m^2+k_3^2+\gamma ^2R^2.`$ (36)
Here$`\tau `$is the relativistic interval and $`R,\phi _0,x_{\left(0\right)}^1,x_{\left(0\right)}^2,\tau _0,k_0,k_3`$ are integration constants. Classical analogs of quantum operators $`\widehat{P}_\mu `$ and $`\widehat{L}_z`$ read
$`P_0=\mathrm{}k_0,P_1=\mathrm{}\gamma R\mathrm{sin}\kappa ,P_2=\mathrm{}\gamma R\mathrm{cos}\kappa ,P_3=\mathrm{}k_3,`$ (37)
$`L_z=\mathrm{}{\displaystyle \frac{\gamma }{2}}\left(R^2R_0^2\right)\mathrm{}\left(l_0+\mu \right),R_0^2=\left(x_{\left(0\right)}^1\right)^2+\left(x_{\left(0\right)}^2\right)^2.`$ (38)
On the plane $`z=\mathrm{const},`$ the trajectories (36) are circles $`\left(x^1x_{\left(0\right)}^1\right)^2+\left(x^2x_{\left(0\right)}^2\right)^2=R^2`$ of radius $`R.`$ The motion along the axis $`z`$ is uniform with the velocity $`v_3=ck_3/k_0.`$ Comparing the classical radial momentum $`P_r^2=P_1^2+P_2^2`$ with the corresponding quantum expressions (11), (13), we get
$$R^2=\frac{2\overline{n}+1}{\gamma }.$$
(39)
This equation relates the principal quantum number to the radius $`R`$ of the classical motion. Comparing $`L_z`$ from (38) with the corresponding quantum expression (10), we find
$$\overline{l}=l+\mu =\frac{\gamma }{2}\left(R^2R_0^2\right).$$
(40)
Thus, we can conclude that classical trajectories with $`l\mu `$ embrace the solenoid ($`R^2>R_0^2`$) and ones with $`l<\mu `$ do not. In quantum theory these conditions are $`l0`$ and $`l<0`$ respectively. A minimal distance $`\mathrm{\Delta }R`$ between a classical trajectory and the solenoid is related to $`l+\mu `$ as follows
$$\mathrm{\Delta }R=\left|RR_0\right|=\frac{2|l+\mu |}{\gamma \left(R+R_0\right)}.$$
(41)
Thus, in fact, the absolute value of $`l`$ specifies the above distance.
Trajectories with $`l=0`$ and$`l=1`$ pass most close to the solenoid. In the first case they embrace the solenoid and in the second one do not. As was already mentioned above, Dirac wave functions with $`l=0`$ are singular at $`r=0.`$ Bearing in mind the classical interpretation of such trajectories, we may treat the existence of the singularity as a result of a superstrong interaction between the electron spin and the solenoid.
The wave functions (28) with $`N_D=[8\pi Lk_0(k_0+m)/\gamma ]^{1/2}`$obey the following orthonormality relations ($`L<z<L`$, $`L\mathrm{}`$)
$$(\mathrm{\Psi }_{n^{},l^{},k_3^{}},\mathrm{\Psi }_{n,l,k_3})=\mathrm{\Psi }_{n^{},l^{},k_3^{}}^+\mathrm{\Psi }_{n,l,k_3}𝑑𝐫=\delta _{n,n^{}}\delta _{l,l^{}}\delta _{k_3,k_3^{}}(v^{}{}_{}{}^{+}v).$$
(42)
## III Matrix elements of electron transitions with one photon radiation
In QED, one-photon radiation intensity caused by electron transitions is given by the expression
$`W_\lambda ={\displaystyle \frac{ce^2}{2\pi }}{\displaystyle 𝑑𝜿\delta \left(\kappa +k_0^ak_0^b\right)\left|\overline{𝜶\text{ }}𝐞_\lambda \right|^2},`$ (43)
$`𝜿=\kappa (\mathrm{sin}\theta \mathrm{cos}\phi ^{},\mathrm{sin}\theta \mathrm{sin}\phi ^{},\mathrm{cos}\theta ),`$ (44)
where $`𝜿`$ is photon wave vector . Spherical angles $`\theta ,\phi ^{}`$ define angular distribution of the emitted photons and $`\kappa =\left|𝜿\right|`$ defines the frequency $`\omega =c\kappa `$ and the energy $`E_{ph}=c\mathrm{}\kappa `$ of a photon. The quantities $`k_0^a,k_0^b`$ are related to electron energies $`E^a,E^b`$ in initial and final states as $`E^{a,b}=chk_0^{a,b}`$. Unit vectors $`𝐞_\lambda `$ characterize radiation polarization, see for example . $`\overline{𝜶}`$ denotes a matrix element of the operator $`𝜶=\left(\alpha ^i=\gamma ^0\gamma ^i\right),`$
$$\overline{𝜶\text{ }}=𝑑𝐱\mathrm{\Psi }_a^+\left(x\right)e^{i𝜿𝐱}𝜶\mathrm{\Psi }_b\left(x\right).$$
(45)
For spinless particle case, one has to replace $`\overline{\alpha }`$ by $`\overline{𝐏},`$
$$\overline{𝐏}=𝑑𝐱\mathrm{\Psi }_a^{}\left(x\right)e^{i𝜿𝐱}\widehat{𝐏}\mathrm{\Psi }_b\left(x\right),$$
(46)
where $`\widehat{𝐏}=\left(\widehat{P}^i\right).`$ To get total intensity of the polarized radiation, we have to sum (44) over all the final states of the electron. In the SR theory , a linear polarization of the radiation is described by $`\sigma `$ and $`\pi `$ components of the operator $`\widehat{𝐏},`$
$$\widehat{P}_\sigma =\widehat{P}_1\mathrm{sin}\phi ^{}\widehat{P}_2\mathrm{cos}\phi ^{},\widehat{P}_\pi =(\widehat{P}_1\mathrm{cos}\phi ^{}+\widehat{P}_2\mathrm{sin}\phi ^{})\mathrm{cos}\theta +\widehat{P}_3\mathrm{sin}\theta .$$
(47)
Using Eq. (6), these components can be written as
$`\widehat{P}_\sigma =i\mathrm{}\sqrt{{\displaystyle \frac{\gamma \rho }{2}}}\left(AB\right),\widehat{P}_\pi =i\mathrm{}\sqrt{{\displaystyle \frac{\gamma \rho }{2}}}\left(A+B\right)+i\mathrm{}\mathrm{sin}\theta {\displaystyle \frac{}{x^3}},`$ (48)
$`A`$ $`=`$ $`e^{i\left(\phi ^{}\phi \right)}\left({\displaystyle \frac{\rho +l_0+\mu i_\phi }{2\rho }}_\rho \right),B=e^{i\left(\phi ^{}\phi \right)}\left({\displaystyle \frac{\rho +l_0+\mu i_\phi }{2\rho }}+_\rho \right)`$ (49)
Consider first the spinless particle case. Thus, we have to substitute (49) and functions (12), (21) into (46). We are going to mark off quantum numbers of final states by primes. The integration over $`x^3`$ in (46) leads to a conservation low for $`z`$component of the momentum
$$k_3k_3^{}=\kappa \mathrm{cos}\theta .$$
(50)
Integrating over $`\kappa `$ in (44), we get a conservation low for the energy,
$$k_0k_0^{}\kappa =0.$$
(51)
Doing integration over $`\phi ,`$ we meet integrals of the form
$`J`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _0^{2\pi }}\mathrm{exp}[i(ll^{})\phi i\kappa r\mathrm{sin}\theta \mathrm{cos}(\phi \phi ^{})]𝑑\phi `$ (52)
$`=`$ $`J_{l^{}l}\left(2\sqrt{q\rho }\right)\mathrm{exp}\left[i\left(ll^{}\right)\left(\phi ^{}+{\displaystyle \frac{\pi }{2}}\right)\right],`$ (53)
where
$$q=\frac{\kappa ^2\mathrm{sin}^2\theta }{2\gamma }.$$
(54)
To make sure that (53) is correct, one can use an integral representation for Bessel functions (, 8.411.1). Integrating over $`\rho ,`$ we meet two integrals containing the Laguerre functions. These integrals can be done exactly as well,
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}I_{\alpha +m,m}(x)I_{\beta +n,n}(x)J_{\alpha +\beta }(2\sqrt{qx})𝑑x=(1)^{n+m}I_{\beta +n,m}(q)I_{\alpha +m,n}(q),`$ (55)
$`0n,mZ,\mathrm{}(\alpha +\beta +1)>0;`$ (58)
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}I_{\alpha +m,m}(x)I_{\beta +n,n}(x)J_{\alpha \beta }(2\sqrt{qx})𝑑x=(1)^{n+m}I_{n,m}(q)I_{\alpha +m,\beta +n}(q),`$
$`\mathrm{\hspace{0.33em}0}n,mZ,\mathrm{}(\alpha +1)>0.`$
Similar integrals can be encountered in (, 7.422.2), however there the calculation was in error.
Spinning particle case can be analyzed in the same manner. Thus, as in the conventional SR theory, matrix elements of electron transitions in the magnetic-solenoid field can be calculated exactly.
## IV Analysis of radiation frequencies
The relations (50) and (51) together with ones (25) or (32) define the frequency of the radiation $`\kappa `$ as a function of initial and final quantum numbers and of the angle $`\theta `$. Due to the axial symmetry of the problem, $`\kappa `$ does not depend on $`\phi ^{}`$. Similar to the conventional SR theory, we introduce a number $`\nu `$ of emitted harmonic as
$$\nu =nn^{}.$$
(59)
For $`\mu =0,`$ the frequency $`\kappa `$ is a function of the principle quantum number $`n`$, of $`\nu ,`$ and of the angle $`\theta `$ (see ). This frequency does not depend on the azimuthal quantum numbers $`l,l^{}`$. For $`\mu 0,`$ this degeneracy is partially lifted. In such a case, the frequency $`\kappa `$ depends on the type of initial and final states. Namely, it depends on the quantum numbers $`j,j^{}`$ in accordance with Eq. (51). Thus $`\kappa =\kappa _{jj^{}}`$. Introducing an effective number $`\overline{\nu }=\overline{\nu }_{jj^{}}`$ of emitted harmonic as
$$\overline{\nu }=\overline{n}\overline{n}^{}=\nu +\mu (j^{}j)=\{\begin{array}{c}\nu ,j=j^{}\hfill \\ \nu +\mu ,j=1,j^{}=2,\overline{\nu }>0\hfill \\ \nu \mu ,j=2,j^{}=1\hfill \end{array},$$
(60)
one can easily get for spinless particle case
$$\kappa _{jj^{}}=\frac{k_{0j}}{\mathrm{sin}^2\theta }\left(1\sqrt{1\beta _j^2\frac{2\overline{\nu }}{2\overline{n}+1}\mathrm{sin}^2\theta }\right),$$
(61)
where
$$\beta _j^2=1\left(\frac{m}{k_{0j}}\right)^2=1\left(\frac{m_0c^2}{E_j}\right)^2.$$
(62)
Similar formula takes place for spinning particle case,
$$\kappa _{jj^{}}=\frac{k_{0j}}{\mathrm{sin}^2\theta }\left(1\sqrt{1\beta _j^2\frac{\overline{\nu }}{\overline{n}}\mathrm{sin}^2\theta }\right).$$
(63)
The expressions (61) and (63) are obtained for initial states with $`k_3=0`$ (next we use the same supposition). Both expressions can be written in the form
$$\kappa _{jj^{}}=\frac{2\gamma \overline{\nu }}{k_{0j}+\sqrt{k_{0j}^22\gamma \overline{\nu }\mathrm{sin}^2\theta }}.$$
(64)
We can also get the following formulas
$`\kappa ={\displaystyle \frac{\gamma }{k_{0j}}}(\overline{\nu }+q),\mathrm{sin}\theta =\sqrt{{\displaystyle \frac{2}{\gamma }}}k_{0j}{\displaystyle \frac{\sqrt{q}}{\overline{\nu }+q}},`$ (65)
$`\sqrt{k_{0j}^22\gamma \overline{\nu }\mathrm{sin}^2\theta }=k_{0j}{\displaystyle \frac{\overline{\nu }q}{\overline{\nu }+q}},`$ (66)
where the quantity $`q`$ was defined by Eq. (54).
Thus, for $`\mu 0,`$ there appear two spectral series: one results from transitions without any change of the quantum number $`j`$ and another one results from transitions with the change of $`j`$. In the former case $`\overline{\nu }=\nu `$ whereas in the latter case $`\overline{\nu }=\nu \pm \mu `$ (effective numbers of emitted harmonics are not integer anymore). Whenever $`\nu >0`$ and $`n`$ are fixed, we gets an inequality
$$\kappa _{21}<\kappa _{11}<\kappa _{22}<\kappa _{12},$$
(67)
which becomes the equality for $`\mu =0`$. The difference between the frequencies $`\kappa _{11}`$ and $`\kappa _{22}`$ can be easily estimated for laboratory magnetic fields $`HH_0`$, where $`H_0=m_0^2c^3/e\mathrm{}4,41\times 10^{13}\mathrm{gauss}`$ is a critical field. This difference is proportional to $`\mu `$,
$$\kappa _{22}\kappa _{11}\mu \kappa _{22}\delta ,\delta =\frac{\gamma }{k_0^2}=\frac{\gamma }{m^2}\frac{m^2}{k_0^2}=\frac{H}{H_0}\left(\frac{m_0c^2}{E_1}\right)^2.$$
(68)
One can see that $`\delta <10^9`$ for not very high electron energies and for typical (those which are realized in accelerators) magnetic fields $`H10^4\mathrm{gauss}`$ . In such a case, the frequency difference reads
$$\kappa _{12}\kappa _{21}2\kappa _{12}\frac{\mu }{\nu }=2\frac{\omega }{c}\mu ,\omega =\frac{\left|ecH\right|}{E_j},$$
(69)
where $`\omega `$ is the synchrotron frequency. Thus, this difference becomes quite noticeable for harmonics with small numbers.
For $`\mu 0,`$ there exist a radiation of a harmonic $`\nu =0`$ due to $`j=1`$ $`j^{}=2`$ transitions. For magnetic fields $`HH_0`$ the frequency of such a harmonic is
$$\omega _{12}=\mu \omega =\mu \frac{ecH}{E_1}.$$
(70)
$`j=2`$ $`j^{}=1`$ transitions cause a radiation of a harmonic $`\nu =1.`$ For the above magnetic fields the corresponding frequency reads
$$\omega _{21}=(1\mu )\omega .$$
(71)
Such a radiation does not exist in pure magnetic field. Both frequencies $`\omega _{12},\omega _{12}`$ are less than $`\omega `$ ($`\omega `$ is the least radiated frequency in pure magnetic field).
## V Radiation of spinless particle
### A Exact expression for radiation intensity
As was discussed above, transition matrix elements that define one-photon radiation in the magnetic-solenoid field can be calculated exactly. For spinless particle case, the differential (with respect to the polarization) radiation intensity has the form:
$$W_j=W_0\frac{H}{H_0}\left(\frac{\gamma }{k_{0j}^2}\right)^2\underset{\nu ,j^{}}{}\frac{1}{2\pi }_0^{2\pi }𝑑\phi ^{}\frac{1}{2}_0^\pi 𝑑\theta \mathrm{sin}\theta \frac{(\overline{\nu }+q)^3}{\overline{\nu }q}Q_{jj^{}}|F_{jj^{}}|^2,$$
(72)
where
$`W_0={\displaystyle \frac{e^2m_0^2c^3}{\mathrm{}^2}},F_{jj^{}}=2l_2\sqrt{q}I_{jj^{}}^{}(q)+l_3\mathrm{cot}\theta \sqrt{{\displaystyle \frac{2k_{0j}^2}{\gamma }}}I_{jj^{}}(q),`$ (73)
$`I_{1j^{}}(q)=I_{\overline{n},\overline{n}^{}}(q),I_{2j^{}}(q)=I_{\overline{n}^{},\overline{n}}(q).`$ (74)
The quantities $`F_{jj^{}}`$ do not depend on the azimuthal quantum number $`l.`$ These quantities are completely defined by the quantum numbers $`\overline{n},\overline{\nu }`$,$`j,j^{}`$ and by the polarization of the radiation. The polarization is characterized by quantities $`l_2`$ and $`l_3`$ (see ). For $`l_2=1,l_3=0,`$ we get so called $`\sigma `$-component of the linear polarization; for $`l_2=0,l_3=1`$ we get so called $`\pi `$-component of the linear polarization; for $`l_2=\pm l_3=1/\sqrt{2},`$ we get right (left) circular polarization, and, finally, for $`l_2^2=l_3^2=1,l_2l_3=0,`$ we get total intensity of non-polarized radiation. The quantities $`Q_{jj^{}}`$ depend on initial quantum number $`l`$ only,
$$Q_{1j^{}}=\underset{l^{}}{}I_{\overline{n}^{}\overline{l}^{},\overline{n}\overline{l}}^2(q),Q_{2j^{}}=\underset{l^{}}{}I_{\overline{n}\overline{l},\overline{n}^{}\overline{l}^{}}^2(q).$$
(75)
Limits of the summation over final quantum numbers $`l^{}`$ depend on $`j^{}`$. Namely, $`0l^{}n\nu `$ whenever $`j^{}=1,`$ and $`\mathrm{}<l^{}1`$ whenever $`j^{}=2`$.
The integrand in (72) does not depend on $`\phi ^{}`$. Thus, the integration over this angle is trivial. The corresponding factor will be taken into account in following expressions.
Integrating over $`\theta ,`$ we get zero for total circular polarization. The reason is that dominant circular polarizations in the upper ($`0\theta \pi /2`$) and in the lower ($`\pi /2\theta \pi `$) half-planes have opposite signs and compensate each other exactly. If we are interested in linear polarization only, then we can always set $`l_2l_3=0`$ in Eq. (72).
Now we are going to fulfil summation in the intensity of the radiation over final azimuthal quantum numbers .
Consider first the case $`\mu =0.`$ Here the quantity $`|F|^2`$ in Eq. (72) does not depend on the type of the final state since the property (18) is valid in this case. Effective quantum numbers coincide with ordinary ones, $`\overline{n}=n,\overline{l}=l`$. Thus, taking into account (75), we get
$$\underset{j^{}}{}Q_{jj^{}}=\underset{k=0}{\overset{\mathrm{}}{}}I_{k,s}^2(q)=1,s=nl.$$
(76)
That is a well-known result in SR theory . In other words, the radiation intensity does not depend on the initial azimuthal quantum number $`l`$.
In magnetic-solenoid field with $`\mu 0,`$ the quantities $`Q_{jj^{}}`$ depend on $`l,n,\nu `$. Thus, the degeneracy with respect to $`l`$ is lifted completely. That may be interpreted as follows: According to Eq. (41), the quantum numbers $`l`$ define distances between classical trajectories and the solenoid. At the same time, $`l`$ defines the type of trajectories. Clearly, that the radiation intensity depends on the distances as well as on the type of states. For $`\mu =0,`$ the origin is not fixed anymore by the presence of the solenoid. Thus, the $`l`$ dependence of the radiation intensity dies out.
Let us return to Eq.(75), which define $`Q_{jj^{}}`$. Using properties of the Laguerre functions, we can get the following expression for a derivative of $`Q_{jj^{}}`$
$`{\displaystyle \frac{d}{dq}}Q_{jj^{}}(q)=(1)^{1+j+j^{}}\sqrt{{\displaystyle \frac{k+1}{q}}}[(2j)I_{k+1,s}(q)I_{k,s}(q)`$ (77)
$`+(j1)I_{s,k+1}(q)I_{s,k}(q)],k=\overline{n}^{}\mu ,s=\overline{n}\overline{l}.`$ (78)
Then, taking into account the behavior of $`Q_{jj^{}}`$ at $`q=0`$ and at $`q=\mathrm{}`$, we obtain
$`Q_{jj^{}}(q)=j^{}1+(1)^{j^{}1}[(2j){\displaystyle _q^{\mathrm{}}}\sqrt{{\displaystyle \frac{k+1}{y}}}I_{k+1,s}(y)I_{k,s}(y)dy`$ (79)
$`+(j1){\displaystyle _0^q}\sqrt{{\displaystyle \frac{k+1}{y}}}I_{s,k+1}(y)I_{s,k}(y)dy].`$ (80)
The result (76) follows from (80) as $`\mu 0`$.
### B Radiation in weak magnetic field approximation
Consider here the magnetic-solenoid field with $`H`$ obeying the condition $`HH_0`$ . Besides, we suppose that
$$2\gamma \overline{\nu }k_{0j}^2=2\frac{H}{H_0}\left(\frac{m_0c^2}{E_j}\right)^2\overline{\nu }1.$$
(81)
It is known that the only $`\overline{\nu }(E_j/m_0c^2)^3`$ harmonics are effectively emitted in the relativistic case. For such harmonics Eq. (81) results in
$$2\frac{H}{H_0}\frac{E_j}{m_0c^2}1.$$
(82)
It was demonstrated in that the condition (82) implies insignificance of quantum corrections in the relativistic case. That may be not true in the non-relativistic approximation since the only harmonic $`\overline{\nu }1`$ is emitted effectively and the condition (81) always holds for $`HH_0`$. Practically, the condition (82) always holds for real laboratory magnetic fields and electron energies. In the above suppositions, the quantity (54) reads
$$q=\frac{1}{2}\frac{H}{H_0}\left(\frac{m_0c^2}{E_j}\right)^2\overline{\nu }^2\mathrm{sin}^2\theta .$$
(83)
Suppose that the numbers of harmonics are not very big, then we can expect that
$$q1.$$
(84)
For the ultra-relativistic case $`\nu (E_j/m_0c^2)^3,`$ we find
$`q{\displaystyle \frac{H}{H_0}}\left({\displaystyle \frac{E_j}{m_0c^2}}\right)^4.`$
That means that in the latter case $`q`$ can be not small. Thus, namely the condition (84) defines the non-relativistic case in the weak magnetic field approximation. In other words, the condition (84) corresponds to CR in weak magnetic field approximation. Namely such a radiation is of concern to us in this Section. Below we suppose that (84) takes place.
In the case under consideration, we can present the radiation intensity in the following form
$`W_j=W_j^{\mathrm{cl}}\overline{W}_j,W_j^{\mathrm{cl}}={\displaystyle \frac{2}{3}}{\displaystyle \frac{e^4H^2\beta _j^2(1\beta _j^2)}{m_0^2c^3}},`$ (85)
$`\overline{W}_j={\displaystyle \underset{j^{}}{}}\overline{W}_{jj^{}},\overline{W}_{jj^{}}={\displaystyle \frac{3}{4(2\overline{n}+1)}}{\displaystyle _0^p}\mathrm{sin}\theta d\theta S_{jj^{}}^2{\displaystyle \underset{\nu }{}}\overline{\nu }^2R_{jj^{}},`$ (86)
$`S_{11}=S_{22}=S_{12}=l_2+l_3\mathrm{cos}\theta =S,S_{21}=l_2l_3\mathrm{cos}\theta =\overline{S}.`$ (87)
The quantity $`W_j^{\mathrm{cl}}`$ is the radiation intensity of a first harmonic in the semiclassical approximation (see ). The radiation polarization is characterized by the factor $`S_{jj^{}}`$. In particular, one can see that for transitions $`j=2j^{}=1`$ the sign of the radiation circular polarization is opposite to the sign of the circular polarization for all other transitions. That observation can be useful to identify the radiation related to $`j=2j^{}=1`$ transitions. The quantities $`R_{jj^{}}`$ can be calculated in lowest order of $`q`$ using exact expressions (72), (74), and (80).
1. For transitions $`j=1j^{}=1,`$ we find :
$`R_{11}={\displaystyle \frac{\mathrm{\Gamma }(n+\mu +1)q^{\nu 1}}{\mathrm{\Gamma }(n+\mu +1\nu )\mathrm{\Gamma }^2(\nu )}},\mathrm{\hspace{0.33em}\hspace{0.33em}1}\nu l,`$ (88)
$`R_{11}={\displaystyle \frac{\mathrm{\Gamma }(nl+1)\mathrm{\Gamma }(n+\mu +1)q^{2\nu l1}}{\mathrm{\Gamma }(n\nu +1)\mathrm{\Gamma }(n\nu +\mu +1\mathrm{\Gamma }^2(\nu l+1))\mathrm{\Gamma }^2(\nu )}},l\nu n.`$ (89)
One can see that the only harmonic $`\nu =1`$ is effectively emitted in these transitions. At the same time, the radiation intensity does not depend on $`l`$ whenever $`1ln`$. The quantity $`\overline{W}_{11}`$ can be easily calculated,
$$\overline{W}_{11}=\frac{2(n+\mu )}{2(n+\mu )+1}\overline{S^2},\overline{S^2}=\frac{3}{4}l_2^2+\frac{1}{4}l_3^2.$$
(90)
Thus, the radiation is polarized similarly to the conventional ($`\mu =0`$) SR case . The quantity $`\overline{W}_{11}`$ increases as $`n\mathrm{}`$, in particular, $`lim_n\mathrm{}\overline{W}_{11}=\overline{S^2}`$ .
For initial states with $`l=0,`$ transition probabilities are of order $`q`$. Let $`\nu =1`$, then we find for such transitions
$$\overline{W}_{11}=\frac{H}{H_0}\left(\frac{m_0c^2}{E_1}\right)^2\frac{3n(n+\mu )}{5(2n+2\mu +1)}\left(\frac{5}{6}l_2^2+\frac{1}{6}l_3^2\right).$$
(91)
Here the linear polarization of the radiation is greater than in (90). However, these transitions contribute insignificantly to the radiation compared to all other transitions.
2. For transitions $`j=2j^{}=2,`$ we find
$$R_{22}=\frac{\mathrm{\Gamma }(n+1)q^{\nu 1}}{\mathrm{\Gamma }(n+1\nu )\mathrm{\Gamma }^2(\nu )}.$$
(92)
As before, we see that the only first harmonic is effectively emitted. For this harmonic $`R_{22}=n`$ and
$$\overline{W}_{22}=\frac{2n}{2n+1}\overline{S^2}.$$
(93)
(Eq. (93) follows from (90) as $`\mu 0`$.) The radiation intensity does not depend on $`l<0`$ in the least order of $`q`$.
3. For transitions $`j=2j^{}=1,`$ we find
$$R_{21}=\frac{\mathrm{\Gamma }(n+|l|+1\mu )\mathrm{\Gamma }(n\nu +1+\mu )\mathrm{\Gamma }^2(1+\nu \mu )\mu ^2(1\mu )^2q^{|l|1}}{\mathrm{\Gamma }(n+1\nu )\mathrm{\Gamma }(n+1)\mathrm{\Gamma }^2(|l|+\nu +1\mu )}f^2(\mu ).$$
(94)
We have introduced here a function $`f(\mu )`$, $`0\mu 1,`$
$`f(\mu )={\displaystyle \frac{\mathrm{sin}\mu \pi }{\mu (1\mu )\pi }},f(\mu )=f(1\mu ),f(0)=f(1)=1,`$ (95)
$`f_{\mathrm{max}}(\mu =1/2)=4/\pi >1,\mathrm{\hspace{0.33em}1}f(\mu )4/\pi ,`$ (96)
which differs insignificantly from the unit whenever $`\mu 0`$.
In the transitions under consideration, we meet a situation, which is completely different from the one considered before. Here the only transitions from states with $`l=1`$ really contribute to the radiation. That fact has a natural physical explanation: For $`l=1,j=2,`$ classical trajectories do not embrace the solenoid but pass maximally close to the latter. A transition to trajectories embracing the solenoid is more likely namely from such states. It is important to stress that no restrictions exist on numbers of emitted harmonics. For $`l=1,`$ we get
$$R_{21}=\frac{\mathrm{\Gamma }(n+2\mu )\mathrm{\Gamma }(n\nu +1+\mu )\mu ^2(1\mu )^2}{\mathrm{\Gamma }(n+1)\mathrm{\Gamma }(n\nu +1)(\nu +1\mu )^2}f^2(\mu ),$$
(97)
and
$`\overline{W}_{21}={\displaystyle \frac{2\mathrm{\Gamma }(n+2\mu )\mu ^2(1\mu )^2M^{21}f^2(\mu )}{(2n+1)\mathrm{\Gamma }(n+1)}}\overline{S^2},`$ (98)
$`M^{21}={\displaystyle \underset{\nu =1}{\overset{n}{}}}M_\nu ^{21},M_\nu ^{21}={\displaystyle \frac{\mathrm{\Gamma }(n\nu +1+\mu )}{\mathrm{\Gamma }(n\nu +1)}}\left({\displaystyle \frac{\nu \mu }{\nu \mu +1}}\right)^2.`$ (99)
For example, for $`n=1`$ we obtain
$$\overline{W}_{21}(n=1)=\frac{2\mu ^2(1\mu )^4}{3(2\mu )}f(\mu )\overline{S^2}.$$
(100)
Expressions for big $`n`$ can be calculated approximately. Let us demonstrate how one can get an estimation for a typical sum.
The sum can be written as
$`{\displaystyle \underset{\nu =0}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n+2\mu \nu )}{\mathrm{\Gamma }(n\nu +1)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2={\displaystyle \frac{\mathrm{\Gamma }(n+2\mu )\mu ^2}{\mathrm{\Gamma }(n+1)(1\mu )^2}}`$ (101)
$`+{\displaystyle \frac{\mathrm{\Gamma }(n+1\mu )n(1+\mu )^2}{\mathrm{\Gamma }(n+1)\mu ^2}}+{\displaystyle \frac{\mathrm{\Gamma }(n\mu )n(n1)}{\mathrm{\Gamma }(n+1)}}\left({\displaystyle \frac{2+\mu }{1+\mu }}\right)^2`$ (102)
$`+{\displaystyle \underset{\nu =3}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n+2\mu \nu )}{\mathrm{\Gamma }(n\nu +1)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2.`$ (103)
For $`\nu =3`$ we have an inequality
$$1<\left(\frac{\nu +\mu }{\nu +\mu 1}\right)^2<\left(\frac{3+\mu }{2+\mu }\right)^2,$$
(104)
which allows us to write a relation
$$\underset{\nu =3}{\overset{n}{}}\frac{\mathrm{\Gamma }(n+2\mu \nu )}{\mathrm{\Gamma }(n\nu +1)}\left(\frac{\nu +\mu }{\nu +\mu 1}\right)^2=\delta \underset{\nu =3}{\overset{n}{}}\frac{\mathrm{\Gamma }(n+2\mu \nu )}{\mathrm{\Gamma }(n\nu +1)},$$
(105)
where $`\delta `$ can be estimated as
$$1<\delta <\left(\frac{3+\mu }{2+\mu }\right)^2.$$
(106)
The latter sum can be calculated exactly using the following well-known relation
$$\underset{\nu =1}{\overset{n}{}}\frac{\mathrm{\Gamma }(n+1+\mu \nu )}{\mathrm{\Gamma }(n\nu +1)}=\frac{\mathrm{\Gamma }(n+1+\mu )}{\left(1+\mu \right)\mathrm{\Gamma }(n)}.$$
(107)
This estimation can be improved if we write separately four or more terms in (103).
In the same manner, we get the following expression for the radiation intensity
$$\overline{W}_{21}=\frac{2n}{2n+1}R_n\left(\mu \right)\mu ^2\left(1\mu \right)^2\left[\left(\frac{1\mu }{2\mu }\right)^2+\left(n1\right)\delta \right]\overline{S^2},\frac{1}{2}<\delta <1,$$
(108)
where
$`R_n\left(\mu \right)={\displaystyle \frac{\mathrm{\Gamma }(n+\mu )\mathrm{\Gamma }(n+2\mu )}{\mathrm{\Gamma }^2(n+1)}}f^2\left(\mu \right),R_0\left(\mu \right)={\displaystyle \frac{f\left(\mu \right)}{\mu }},`$ (109)
$`R_1\left(\mu \right)=\left(2\mu \right)f\left(\mu \right),R_2\left(\mu \right)={\displaystyle \frac{1}{4}}\left(1+\mu \right)\left(2\mu \right)\left(3\mu \right)f\left(\mu \right),\mathrm{},`$ (110)
$`{\displaystyle \frac{n+1}{n}}f^2\left(\mu \right)R_n\left(\mu \right)R_n\left(1\right)=1,\underset{n\mathrm{}}{lim}R_n\left(\mu \right)=f^2\left(\mu \right).`$ (111)
One can see that the quantities $`M^{21}`$ from (99) change slightly as $`\nu `$ changes. Thus, at least a whole succession of first harmonics has equal probabilities of the radiation. In this approximation, such harmonics are not emitted for $`\mu =0`$ (they appear only in higher orders of $`q).`$ For big $`n,`$ one can see that $`M^{21}n.`$ The can serve as an additional argument in the favor of the above observation.
The case $`\nu =1`$ deserves to be considered especially. As was already remarked before, the corresponding radiation frequency (71) is less than the cyclotron one. It follows from (99) that
$$\overline{W}_{21}(\nu =1)=\frac{2n}{2n+1}R_n\left(\mu \right)\left[\frac{\mu (1\mu )^2}{2\mu }\right]^2\overline{S^2}.$$
(112)
Thus, in this case, the radiation intensity is approximately equal to the classical one multiplied by the factor
$$\left[\frac{\mu (1\mu )^2f\left(\mu \right)}{2\mu }\right]^2.$$
(113)
4. Finally, consider transitions $`j=1j^{}=2`$. In this case we get
$$R_{12}=\frac{\mathrm{\Gamma }(n\nu \mu +2)\mathrm{\Gamma }(n+\mu +1)(\nu +\mu )^2q^l}{\mathrm{\Gamma }(nl+1)\mathrm{\Gamma }(n\nu +1)\mathrm{\Gamma }^2(l\nu +2\mu )\mathrm{\Gamma }^2(\nu +1+\mu )}.$$
(114)
We see that the only transitions from states with $`l=0`$ contribute effectively to the radiation. Classical trajectories, which pass maximally close to the solenoid (embracing it), correspond to such initial states. Then, a possible physical interpretation is similar to the one given above. Thus, for $`l=0`$ we get
$`\overline{W}_{12}={\displaystyle \frac{2\mathrm{\Gamma }(n+1+\mu )\mu ^2(1\mu )^2M^{12}f^2(\mu )}{(2n+2\mu +1)\mathrm{\Gamma }(n+1)}}\overline{S^2},`$ (115)
$`M^{12}={\displaystyle \underset{\nu =0}{\overset{n}{}}}M_\nu ^{12},M_\nu ^{12}={\displaystyle \frac{\mathrm{\Gamma }(n\nu +2\mu )}{\mathrm{\Gamma }(n\nu +1)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2.`$ (116)
As before, we have here a radiation of $`\nu =0`$ harmonic (even for $`n=0`$ in the initial state). Such a radiation is forbidden for $`\mu =0`$. The frequency of the corresponding radiation is given by the expression (70). For $`\nu =0,`$ one finds
$$\overline{W}_{12}(\nu =0)=\mu ^4\frac{2(n+\mu )}{2(n+\mu )+1}R_n\left(\mu \right)\overline{S^2}.$$
(117)
In particular, for $`n=0`$ we find
$$\overline{W}_{12}(\nu =n=0)=\frac{2\mu ^4f(\mu )}{2\mu +1}\overline{S^2}.$$
(118)
Thus, the radiation intensity of such a harmonic is approximately equal to the classical intensity reduced by the factor $`\mu ^4f(\mu )`$.
In the transitions $`j=1j^{}=2`$, all the harmonics with different numbers $`\nu `$ contribute almost equally to the radiation intensity since $`M_\nu ^{12}`$ from (116) does not change significantly as $`\nu `$ varies. One can find the following estimation for $`\overline{W}_{12}`$ (taking into account the estimation (108) for $`\delta `$)
$$\overline{W}_{12}=\frac{2(n+\mu )}{2\left(n+\mu \right)+1}R_n\left(\mu \right)\left[\mu ^4+\frac{n\left(1\mu ^2\right)^2}{n+1\mu }+\frac{n\left(n1\right)\mu ^2\left(1\mu ^2\right)^2\delta }{n+1\mu }\right]\overline{S^2}.$$
(119)
The existence of the transitions under consideration may be treated as a manifestation of the AB effect. Indeed, in the absence of the solenoid field (more exactly for $`\mu =0)`$ and in the approximation under consideration, the only $`\nu =1`$ harmonic survives.
### C Peculiarities of radiation angular distribution
As it is known , in the relativistic case the intensity of the conventional SR is concentrated in the vicinity of the orbit plane within a small angular interval
$$\mathrm{\Delta }\theta \frac{m_0c^2}{E}=\sqrt{1\beta ^2}.$$
(120)
In such a case, the radiation intensity is maximal for harmonics with big numbers
$$\nu \left(\frac{E}{m_0c^2}\right)^3$$
(121)
Thus, it is widely believed that low number harmonics cannot practically be isolated against the background of intensive high frequency radiation.
However, there exist one exclusion from this rule. Indeed, we can easily see that in the conventional SR the intensity of all the harmonics with $`\nu 2`$ is exactly zero in the directions $`\theta =0,\pi `$ (along the magnetic field). Besides, the radiation intensity of the first harmonic ($`\nu =1`$) is maximal along the magnetic field for any particle energy. Moreover, the latter radiation has total circular polarization and, thus, can be easily identified.
The presence of the solenoid field modifies both the spectrum and angular distribution of SR. Consider, for example, the intensity of SR in the magnetic-solenoid field in the directions $`\theta =0,\pi `$ and within the infinitesimal solid angle $`d\mathrm{\Omega }=\mathrm{sin}\theta d\theta d\phi ^{}`$ . The expressions (72), (74), and (75) allow us to get the following exact result
$$4\pi \frac{dW_{jj^{}}}{d\mathrm{\Omega }}|_{\theta =0,\pi }=W^{\mathrm{cl}}G_{jj^{}}(l,n,\nu ;\mu ).$$
(122)
The quantity $`W^{\mathrm{cl}}`$ is defined by Eq. (85) and
$`G_{11}={\displaystyle \frac{3\left(n+\mu \right)\left(1\delta _{l,0}\right)\delta _{\nu ,1}}{2n+2\mu +1}},G_{22}={\displaystyle \frac{3n\delta _{\nu ,1}}{2n+1}},`$ (123)
$`G_{12}={\displaystyle \frac{3\left(n+\mu \right)R_n\left(\mu \right)\delta _{l,0}}{2n+2\mu +1}}[\mu ^4\delta _{\nu ,0}+{\displaystyle \frac{n\left(1\mu ^2\right)^2\delta _{\nu ,1}}{n\mu +1}}`$ (124)
$`+{\displaystyle \frac{\mu ^2\left(1\mu \right)^2\mathrm{\Gamma }\left(n+1\right)}{\mathrm{\Gamma }\left(n+2\mu \right)}}{\displaystyle \underset{\nu =2}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+2\mu \nu \right)}{\mathrm{\Gamma }\left(n+1\nu \right)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2],`$ (125)
$`G_{21}={\displaystyle \frac{3\mu ^2\left(1\mu \right)^2nR_n\left(\mu \right)\mathrm{\Gamma }\left(n\right)\delta _{l,1}}{\left(2n+1\right)\mathrm{\Gamma }\left(n+\mu \right)}}{\displaystyle \underset{\nu =1}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+1+\mu \nu \right)}{\mathrm{\Gamma }\left(n+1\nu \right)}}\left({\displaystyle \frac{\nu \mu }{\nu \mu +1}}\right)^2.`$ (126)
The function $`R_n\left(\mu \right)`$ is given by Eq. (111).
Let us briefly run through some of consequences of the above expression.
Transitions without a change of the type of the initial state (without a change of $`j`$) cause the only first harmonic ($`\nu =1`$) radiation in the directions $`\theta =0,\pi `$ whenever $`l0.`$ This fact does not depend on particle energies. One can see that the quantities $`G_{jj}`$ grow slightly and tend to finite constant values as $`n\mathrm{}.`$ Transitions from initial states with $`l=0`$ without a change of $`j`$ do not cause any radiation in $`\theta =0,\pi `$ directions.
Transitions with a change of the type of the initial state (with a change of $`j`$) cause a radiation in the directions $`\theta =0,\pi `$ solely for $`l=0,1`$(the solenoid is situated maximally close to a classical trajectory). In such cases all possible harmonics ($`0\nu n`$) are emitted with approximately equal intensities since the quantities $`G_{12},`$ $`G_{21}`$ grow proportionally to $`n.`$ For $`\mu =0,`$ the only first harmonic radiation survives.
Expressions (126) allow us to conclude that all the transitions cause totally circular polarized radiation in the directions $`\theta =0,\pi .`$ Moreover, as it follows from (87), the sign of the circular polarization for $`j=2j^{}=1`$ transitions is opposite to the one for all other transitions.
We believe that the peculiarities of the angular distribution of the radiation open up possibilities for experimental observation of superlow frequencies (70), (71) and of frequencies that are not multiple of the synchrotron one.
Note, that the expressions (126) (and the above mentioned consequences from them) were not known before even in the absence of the solenoid field (for $`\mu =0`$).
### D Semiclassical approximation
It was shown in the conventional SR theory that a semiclassical expansion of the radiation intensity can be done in terms of a small parameter $`\nu /n`$. Practically, to this end the formula
$$\underset{p\mathrm{}}{lim}I_{p+\alpha ,p+\beta }\left(\frac{x^2}{4p}\right)=J_{\alpha \beta }(x)$$
(127)
was used. Here $`J_\alpha (x)`$ are Bessel functions. It is natural to believe that for the case $`\mu 0`$ we can use the same parameter to perform the semiclassical expansion. Thus, we get a classical part of the intensity
$`W_j^{\mathrm{cl}}={\displaystyle \frac{e^4H^2(1\beta _j^2)}{m_0^2c^3}}{\displaystyle \underset{\nu ,j^{}}{}}{\displaystyle _0^p}\mathrm{sin}\theta d\theta \overline{\nu }^2Q_{jj^{}}^{\mathrm{cl}}|F^{\mathrm{cl}}|^2,F^{\mathrm{cl}}=l_2\beta _jJ_{jj^{}}^{}(z)+l_3\mathrm{cot}\theta J_{jj^{}}(z),`$ (128)
$`z=\overline{\nu }\beta _j\mathrm{sin}\theta ,J_{11}=J_{22}=J_{12}=J_{\overline{\nu }}(z),J_{21}=J_{\overline{\nu }}(z),`$ (130)
$`Q_{jj^{}}^{\mathrm{cl}}={\displaystyle \frac{1}{2}}+(1)^{j+j^{}}{\displaystyle _z^{\mathrm{}}}𝑑y(2j)J_{l\overline{\nu }+1}(y)J_{l\overline{\nu }}(y)+(j1)J_{|l|+\overline{\nu }}(y)J_{|l|+\overline{\nu }1}(y).`$
The components of $`Q_{jj^{}}^{\mathrm{cl}}`$ have the form
$`Q_{11}^{\mathrm{cl}}=\{\begin{array}{c}1_0^zJ_{l\nu +1}(y)J_{l\nu }(y)𝑑y,l\nu \hfill \\ _0^zJ_{\nu l1}(y)J_{\nu l}(y)𝑑y,l<\nu ,\hfill \end{array}Q_{22}^{\mathrm{cl}}=1{\displaystyle _0^z}J_{|l|+\nu }(y)J_{|l|+\nu l}(y)𝑑y,`$ (133)
$`Q_{12}^{\mathrm{cl}}=\{\begin{array}{c}_0^zJ_{l\nu +1\mu }(y)J_{l\nu \mu }(y)𝑑y,l\nu \hfill \\ \frac{1}{2}_z^{\mathrm{}}J_{l\nu +1\mu }(y)J_{l\nu \mu }(y)𝑑y,l<\nu ,\hfill \end{array}Q_{21}^{\mathrm{cl}}={\displaystyle _0^z}J_{|l|+\nu \mu }(y)J_{|l|+\nu l\mu }(y)𝑑y.`$ (136)
For $`\mu =0,`$ we see that $`|F^{\mathrm{cl}}|^2`$ does not depend on $`j^{}`$ and $`_j^{}Q_{jj^{}}^{\mathrm{cl}}=1.`$ Then the expression (130) presents the well-known classical SR differential intensity.
In the non-relativistic approximation, we get
$$W_j^{\mathrm{cl}}=W^{\mathrm{cl}}\overline{S^2}\underset{j^{}}{}\overline{W}_{jj^{}}^{\mathrm{cl}},\overline{W}_{jj}^{\mathrm{cl}}=1,\overline{W}_{12}^{\mathrm{cl}}=\mu ^4f^2(\mu ),\overline{W}_{21}^{\mathrm{cl}}=\left[\frac{\mu (1\mu )^2f(\mu )}{2\mu }\right]^2.$$
(137)
Here $`l=\nu =0`$ for $`j=1j^{}=2`$ transitions, and $`|l|=\nu =1`$ for $`j=2j^{}=1`$ transitions. Eqs. (137) follow from the exact quantum expressions (93), (113), and (117) as $`n\mathrm{}`$.
Semiclassical expressions (136) depend essentially on the initial azimuthal quantum number $`l`$ whenever $`\mu 0`$. Taking into account Eq. (40), we can express $`l`$ in terms of the pure classical quantity $`R^2R_0^2.`$ If $`R^2R_0^2`$ is fixed, we get $`|l|1/\mathrm{}\mathrm{}`$ as $`\mathrm{}0`$. Then, it follows from (136) that
$$Q_{jj^{}}^{\mathrm{cl}}=\delta _{jj^{}}.$$
(138)
Such a result seams to be natural. In classical theory of radiation trajectories of particles are fixed (there is no back reaction from emitted photons) and transitions with a change of initial states are not considered.
From Eqs. (136), we find the following relations
$$Q_{11}^{\mathrm{cl}}=\underset{k=\mathrm{}}{\overset{l\nu }{}}J_k^2(z),Q_{12}^{\mathrm{cl}}=\underset{k=l\nu +1}{\overset{\mathrm{}}{}}J_{k\mu }^2(z),Q_{22}^{\mathrm{cl}}=\underset{k=\mathrm{}}{\overset{|l|+\nu 1}{}}J_k^2(z),Q_{21}^{\mathrm{cl}}=\underset{k=|l|+\nu }{\overset{\mathrm{}}{}}J_{k\mu }^2(z).$$
(139)
It is clear that $`Q_{jj}^{\mathrm{cl}}`$ are monotonically increasing functions and $`Q_{jj^{}}^{\mathrm{cl}}`$ ($`jj^{}`$) are monotonically decreasing functions of $`\left|l\right|`$ such that
$$\underset{\left|l\right|\mathrm{}}{lim}Q_{jj^{}}^{\mathrm{cl}}=\delta _{_{jj^{}}}.$$
(140)
Thus, manifestations of the AB effect in SR are maximal for initial states with $`l=0,1.`$
All the angular distribution peculiarities, which were noted for the general case in the previous Section, take place in the approximation under consideration as well. Calculating the quantity $`4\pi \frac{dW_{jj^{}}}{d\mathrm{\Omega }}|_{\theta =0,\pi }`$ by the help of Eqs. (130) and (139), we can see that the representation (122) holds provided
$`G_{11}={\displaystyle \frac{3}{2}}(1\delta _{l,0})\delta _{\nu ,1},G_{22}={\displaystyle \frac{3}{2}}\delta _{\nu ,1},G_{12}={\displaystyle \frac{3}{2}}f^2\left(\mu \right)\delta _{l,0}[\mu ^4\delta _{\nu ,0}+(1\mu ^2)^2\delta _{\nu ,1}`$ (141)
$`+\mu ^2(1\mu )^2{\displaystyle \underset{\nu =2}{\overset{n}{}}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2],G_{21}={\displaystyle \frac{3}{2}}f^2(\mu \left)\mu ^2\right(1\mu \left)^2\delta _{l,1}{\displaystyle }_{\nu =1}^n\right({\displaystyle \frac{\nu \mu }{\nu \mu +1}})^2.`$ (142)
The latter quantities can be obtained from (126) in the limit $`n\mathrm{}.`$
## VI Radiation of spinning particle
### A Exact expression for radiation intensity
The analysis of radiation frequencies presented in Sect.IV was done for spinless particle case. However, all the qualitative results of this analysis remain valid for the spinning particle case. That can be seen from the corresponding exact formula (63).
In the spinning particle case, we get the following exact expression for the differential radiation intensity:
$$W_j=W_0\left(\frac{H}{H_0}\right)^2\epsilon _j_0^\pi 𝑑\theta \mathrm{sin}\theta \underset{\nu ,j^{},\zeta ^{}}{}Q_{jj^{}}\left[12p_j\left(\overline{\nu }+q\right)\right]^1\frac{(\overline{\nu }+q)^3}{\overline{\nu }q}|F_{jj^{}}|^2.$$
(143)
Here
$$\epsilon _j=1\beta _j^2=\left(\frac{m}{k_0}\right)^2=\left(\frac{m_0c^2}{E_j}\right)^2,p_j=\frac{1}{2}\frac{H}{H_0}\frac{\epsilon _j}{1+\sqrt{\epsilon _j}},$$
(144)
and the constant $`W_0`$ was introduced in (74). Remarkable that the quantities $`Q_{jj^{}}`$ are given by the same expressions (75) as in spinless particle case. Thus, all the previous conclusions related to these quantities are applied here as well.
In the spinning particle case, the quantities $`F_{jj^{}}`$ have the form
$`F_{jj^{}}=l_2F_{jj^{}}^{\left(2\right)}+l_3F_{jj^{}}^{\left(3\right)},F_{jj^{}}^{\left(2\right)}=\sqrt{{\displaystyle \frac{H}{H_0}}{\displaystyle \frac{\epsilon _j}{2q}}}[\delta _{\zeta ,\zeta ^{}}(1)^{j1}({\displaystyle \frac{1+\zeta }{2}}A_+^j+{\displaystyle \frac{1\zeta }{2}}A_{}^j)`$ (145)
$`+\delta _{\zeta ,\zeta ^{}}q\mathrm{cot}\theta ({\displaystyle \frac{1+\zeta }{2}}\chi _+^j+{\displaystyle \frac{1\zeta }{2}}\chi _{}^j)],F_{jj^{}}^{\left(3\right)}=\delta _{\zeta ,\zeta ^{}}(1)^{j1}\mathrm{cot}\theta `$ (146)
$`\times \left({\displaystyle \frac{1+\zeta }{2}}B_+^j+{\displaystyle \frac{1\zeta }{2}}B_{}^j\right)+\delta _{\zeta ,\zeta ^{}}p_j\left[\left(1\right)^{j1}{\displaystyle \frac{1+\zeta }{2}}C_+^j+{\displaystyle \frac{1\zeta }{2}}C_{}^j\right],`$ (147)
where
$`A_\pm ^j`$ $`=`$ $`2q\left[1p_j\left(\overline{\nu }+q\right)\right]{\displaystyle \frac{d\phi _\pm ^j}{dq}}\left[qp_j\left(\overline{\nu }+q\right)^2\right]\phi _\pm ^j,`$
$`B_\pm ^j`$ $`=`$ $`\left[1p_j\left(\overline{\nu }+q\right)\right]\phi _\pm ^j+2qp_j{\displaystyle \frac{d\phi _\pm ^j}{dq}},C_\pm ^j=\left[1\pm \sqrt{\epsilon _j}\left(\overline{\nu }+q\right)\right]\chi _\pm ^j+2q{\displaystyle \frac{d\chi _\pm ^j}{dq}},`$
$`\phi _+^1`$ $`=`$ $`I_{\overline{n}1,\overline{n}^{}1}(q),\phi _+^2=I_{\overline{n}^{}1,\overline{n}1}(q),\phi _{}^1=I_{\overline{n},\overline{n}^{}}(q),\phi _{}^2=I_{\overline{n}^{},\overline{n}}(q),`$
$`\chi _+^1`$ $`=`$ $`I_{\overline{n}1,\overline{n}^{}}(q),\chi _+^2=I_{\overline{n}^{},\overline{n}1}(q),\chi _{}^1=I_{\overline{n},\overline{n}^{}1}(q),\chi _{}^2=I_{\overline{n}^{}1,\overline{n}}(q),`$
and $`I_{n,n^{}}(x)`$ are the Laguerre functions. All the final quantum numbers are primed here.The quantities $`l_2`$ and $`l_3`$ characterize the radiation polarization. Contributions from transitions without ($`\delta _{\zeta ,\zeta ^{}}`$) and with ($`\delta _{\zeta ,\zeta ^{}}`$) spin-flip are separated.
The states with $`n=0`$ are a special case. For $`j=2,`$ there exist the only one (opposite to the magnetic field) spin orientation. Thus, all the transitions from any states with $`\zeta =1`$ to $`n=0`$, $`j=2`$ states do not cause a spin-flip ($`A_+^j=B_+^j=\phi _+^j=0`$), and all the transitions from any states with $`\zeta =1`$ to $`n=0`$, $`j=2`$ states do cause a spin-flip ($`C_{}^j=\chi _{}^j=0`$). One of such transitions is studied below. States with $`n=0`$, $`j=1`$ are singular at $`r=0.`$ However, they still can be normalized to a $`\delta `$function.
As in the spinless particle case, we can conclude that the radiation intensity depends on the mantissa $`\mu `$ only but not on the total solenoid magnetic flux $`\mathrm{\Phi }.`$
The radiation in question has not a preferential circular polarization. Total (integrated over all angles) intensities of the right and left circular polarizations are equal as in the scalar particle case. However, for transitions $`j=2j^{}=1`$ and $`j=2,\zeta =1j^{}=2,\zeta =1`$, the sign of the circular polarization is opposite to the one for all other transitions. Further, we are going to analyze the linear polarization only.
The angular distribution of the electron radiation intensity in the magnetic-solenoid field is quite similar to the one for spinless particle. The corresponding exact formula reads
$$4\pi \frac{dW_{jj^{}}}{d\mathrm{\Omega }}|_{\theta =0,\pi }=W^{\mathrm{cl}}\frac{3}{2}\left(\frac{1+\zeta }{2}G_{jj^{}}^++\frac{1\zeta }{2}G_{jj^{}}^{}\right),$$
(148)
where
$`G_{11}^+={\displaystyle \frac{\left(n+\mu 1\right)\left(1\delta _{l,0}\right)\delta _{\nu ,1}}{n+\mu }},G_{11}^{}=\left(1\delta _{l,0}\right)\delta _{\nu ,1},G_{22}^+={\displaystyle \frac{n1}{n}}\delta _{\nu ,1},G_{22}^{}=\delta _{\nu ,1},`$ (149)
$`G_{12}^+={\displaystyle \frac{nR_n\left(\mu \right)\delta _{l,0}}{n+\mu }}[\mu ^4\delta _{\nu ,0}+{\displaystyle \frac{(n1)\left(1\mu ^2\right)^2\delta _{\nu ,1}}{n\mu +1}}+{\displaystyle \frac{\mu ^2\left(1\mu \right)^2\mathrm{\Gamma }\left(n\right)}{\mathrm{\Gamma }\left(n+2\mu \right)}}`$ (150)
$`\times {\displaystyle \underset{\nu =2}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+2\mu \nu \right)}{\mathrm{\Gamma }\left(n\nu \right)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2],G_{12}^{}=R_n(\mu \left)\delta _{l,0}\right[\mu ^4\delta _{\nu ,0}+{\displaystyle \frac{n\left(1\mu ^2\right)^2\delta _{\nu ,1}}{n\mu +1}}`$ (151)
$`+{\displaystyle \frac{\mu ^2\left(1\mu \right)^2\mathrm{\Gamma }\left(n+1\right)}{\mathrm{\Gamma }\left(n+2\mu \right)}}{\displaystyle \underset{\nu =2}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+2\mu \nu \right)}{\mathrm{\Gamma }\left(n+1\nu \right)}}\left({\displaystyle \frac{\nu +\mu }{\nu +\mu 1}}\right)^2],`$ (152)
$`G_{21}^+={\displaystyle \frac{R_n\left(\mu \right)\mu ^2\left(1\mu \right)^2\mathrm{\Gamma }\left(n+1\right)\delta _{l,1}}{\mathrm{\Gamma }\left(n+\mu \right)}}{\displaystyle \underset{\nu =1}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+\mu \nu \right)}{\mathrm{\Gamma }\left(n+1\nu \right)}}\left({\displaystyle \frac{\nu \mu }{\nu \mu +1}}\right)^2,`$ (153)
$`G_{21}^{}={\displaystyle \frac{R_n\left(\mu \right)\mu ^2\left(1\mu \right)^2\mathrm{\Gamma }\left(n\right)\delta _{l,1}}{\mathrm{\Gamma }\left(n+\mu \right)}}{\displaystyle \underset{\nu =1}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }\left(n+1+\mu \nu \right)}{\mathrm{\Gamma }\left(n+1\nu \right)}}\left({\displaystyle \frac{\nu \mu }{\nu \mu +1}}\right)^2.`$ (154)
### B Radiation in weak magnetic field approximation
Here we suppose that the magnetic field is weak, i.e. $`HH_0`$ (more exactly $`q1),`$ and that initial quantum numbers are not very big (thus the particle remains non-relativistic). In this approximation, the main contributions to the radiation are due to transitions without a spin-flip. Below we present the only first terms in the $`H/H_0`$ decomposition for the radiation intensity.
First consider the radiation caused by transitions without a change of the state type ($`j=j^{}`$). Here the transitions $`\nu =1,l^{}=l1`$play the main role. For initial states with $`l0`$ the radiation intensity reads
$`W_j=W_j^{\mathrm{cl}}\{\delta _{\zeta ,\zeta ^{}}({\displaystyle \frac{1+\zeta }{2}}{\displaystyle \frac{\overline{n}1}{\overline{n}}}+{\displaystyle \frac{1+\zeta }{2}})S_0`$ (155)
$`+\delta _{\zeta ,\zeta ^{}}[{\displaystyle \frac{1+\zeta }{2}}{\displaystyle \frac{H}{H_0}}{\displaystyle \frac{S_1}{2\overline{n}}}+{\displaystyle \frac{1\zeta }{2}}\left({\displaystyle \frac{H}{H_0}}\right)^3{\displaystyle \frac{\overline{n}1}{70}}S_2]\}.`$ (156)
The quantity $`W_j^{\mathrm{cl}}`$ is the radiation intensity of the first harmonic in the semiclassical approximation (see (85)). The linear radiation polarization is characterized by the factors
$$S_0=\frac{3}{4}l_2^2+\frac{1}{4}l_3^2,S_1=\frac{1}{4}l_2^2+\frac{3}{4}l_3^2,S_2=\frac{1}{8}l_2^2+\frac{7}{8}l_3^2.$$
(157)
Whenever the initial state has the spin along the field ($`\zeta =1),`$ the ratio between the transitions with and without a spin-flip is of the order $`H/H_0.`$The same ratio is of the order$`\left(H/H_0\right)^3`$ for $`\zeta =1.`$ Thus, states with $`\zeta =1`$ are more stable than ones with $`\zeta =1.`$ That is the reason of the self-polarization effect in SR. The presence of the solenoid affects the only effective quantum numbers $`\overline{n},`$ the latter are not always integer, for example, $`\overline{n}=n+\mu `$ for $`j=1.`$ The radiation has a preferential linear polarization. For $`\zeta =1`$ initial states, transitions with and without a spin-flip cause radiation intensities of almost (with the interchange of $`l_2`$ and $`l_3)`$ the same form . For $`\zeta =1`$ initial states, transitions with a spin-flip cause almost (with the interchange of $`\sigma `$ and $`\pi `$ components) the same linear polarization of the radiation intensity as SR has in the relativistic case . For $`\mu =0,`$ the expression (156) (without the polarization specification) coincides with a corresponding expression presented in .
For $`l=0`$ in initial states, transitions without any change of $`j`$ are suppressed; these transitions contribute to the radiation intensity in higher orders of $`H/H_0`$ only. For example, for such transitions with a spin-flip, we find
$$W=W^{\mathrm{cl}}\frac{H}{H_0}\frac{3n\left(1\beta ^2\right)}{10}\left(\frac{1+\zeta }{2}\frac{n+\mu 1}{n+\mu }+\frac{1\zeta }{2}\right)\left(\frac{5}{6}l_2^2+\frac{1}{6}l_3^2\right).$$
(158)
We see that for the latter transitions, the polarization is distinctive and may serve to isolate such transitions.
Transitions with a change of the state type (with a change of $`j`$) are of particular interest from the AB effect point of view. Leading (with respect to $`H/H_0)`$ contributions correspond to $`j=1,l=0j^{}=2,l^{}=1`$ and to $`j=2,l=1j^{}=1,l^{}=0`$ transitions. In such transitions a whole set of successive harmonics is emitted, all these harmonics have approximately equal probabilities. The same situation was discovered by us in the spinless particle case.
As an example, we considered the radiation intensity for $`j=1j^{}=2`$ transitions without a spin-flip in detail. In such a case, this intensity has the form
$$W=W^{\mathrm{cl}}S_0M_{12}\delta _{\zeta ,\zeta ^{}},$$
(159)
where $`S_0`$ is defined by (157) and the quantity $`M_{12}`$ is a function of initial quantum numbers $`n`$ and $`\zeta ,`$
$`M_{12}=\mu ^2(1\mu )^2f^2(\mu ){\displaystyle \frac{\mathrm{\Gamma }(n+\mu )}{\mathrm{\Gamma }(n+1)}}`$ (161)
$`\times {\displaystyle \underset{\nu =0}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(n+2\mu \nu )}{\mathrm{\Gamma }(n+1\nu )}}\left({\displaystyle \frac{\mu +\nu }{\mu +\nu 1}}\right)^2({\displaystyle \frac{1+\zeta }{2}}{\displaystyle \frac{n\nu }{n+\mu }}+{\displaystyle \frac{1\zeta }{2}}).`$
( $`f(\mu )`$ was defined by (96).) In particular, here there is a possibility for $`\nu =0`$ transition with the emission of a superlow frequency (70). For the latter transition,
$$M_{12}(\nu =0)=\mu ^4R_n(\mu )\left(\frac{1+\zeta }{2}\frac{n}{n+\mu }+\frac{1\zeta }{2}\right),$$
(162)
where $`R_n\left(\mu \right)`$ is given by (111). Similar transition is possible even from $`n=0`$ states. Then
$$M_{12}(n=0,\zeta ,\mu )=\frac{1\zeta }{2}\mu ^3f(\mu ).$$
(163)
Taking into account that $`\beta ^2=2\mu H/H_0`$ for a state $`j=1,n=0,`$ we find from (159)
$$W(n=0)=\frac{4}{3}\frac{1\zeta }{2}W_0\left(\frac{H}{H_0}\right)^3\mu ^4f(\mu )S_0.$$
(164)
This results fitted well with Eq. (118) since in the spinless particle case
$$\beta ^2=\left(1+2\mu \right)H/H_0$$
(165)
for the state under consideration.
As before, for big $`n,`$ one can easily obtain estimations for emerged sums. Thus, for the radiation intensity caused by transitions with a change of $`j`$ and without any spin-flip, we get the following expression :
$$W=W^{\mathrm{cl}}R_n(\mu )\left(\frac{1+\zeta }{2}M^++\frac{1\zeta }{2}M^{}\right)S_0,$$
(166)
where
$`M_{12}^+={\displaystyle \frac{n}{n+\mu }}\left[\mu ^4+{\displaystyle \frac{\left(n1\right)\left(1\mu ^2\right)^2}{n+1\mu }}+{\displaystyle \frac{\left(n1\right)\left(n2\right)\mu ^2\left(1\mu \right)^2}{n+1\mu }}\delta _1\right],`$
$`M_{12}^{}=\mu ^4+{\displaystyle \frac{n\left(1\mu ^2\right)^2}{n+1\mu }}+{\displaystyle \frac{n\left(n1\right)\mu ^2\left(1\mu \right)^2}{n+1\mu }}\delta _1,M_{21}^+={\displaystyle \frac{\mu \left(1\mu \right)^2n}{n+\mu 1}}`$
$`\times [\mu \left({\displaystyle \frac{1\mu }{2\mu }}\right)^2+(n1)\delta _2],M_{21}^{}=\mu ^2(1\mu )^2[\left({\displaystyle \frac{1\mu }{2\mu }}\right)^2+(n1)\delta _2].`$
For $`\delta _k,k=1,2,`$ we have the estimation$`\mathrm{\hspace{0.33em}1}<\delta _1<2,\mathrm{\hspace{0.33em}1}/2<\delta _2<1.`$ The quantities $`M_{jj^{}}^\pm `$ are linearly increasing functions of $`n`$ whenever $`\mu 0`$ (the same property takes places in spinless particle case). For $`\mu =0,`$ we get another behavior
$$\underset{n\mathrm{}}{lim}M_{12}^\pm (\mu =0)=\delta _{\nu ,1},M_{21}^\pm (\mu =0,1)=0.$$
(167)
Consider transitions that cause non zero contributions to $`M_{21}^\pm `$ for $`\mu 0.`$ We know that the contribution of such transitions to the radiation intensity is of higher order of $`H/H_0`$ whenever $`\mu =0`$. Thus, only in the presence of the solenoid (with $`\mu 0)`$ a whole set of successive harmonics is emitted with approximately equal probabilities. The number of harmonics in the set is comparable with the number of the energy level.
For the radiation intensity caused by transitions with a change of the state type and with the spin-flip, we get the following results:
For $`j=1j^{}=2`$ transitions
$`W_{12}=W^{\mathrm{cl}}R_n(\mu )\left[{\displaystyle \frac{1+\zeta }{2}}{\displaystyle \frac{H}{H_0}}{\displaystyle \frac{S_1}{2\left(n+\mu \right)}}N_{12}^++{\displaystyle \frac{1\zeta }{2}}\left({\displaystyle \frac{H}{H_0}}\right)^3{\displaystyle \frac{2n}{35}}N_{12}^{}\right],`$ (168)
$`N_{12}^+=\mu ^6+{\displaystyle \frac{n(1+\mu )^2(1\mu ^2)^2}{n+1\mu }}+{\displaystyle \frac{n\left(n1\right)(n+2)^2\mu ^2(1\mu )^2}{n+1\mu }}\delta _1,`$ (169)
$`N_{12}^{}={\displaystyle \frac{\mu ^6}{\left(1+\mu \right)^2}}\left({\displaystyle \frac{\mu ^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right)+{\displaystyle \frac{\left(n1\right)(1\mu ^2)^2(1+\mu )^2}{\left(n+1\mu \right)\left(2+\mu \right)^2}}\left[{\displaystyle \frac{(1+\mu )^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right]`$ (170)
$`+{\displaystyle \frac{\left(n1\right)\left(n2\right)\mu ^2\left(1\mu \right)^2\delta _1}{n+1\mu }}\left[{\displaystyle \frac{(n+2)^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right].`$ (171)
For $`j=2j^{}=1`$ transitions
$`W_{21}=W^{\mathrm{cl}}R_n(\mu )\left[{\displaystyle \frac{1+\zeta }{2}}\left({\displaystyle \frac{H}{H_0}}\right)^3{\displaystyle \frac{2n}{35}}N_{21}^++{\displaystyle \frac{1\zeta }{2}}{\displaystyle \frac{H}{H_0}}{\displaystyle \frac{\mu (1\mu )^2S_1}{2(n+\mu 1)}}N_{21}^{}\right],`$ (172)
$`N_{21}^+={\displaystyle \frac{(1\mu )^6}{\left(2\mu \right)^2}}\left[{\displaystyle \frac{(1\mu )^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right]+{\displaystyle \frac{\left(n1\right)\mu ^2(2\mu )^4}{\left(n+1\mu \right)\left(3\mu \right)^2}}\left[{\displaystyle \frac{(2\mu )^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right]`$ (173)
$`+{\displaystyle \frac{\left(n1\right)\left(n2\right)\mu ^2\left(1\mu \right)^2\delta _1}{n+1\mu }}\left[{\displaystyle \frac{(n+2)^2}{8}}l_2^2+{\displaystyle \frac{7}{8}}l_3^2\right],`$ (174)
$`N_{21}^{}={\displaystyle \frac{\mu (1\mu )^4}{\left(2\mu \right)^2}}+\left(n1\right)(n+2)^2\delta _{2.}`$ (175)
Here the radiation intensity grows as $`n^4`$ whenever $`\mu 0`$, and the radiation polarization depends essentially on $`\mu `$ and $`n.`$
Of special note is the loss of spin $`\zeta =1`$ stability in transitions $`j=2j^{}=1`$. It follows from (175), for
$$\delta _3<\mu <1\delta _3,\delta _3=\frac{n}{3}\frac{H}{H_0},$$
(176)
that the spin $`\zeta =1`$ is more stable in the transitions under considerations. Let an initial state be of second ($`j=2)`$ type and the condition (176) holds, then the radiation creates a two-phase system of final electron states. Final electron states of second type have in the main negative spin orientation and final electron states of first type have in the main positive spin orientation. Thus, the presence of the solenoid field with $`\mu 0`$ plays a role of a depolarization factor in the above mentioned self-polarization effect.
### C Semiclassical approximation
Consider here the radiation intensity in the semiclassical approximation. From the previous discussion, we know that such an approximation corresponds to the condition $`v/n1.`$ Similar to the spinless particle case, we can approximate the Laguerre functions by the Bessel ones to get the following expression for the radiation intensity
$$W_j=W_0\left(\frac{H}{H_0}\right)^2(1\beta _j^2)\underset{\nu ,j^{}}{}_0^\pi Q_{jj^{}}^{\mathrm{cl}}|F_{jj^{}}^{\mathrm{cl}}|^2\mathrm{sin}\theta d\theta .$$
(177)
The quantities $`Q_{jj^{}}^{\mathrm{cl}}`$ are defined by Eq. (136) and $`F_{jj^{}}^{\mathrm{cl}}`$ have the form
$`F_{jj^{}}^{\mathrm{cl}}=\beta \delta _{\zeta ,\zeta ^{}}F_{jj^{}}^{\left(0\right)\mathrm{cl}}+\delta _{\zeta ,\zeta ^{}}{\displaystyle \frac{H}{H_0}}{\displaystyle \frac{1\beta ^2}{2}}\overline{\nu }F_{jj^{}}^{\left(1\right)\mathrm{cl}},F_{jj^{}}^{\left(0\right)\mathrm{cl}}=l_2I_{jj^{}}^{}(x)+l_3\mathrm{cos}\theta {\displaystyle \frac{I_{jj^{}}(x)}{\beta \mathrm{sin}\theta }},`$ (179)
$`F_{jj^{}}^{\left(1\right)\mathrm{cl}}=\left(\zeta \right)^jl_2\mathrm{cos}\theta \left[{\displaystyle \frac{I_{jj^{}}(x)}{\beta \mathrm{sin}\theta }}+\zeta I_{jj^{}}^{}(x)\right]l_3\left[{\displaystyle \frac{aI_{jj^{}}(x)}{\beta \mathrm{sin}\theta }}+\zeta I_{jj^{}}^{}(x)\right],`$
$`I_{11}(x)=I_{22}(x)=I_{12}(x)=J_{\overline{\nu }}(\overline{\nu }\beta \mathrm{sin}\theta ),I_{21}(x)=J_{\overline{\nu }}(\overline{\nu }\beta \mathrm{sin}\theta ),`$ (181)
$`x=\overline{\nu }\beta \mathrm{sin}\theta ,a=\mathrm{cos}^2\theta +\sqrt{1\beta ^2}\mathrm{sin}^2\theta .`$
In the non-relativistic approximation $`\beta ^2=2\overline{n}H/H_0`$, then results of the previous Section follow from (177).
It follows from (181) that the solenoid field with $`\mu 0`$ suppresses the electron self-polarization effect due to transitions $`j=2j^{}=1`$. This suppression can be considered as a manifestation of AB effect in SR. For $`\mu =0`$ such a manifestation disappears due to the property
$`J_\nu (x)=\left(1\right)^\nu J_\nu (x),`$
which takes place whenever $`\nu `$ are integer.
Similarly to the spinless particle case, the degeneracy of the radiation intensity with respect to the azimuthal quantum number is lifted here completely. That can be also considered as one of manifestations of AB effect in SR.
### D Electron transitions from zero energy levels with a change of state type
Consider here the radiation intensity caused by electron transitions from $`n=0`$ energy level with a change of the type of state (namely $`n=0,j=1j^{}=2`$ transitions). In this case a superlow frequency (70) is emitted. One can get an exact expression for the quantity $`Q_{12},`$
$$Q_{12}=\frac{q^{1\mu }\mathrm{exp}\left(q\right)\mathrm{\Phi }(1,2\mu ;q)}{\mathrm{\Gamma }\left(2\mu \right)},q=\mu \frac{1\sqrt{p}}{1+\sqrt{p}},p=1\alpha (1x^2),$$
(182)
where $`\mathrm{\Phi }(\alpha ,\gamma ;x)`$ is the confluent hypergeometric function. In the case under consideration, we can express $`\mathrm{\Phi }(\alpha ,\gamma ;x)`$ via the incomplete $`\mathrm{\Gamma }`$function and get the following expression
$$\mathrm{\Phi }(1,2\mu ;x)=\left(1\mu \right)x^{\mu 1}e^x_0^xe^xy^\mu 𝑑y,\mu <1.$$
(183)
For the transitions under consideration, the radiation intensity has the form
$`W=W_0{\displaystyle \frac{H}{H_0}}\left(\alpha \mu \right)^2f(\mu )G(\alpha ,\mu ),G(\alpha ,\mu )={\displaystyle _0^1}{\displaystyle \frac{\sqrt{p}+\sqrt{1\alpha }}{\sqrt{p}\left(1+\sqrt{p}\right)^3}}e^{2q}\mathrm{\Phi }(1,2\mu ;q)F\left(x\right)𝑑x,`$ (184)
$`F\left(x\right)=\delta _{\zeta ,\zeta ^{}}{\displaystyle \frac{1\zeta }{2}}\left[l_2^2+l_3^2\psi \left(x\right)\right]+\delta _{\zeta ,\zeta ^{}}{\displaystyle \frac{\alpha }{\left(1+\sqrt{1\alpha }\right)^2}}{\displaystyle \frac{1+\zeta }{2}}\left[l_2^2\psi \left(x\right)+l_3^2\right],`$ (186)
$`\alpha ={\displaystyle \frac{2\mu H}{H_0+2\mu H}},\psi \left(x\right)={\displaystyle \frac{x^2\left(1+\sqrt{1\alpha }\right)^2}{\left(\sqrt{p}+\sqrt{1\alpha }\right)^2}}`$
( $`f(\mu )`$ was defined in (96)). The function $`G(\alpha ,\mu )`$ depends on the magnetic field via the quantity $`\alpha ,`$
$$0<\alpha <1,\alpha 2\mu \frac{H}{H_0}\left(\frac{H}{H_0}1\right),\underset{H\mathrm{}}{lim}\alpha =1.$$
(187)
It is easy to see that
$$\alpha =1\left(m/k_0\right)^2=\beta ^2.$$
(188)
However, for such quantum states ($`n=0`$), we cannot use a classical interpretation for $`\beta .`$
It follows from (186) that transitions with and without spin-flip have almost (with the interchange of $`\sigma `$ and $`\pi `$ components) the same linear polarization of the radiation intensity.
Doing summation over photon polarization states, over final electron spin states, and averaging over initial spin states, we get total radiation intensity for a non-polarized electron
$$\overline{W}=2W_0\frac{H}{H_0}\alpha \mu ^2f(\mu )_0^1\frac{\sqrt{p}+\alpha 1}{\sqrt{p}\left(1+\sqrt{p}\right)^3}e^{2q}\mathrm{\Phi }(1,2\mu ;q)𝑑x.$$
(189)
In the weak magnetic field approximation ($`\alpha 1`$), we obtain from (186)
$$W=\frac{1}{3}W_0\frac{H}{H_0}(\alpha \mu )^2f(\mu )\left(\delta _{\zeta ,\zeta ^{}}\frac{1\zeta }{2}S_0+\delta _{\zeta ,\zeta ^{}}\frac{\alpha }{4}\frac{1+\zeta }{2}S_1\right).$$
(190)
Finally consider the case of superstrong magnetic fields ($`HH_0,\alpha =1)`$). Here $`\psi (x)=1`$ and the radiation intensity has the form
$`W`$ $`=`$ $`{\displaystyle \frac{1}{2}}\overline{W}\left(l_2^2+l_3^2\right)\left(\delta _{\zeta ,\zeta ^{}}{\displaystyle \frac{1\zeta }{2}}+\delta _{\zeta ,\zeta ^{}}{\displaystyle \frac{1+\zeta }{2}}\right),`$ (191)
$`\overline{W}`$ $`=`$ $`W_0{\displaystyle \frac{H}{H_0}}\mu ^2f(\mu )J(\mu ),J(\mu )={\displaystyle _0^1}(1+x)e^{2\mu x}\mathrm{\Phi }(1,2\mu ;\mu x)𝑑x.`$ (192)
$`J(\mu )`$ is a monotonically decreasing function of $`\mu .`$ In particular,$`J(0)=1,5;J(1)=2\frac{3}{e}0,896.`$ Thus, in the superstrong magnetic fields, transitions with and without spin-flip have equal probabilities, the radiation is completely depolarized, and the radiation intensity is linearly increasing function of the magnetic field.
## VII Summary
We have obtained exact solutions of Klein-Gordon and Dirac equations in the magnetic-solenoid field. Employing these solutions, we succeeded to calculate various characteristics of one-photon radiation in such a field. Namely, peculiarities of the radiation related to the presence of the AB solenoid are considered by us as manifestations of AB effect in CR and SR. Below we list the most important results obtained.
1. It is demonstrated that all the peculiarities of the radiation related to the presence of AB solenoid depend on the mantissa $`\mu `$ of the solenoid flux only. For the fluxes with $`\mu =0,`$ these peculiarities disappear.
2. The energy spectrum of charge particles in the magnetic-solenoid field differs essentially from the one in pure magnetic field. In particular, the degeneracy with respect to the azimuthal quantum number is partially lifted. Each magnetic field energy level splits in two ones in the magnetic-solenoid field. In turn, this complicates the radiation spectrum. In particular, the degeneracy of the radiation intensity with respect to the azimuthal quantum number is lifted completely.
3. New lines in the radiation spectrum appear, they do not have an analog in the pure magnetic field case. These lines consist of two series of harmonics (the latter are not multiple of the basic synchrotron frequency) and of two superlow frequency harmonics (their frequencies are less than the basic synchrotron frequency).
4. It is shown that the only one basic synchrotron harmonic and the new frequencies are irradiated along the magnetic field. We stress important peculiarities of the radiation along the magnetic field. The basic synchrotron harmonic has total circular polarization; the radiation intensity of superlow harmonics has maximum in the magnetic field direction; all the harmonics from the two above mentioned series have approximately equal radiation intensities. The latter property of the radiation is not typical for the conventional CR and SR. We believe that a considerable relative shift between new harmonics and the basic synchrotron one as well as the peculiarities of the angular distribution of the radiation intensity open up possibilities for experimental observation of AB effect in CR and SR.
5. It is discovered that the presence of the solenoid field can suppress the well-known in SR electron self-polarization effect.
Acknowledgement The authors (V.G.B, D.M.G, and A.L) are thankful to FAPESP for support. (D.M.G) thanks also CNPq for permanent support and (V.G.B.) thanks Russian Science Ministry Foundation and RFFI for partial support. |
warning/0001/nucl-th0001017.html | ar5iv | text | # The neutron star in the Relativistic Mean-Field Theory
## Introduction
The physics of compact objects like neutron stars offers an intriguing interplay between nuclear processes and astrophysical observables. Neutron stars exhibit conditions far from those encountered on earth. The determination of an equation of state (EoS) for dense matter is essential to calculations of neutron star properties.
This paper presents a basic model of neutron star matter including interactions among nucleons in the relativistic mean field approximation . Especially the Walecka model (QHD) and its nonlinear extensions have been quite successful and widely used for the description of hadronic matter and finite nuclei. Increasing interest in neutron matter at finite temperature has been observed recently in relation to the problems of hot neutron stars and of protoneutron stars and their evolutions in particular. Theories concerning protoneutron stars are being discussed in works by Prakash et. al. . Recently a detail calculation has been done with different models to study the properties of neutron stars . Glendenning has studied the properties of neutron star in the framework of nuclear relativistic field theory. In our calculations, we used the TM1 parameter set, which has a capability to reproduce the known results of finite nuclei as well as of normal nuclear matter. The TM1 model possess a distinctively stiffer EOS.
## The Relativistic Mean Field Theory
The fields of the model RMF for $`\sigma ,\omega `$ and $`\rho `$-mesons are denoted as $`\phi `$, $`\omega _\mu `$, $`\rho _\mu `$. The Lagrange density function for this model has the following form
$``$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \phi ^\mu \phi {\displaystyle \frac{1}{4}}R_{\mu \nu }^aR^{a\mu \nu }{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{1}{2}}M_\omega ^2\omega _\mu \omega ^\mu +{\displaystyle \frac{1}{2}}M_\rho ^2\rho _\mu ^a\rho ^{a\mu }`$ (2)
$``$ $`U(\phi )+{\displaystyle \frac{1}{4}}c_3(\omega _\mu \omega ^\mu )^2+i\overline{\psi }\gamma ^\mu D_\mu \psi \overline{\psi }(Mg_s\phi )\psi +`$
$`i{\displaystyle \underset{f=1}{\overset{2}{}}}\overline{L_f}\gamma ^\mu _\mu L_f{\displaystyle \underset{f=1}{\overset{2}{}}}g_f(\overline{L}_fHe_{Rf}+h.c.)`$
where
$$R_{\mu \nu }^a=_\mu \rho _\nu ^a_\nu \rho _\mu ^a+g\epsilon _{abc}\rho _\mu ^b\rho _\nu ^c$$
(3)
$$F_{\mu \nu }=_\mu \omega _\nu _\nu \omega _\mu $$
(4)
$$D_\mu =_\mu +\frac{1}{2}ig_\rho \rho _\mu ^a\sigma ^a+ig_\omega \omega _\mu $$
(5)
The potential is given by
$$U(\phi )=\frac{1}{2}m_s^2\phi ^2\frac{1}{3}g_2\phi ^3\frac{1}{4}g_3\phi ^4=\frac{1}{2}m_s^2\phi ^2+\frac{1}{3!}\kappa \phi ^3+\frac{1}{4!}\lambda \phi ^4$$
(6)
The fermion fields are composed of protons, neutrons and electrons, muons and neutrinos
$$\psi =\left(\begin{array}{c}\psi _p\hfill \\ \psi _n\hfill \end{array}\right),L_1=\left[\begin{array}{c}\nu _e\hfill \\ e^{}\hfill \end{array}\right]_L,L_2=\left[\begin{array}{c}\nu _\mu \hfill \\ \mu ^{}\hfill \end{array}\right]_L,e_{Rf}=(e_R^{},\mu _R^{}).$$
(7)
$`M`$ is the nucleon mass and $`m_s`$, $`M_\omega `$, $`M_\rho `$ are masses assigned to the mesons fields, $`g`$, $`g^{}`$ and $`g_s`$ are coupling constants. The Lagrangian function includes also the nonlinear term $`\frac{1}{4}c_3(\omega _\mu \omega ^\mu )^2`$ which affects remarkably the form of the equation of state. the Higgs field $`H`$ takes the form of
$$H=\frac{1}{\sqrt{2}}\left(\begin{array}{c}0\hfill \\ V\hfill \end{array}\right)$$
(8)
takes here the residual form.
The parameters used in NBL model are $`m_s=500`$ MeV with $`\kappa =800`$ MeV and $`\lambda =10`$.
The Euler equation for $`\mathrm{\Phi }_A=\{`$$`\phi `$, $`\omega _\mu `$, $`\rho _\mu `$,$`\psi \}`$ fields are
$$\mathrm{}\phi =m_s^2\phi +g_2\phi ^2+g_3\phi ^3g_s\overline{\psi }\psi $$
(9)
$$_\mu F^{\mu \nu }=M_\omega ^2\omega ^\nu +c_3(\omega _\mu \omega ^\mu )\omega ^\nu g_\omega J_B^\nu $$
(10)
where
$$J_B^\nu =\overline{\psi }\gamma ^\nu \psi $$
(11)
is the baryon current
$$D_\mu R^{\mu \nu ,a}=M_\rho ^2\rho ^{\nu ,a}g_\rho J_3^\nu $$
(12)
and
$$J_3^\nu =\frac{1}{2}\overline{\psi }\gamma ^\nu \sigma ^3\psi $$
(13)
is the isospin current. In the system we have conservation of baryon charge
$$Q_B=d^3xJ^0,$$
and the isospin charge
$$Q_3=d^3xJ_3^0.$$
The last is the Dirac equation
$$i\gamma ^\mu D_\mu \psi (Mg_s\phi )\psi =0.$$
(14)
The physical system is totally defined by the thermodynamic potential
$$\mathrm{\Omega }=kTlnTr(e^{\beta (H\mu Q_B\mu _3Q_3)})$$
(15)
where H is the Hamiltonian of the physical system
$$H=\underset{A}{}d^3x\{_0\mathrm{\Phi }_A\pi _\mathrm{\Phi }^A\}$$
(16)
and $`\pi ^A=\frac{}{(_0\mathrm{\Phi }_A)}`$ is a momentum connected to the field $`\mathrm{\Phi }_A`$. The fields $`\mathrm{\Phi }_A=\{`$$`\phi `$, $`\omega _\mu `$, $`\rho _\mu `$,$`\psi \}`$ denote all fields in the system. The average charges
$$\frac{\mathrm{\Omega }}{\mu }=<Q_B>,\frac{\mathrm{\Omega }}{\mu _3}=<Q_3>$$
(17)
can be obtained from the thermodynamic potential, of course they should be conserved. In this paper we shall use the effective potential approach build using the Bogolubov inequality
$$\mathrm{\Omega }\mathrm{\Omega }_1=\mathrm{\Omega }_0(m_B,m_F)+<HH_0>_0$$
(18)
$`\mathrm{\Omega }_0`$ is the thermodynamic potential of the trial system as effectively free quasiparticle system described by the Lagrange function
$`_0(m_B,m_F)`$ $`=`$ $`{\displaystyle \frac{1}{2}}_\mu \overline{\phi }^\mu \overline{\phi }{\displaystyle \frac{1}{2}}m_B^2\overline{\phi }^2{\displaystyle \frac{1}{4}}\overline{G}_{\mu \nu }^a\overline{G}^{a\mu \nu }{\displaystyle \frac{1}{4}}\overline{F}_{\omega .\mu \nu }\overline{F}_\omega ^{\mu \nu }`$ (19)
$`+{\displaystyle \frac{1}{2}}M_\omega ^2\overline{\omega }_\mu \overline{\omega }^\mu +{\displaystyle \frac{1}{2}}M_\rho ^2\overline{\rho }_\mu ^a\overline{\rho }^{a\mu }+\overline{\psi }(i\gamma ^\mu \overline{D}_\mu m_F)\psi `$
Similar to the general case
$$\overline{G}_{\mu \nu }^a=_\mu \overline{\rho }_\nu ^a_\nu \overline{\rho }_\mu ^a$$
and
$$\overline{F}_{\omega ,\mu \nu }=_\mu \overline{\omega }_\nu _\nu \overline{\omega }_\mu .$$
We decompose the $`\mathrm{\Phi }_A`$ field into two components, the effectively free quasiparticle field $`\stackrel{~}{\mathrm{\Phi }}_A`$ and the classical boson condensate $`\xi _A`$
$$\mathrm{\Phi }_A=\stackrel{~}{\mathrm{\Phi }}_A+\xi _A$$
(20)
In the case of the RMF model we have
$$\phi =\overline{\phi }+\sigma $$
(21)
$$\omega _\mu =\overline{\omega }_\mu +w_\mu ,w_\mu =\delta _{\mu ,0}w$$
(22)
$$\rho _\mu ^a=\overline{\rho }^a+r_\mu ^a,r_\mu ^a=\delta ^{a,3}\delta _{\mu ,0}r$$
(23)
The $`\xi _A=\{\sigma ,w,r\}`$ field will be treated as the variational parameters in the effective potential. Also the boson and fermion mass $`m_B,m_F`$ will be treated as as the variational parameters. The covariant derivative for the trial system is
$$\overline{D}_\mu =_\mu +\frac{1}{2}ig_\rho r_\mu ^a\sigma ^a+ig_\omega w_\mu $$
(24)
This introduce the homogenous fermion interaction with boson condensate $`w_\mu ,r_\mu ^a.`$ The fermion quasiparticle will obey the Dirac equation
$$(i\gamma ^\mu \overline{D}_\mu m_F)\psi =0$$
(25)
The constant condensate $`w,r`$ simply shift the chemical potential from $`\mu _i=\mu `$$`{}_{i}{}^{}{}_{}{}^{0}`$ (when $`w=r=0`$) to
$`\mu _n=\mu _n^0+{\displaystyle \frac{1}{2}}g_\rho rg_\omega w`$ (26)
$`\mu _p=\mu _p^0{\displaystyle \frac{1}{2}}g_\rho rg_\omega w`$ (27)
where $`\mu _n=\mu +\frac{1}{2}\mu _3`$ and $`\mu _p=\mu \frac{1}{2}\mu _3`$.
Neutrons, protons and electrons are in $`\beta `$-equilibrium which can be described as a relation among their chemical potentials
$$\mu _p+\mu _e=\mu _n$$
(28)
where $`\mu _p`$, $`\mu _n`$ and $`\mu _e`$ stand for proton, neutron and electron chemical potentials respectively. If the electron Fermi energy is high enough (greater then the muon mass) in the neutron star matter muons start to appear as a result of the following reaction
$$e^{}\mu ^{}+\nu _e+\overline{\nu _\mu }$$
(29)
The chemical equilibrium between muons and electrons can be described by the condition
$$\mu _\mu =\mu _e$$
(30)
When neutrinos are trapped inside the protoneutron star also the neutrino chemical potential should be included
$$\mu _p+\mu _e=\mu _n+\mu _{\nu _e}.$$
(31)
The density of the thermodynamic potential $`f_1=\mathrm{\Omega }_1/V`$ is equal to
$`f_1(m_B,m_F,\sigma ,w,r)=`$ (32)
$`f_0(m_B,m_F)+{\displaystyle \frac{1}{2}}<\overline{\phi }^2>_0(m_s^2m_B^2)+{\displaystyle \frac{1}{8}}\lambda <\overline{\phi }^2>_0^2+`$
$`{\displaystyle \frac{1}{2}}\kappa <\overline{\phi }^2>_0\sigma +{\displaystyle \frac{1}{2}}(m_s^2+{\displaystyle \frac{1}{2}}\lambda <\overline{\phi }^2>_0)\sigma ^2+`$
$`{\displaystyle \frac{1}{3!}}\kappa \sigma ^3+{\displaystyle \frac{1}{4!}}\lambda \sigma ^4+<\overline{\psi }\psi >_0(g\sigma m_F)+..,`$
$$f_0=f_B+f_F$$
(33)
where $`f_B`$ is the boson free energy and $`f_F`$ the fermion one. For boson field the free energy is
$$f_B=\frac{k_BT}{(2\pi )^3}d^3pln(1e^{\beta \omega (p)})$$
(34)
with $`\omega (p)=\sqrt{\text{p}^2+m_B^2}`$. For fermions we have 4 degree of freedom, 2 for spin 1/2 and 2 for particle - antiparticle distinguishing. For one fermion field the free energy is equal to
$$f_F=\underset{i=\{n,p\}}{}\frac{2k_BT}{(2\pi )^3}d^3p\{ln(1+e^{\beta (ϵ(p)\mu _i)})+ln(1+e^{\beta (ϵ(p)+\mu _i)})\}$$
(35)
now with $`ϵ(p)=\sqrt{\text{p}^2+m_F^2}`$. Variation
$$\frac{f_1}{m_B^2}=0,\frac{f_1}{m_F}=0$$
with respect to the trial system $`L_0`$ gives
$$m_B^2=m_s^2+2g_2\sigma +3g_3(\sigma ^2+<\overline{\phi }^2>_0),$$
(36)
$$m_F=M\delta =Mg_s\sigma .$$
(37)
In the local equilibrium inside the star the free energy reaches the minimum at $`\sigma `$.
The same result may be achieved calculating the averages of the equation of motions (9) for the effective system $`_0`$. In the mean field approximation the meson field operators are replaced by their expectation values. We also consider the isotropic system at rest. For the scalar field the equation (9) have is follows
$$(m_s^2+\frac{1}{2}\lambda <\overline{\phi }^2>_0)\sigma +\frac{1}{2}\kappa \sigma ^2+\frac{1}{6}\lambda \sigma ^3=g_s<\overline{\psi }\psi >_0\frac{1}{2}\kappa <\overline{\phi }^2>_0$$
(38)
If we shall notice that
$$\frac{f_B}{m_B^2}=\frac{1}{2}<\overline{\phi }^2>_0,$$
then it easy to obtain
$$<\overline{\phi }^2>_0=\frac{1}{2\pi ^2}\frac{dpp^2}{\sqrt{p^2+m_B^2}}\frac{1}{(exp(\beta \omega (p))1)}$$
This result means that equation (36) is highly nonlinear with respect to the meson mass $`m_B`$ (Fig. 2).
As $`<\overline{\psi }\psi >_0`$ (calculated with respect to $`_0`$ system) depends on the effective nucleon mass $`m_F`$ (or $`\sigma `$ ) the equation (38) is highly nonlinear also with respect to $`\sigma `$. In most of papers temperature dependence of the boson fields is neglected. Contrary to the fermion case the effective mass of boson $`m_B`$ growing with Fermi momentum. In the result the bosons temperature contributions may be neglected.
Calculation similar to the boson case, base on the relation
$$\frac{f_F}{m_F}=<\overline{\psi }\psi >_0$$
gives
$`<\overline{\psi }\psi >_0={\displaystyle \underset{i=\{n,p\}}{}}{\displaystyle \frac{m_F}{\pi ^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{p^2dp}{\sqrt{p^2+m_F^2}}}\{{\displaystyle \frac{1}{\mathrm{exp}(\beta (ϵ_p\mu _i))+1}}+`$ (39)
$`{\displaystyle \frac{1}{\mathrm{exp}(\beta (ϵ_p+\mu _i))+1}}\}.`$
The quantum average $`<\overline{\psi }\psi >_0`$ depends on the neutron and proton chemical potentials $`\mu _p`$, $`\mu _n`$. In the result the effective nucleon effective mass $`m_F`$ also will be dependent he neutron and proton chemical potentials. In the result the solution $`\sigma `$ of the equation (9) also will be dependent on $`\xi `$. The same situation will consider other fields. This model is the simple example of the relativistic mean filed theory .
The effective mass $`m_F`$ (or $`\delta =m_F/M`$ ) dependence on the dimensionless Fermi momentum $`x_F`$ is presented on the Fig.3.
The binding energy
$$E_0=ϵ(x_F,T)/Q_BM$$
(40)
for nucleon symmetric phase is presented on Fig. 2 (see Table 2).
We see that binding energy strongly depends on temperature and above $`T>15`$ MeV is positive.
The temperature dependence of the binding energy in presented on the Fig 4. To calculate the properties of the neutron star we need the energy-momentum tensor. In case of the fermions field it is more convenient to use the reper filed $`e_\mu ^a`$ defined as follows $`g_{\mu \nu }=e_\mu ^ae_\nu ^b\eta _{ab}`$ where $`\eta _{ab}`$ is the flat Minkowski space-time matrix. The general definition
$$T_{\mu \nu }=2\frac{L_B}{g^{\mu \nu }}+e_\mu ^a\frac{L_F}{e^{a\nu }}g_{\mu \nu }L$$
(41)
allows us to calculate the density of energy and pressure. The total Lagrange function $`L=L_B+L_F`$ is divided into boson and fermion part. To calculate the density of energy and pressure we shall average the energy-momentum tensor $`T_{\mu \nu }`$ with respect to the quasi equilibrium configuration defined by the trial system $`L_0`$. We define the density of energy and pressure by the energy - momentum tensor
$$<T_{\mu \nu }>=(P+ϵ)u_\mu u_\nu Pg_{\mu \nu }$$
(42)
where $`u_\mu `$ is a unite vector ($`u_\mu u^\mu =1`$). So, the calculations give
$$ϵ(x_F,T)=\rho c^2=\frac{1}{2}M_\omega ^2w^2\frac{1}{2}M_\rho ^2r^2+g_\omega Q_Bw+\frac{1}{2}g_\rho Q_3r\frac{1}{4}c_3w^4+U(\sigma )+ϵ_F$$
(43)
$$P(x_F,T)=\frac{1}{2}M_\omega ^2w^2+\frac{1}{2}M_\rho ^2r^2+\frac{1}{4}c_3w^4U(\sigma )+P_F$$
(44)
where
$$ϵ_F=ϵ_0\chi (x_F,T)$$
(45)
$$P_F=P_0\varphi (x_F,T)$$
(46)
The fact that neutron mass depends on fermion concentration ( or neutron chemical potential $`\mu `$) now must be included into $`\chi (x_F,T)`$ and $`\varphi (x_F,T)`$,
$`\chi (x_F,T)={\displaystyle \frac{1}{^{\pi ^2}}}{\displaystyle _0^{\mathrm{}}}dzz^2\sqrt{z^2+\delta ^2(x_F)}\{{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{\delta ^2(x_F)+z^2}\mu ^{})/\tau )+1}}`$ (47)
$`+{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{\delta ^2(x_F)+z^2}+\mu ^{})/\tau )+1}},`$
$`\varphi (x_F,T)={\displaystyle \frac{1}{3\pi ^2}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{z^4dz}{\sqrt{z^2+\delta ^2(x_F)}}}\{{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{\delta ^2(x_F)+z^2}\mu ^{})/\tau )+1}}`$ (48)
$`+{\displaystyle \frac{1}{\mathrm{exp}((\sqrt{\delta ^2(x_F)+z^2}+\mu ^{})/\tau )+1}}\}`$
where $`\tau =(k_BT)/M`$,
$$\mu ^{}=\mu /M=\sqrt{\delta ^2(x_F)+x_F^2}$$
(49)
and
$$x_F=k_F/M$$
(50)
Similar to paper we have introduced (49,50) the dimensionless “Fermi” momentum even at finite temperature which exactly corresponds to the Fermi momentum at zero temperature. Both $`ϵ_F`$ and $`P_F`$ depend on the neutron chemical potential $`\mu `$ or Fermi momentum $`x_F`$. This parametric dependence on $`\mu `$ (or $`x_F`$) defines the equation of state. The various equations of state for different parameters sets is presented on fig.5.
The equation of state for the parameters set TM1 for temperature T=10 MeV is presented on fig.6. It is interesting to notice that even for $`Q_B=0`$ due to the presence of the thermal exited particle antiparticle pairs there is finite energy and pressure density.
It is interesting to notice that the meson field $`\phi `$ as the scalar field contributes to the pressure with negative sign while the vector meson fields ($`\omega ,\rho `$) with the positive one.
## The neutron star
In this paper we present numerical results describing the structure of neutron star based on the relativistic mean field theory. It is possible to describe a static spherical star solving the OTV equation.
$$\frac{dP(r)}{dr}=\frac{G}{r^2}(\rho (r)+\frac{P(r)}{c^2})\frac{(m(r)+\frac{4\pi }{c^2}P(r)r^3)}{(1\frac{2Gm(r)}{c^2r})}$$
(51)
$$\frac{dm(r)}{dr}=4\pi r^2\rho (r)$$
(52)
Having solved the OTV equation the pressure $`p(r)`$, mass $`m(r)`$ and density $`\rho (r)`$ were obtained. To obtain the total radius $`R`$ of the star the fulfillment of the condition $`p(R)=0`$ is necessary. This allows to determine the total gravitational mass of the star $`M(R)`$.
Introducing the dimensionless variable $`\xi ,`$ which is connected with the star radius $`r`$ by the relation $`r=a\xi `$ enables to define the functions $`p(r)`$, $`\rho (r)`$ and $`m(r)`$
$$\rho (r)=\rho _0\chi (x(\xi ))$$
(53)
$$P(r)=P_0\phi (x(\xi ))$$
(54)
$$m(r)=M_{}v(\xi )$$
(55)
by $`\xi `$. If we define dimensionless functions.
$$\lambda =\frac{GM_{}\rho _c}{P_0a},\mu =3\frac{M_c}{M_{}},M_c=\frac{4}{3}\pi \rho _0a^3$$
(56)
are also need to achieve the OTV equation of the following form
$$\frac{d\phi }{d\xi }=\lambda (\chi (\xi )+\phi (\xi ))\frac{v(\xi )+\mu \phi (\xi )\xi ^3}{\xi ^2(1\frac{r_g}{a}\frac{v(\xi )}{\xi })}$$
(57)
$$\frac{dv}{d\xi }=\mu \chi (\xi )\xi ^2$$
(58)
with $`r_g`$ being the gravitational radius.The equations (57,58) are easy integrated numerically, For example, for the neutron star with the central density $`\rho _c=\mathrm{9\hspace{0.17em}10}^{14}g/cm^3`$ the star profile in the mean field approach is presented on the Fig.7.
It’s interesting that fermions contribution to density is lower then a half of total density. The biggest contributions come from the nucleons, the gauge boson field $`\omega _\mu `$ and scalar boson field $`\phi `$.
On the surface of star we have non zero value of density ($`\rho 0.25<\rho _0`$) and nucleon mass lower than vacuum nucleon mass ($`m_F\mathrm{\hspace{0.17em}762}MeV<M`$). Inside star with $`\rho _c=9\times 10^{14}g/cm^3`$ (maximal stable configuration for the TM1 parameters set) value of nucleon mass varies from $`300`$ MeV in centrum of star to $`762`$ MeV on the surface (Fig. 8). Fig. 10. shows the stellar masses as a function of the central density. The parameters of the maximum mass configuration are:
$$M_{max}=1.91M_{},R=12.84km.$$
(59)
This fact is easy to notice on the mass-radius diagram (Fig. 9).
When temperature is different from zero, a star is in generar bigger and more massive (see Figs. 9 and 10).
## Conclusion
The structure of static neutron stars can be determined by solving the Tolman-Oppenheimer-Volkoff equations. The equation of state will strong influence on the neutron star properties. The aim of this work has been to present the equation of state in the RMF for the nuclear matter with temperature different from zero. The hot neutron star in the simple L2 parameters set was examined in our previous work . The main objective of our work was to study the influence of the temperature on the main parameters of a neutron star. Neutron matter at finite temperature is of increasing interest in relation to the problems of hot neutron stars. In order to achieve the proper form of the equation of state the relativistic mean field approach was involved.
## Appendix: Nuclear matter properties
The nuclear symmetric matter may be described by the phenomenological equation of state
$$\epsilon (u)=\rho _0u(M+\frac{3}{5}\frac{\mathrm{}^2k_F^2}{2M}+\frac{1}{2}Au^2+\frac{B}{(\sigma +1)}u^\sigma )$$
(60)
where $`u=\rho /\rho _0`$ is a dimensionless density. For symmetric nuclear matter dimensionless $`\rho _0=\mathrm{2.5\hspace{0.17em}10}^{14}`$ $`gcm^3`$ $`=0.15nucleons/fm^3`$ $`=140`$ MeV$`fm^3`$. The parameter of $`\epsilon `$ (60) are
$$\sigma =\frac{K_0+2E_{F,0}}{3E_{F,0}9E_0},$$
(61)
$$B=(\frac{\sigma +1}{\sigma 1})[\frac{1}{3}E_{F,0}E_0],$$
(62)
$$A=E_0\frac{5}{3}E_{F,0}B.$$
(63)
$`E_{F,0}`$ is the nonrelativistic nucleon Fermi energy at the saturation point $`k_{F,0}`$ (see Table 2). For the TM1 parameter set (see Table 1) we have
$$\sigma =1.5485;A=\mathrm{1.58.52}MeV;B=107.689MeV;$$
In this approach the binding energy (see Fig.11 ) is
$$E_0(u)=E_{F,0}u^{\frac{2}{3}}+\frac{1}{2}Au+\frac{B}{(\sigma +1)}u^\sigma .$$
(64)
The pressure is
$$P=u^2\frac{d(\epsilon /Q_B)}{d\rho }|_{\rho _0}=\frac{2}{3}E_{F,0}u^{\frac{5}{3}}+\frac{1}{2}Au^2+\frac{B\sigma }{(\sigma +1)}u^{\sigma +1}.$$
(65)
The minimum of the binding energy determine the equilibrium Fermi momentum $`k_{F,0}`$, density $`\rho _0`$ and incompressibility factor
$$K=9\rho _0\frac{d^2(\epsilon /Q_B)}{d\rho ^2}|_{\rho _0}.$$
(66)
The empirical value of $`K`$ is $`210\pm 30MeV`$ . The incompressibility factor is equal to
$$K(u)=9\frac{d(P/\rho )}{du}|_{\rho _0}=10E_{F,0}u^{\frac{2}{3}}+9(Au+B\sigma u^\sigma ).$$
(67)
The chemical potential defined for particle species $`i`$ is given by
$$\mu _i(u)=\frac{d(\epsilon )}{d\rho _i.}.$$
(68)
For example, for nucleon symmetric saturation point $`\mu =912.73MeV`$ for nucleons. The adiabatic sound speed
$$(\frac{v}{c})^2=\frac{dP}{d\epsilon }=\frac{K}{9\mu }$$
(69)
is presented on the Fig.12. At the saturation point we have the sound speed $`v=0.184c`$. For TM1 parameters, for $`\rho `$ above $`5\rho _0`$ the theory looses its sense. For other parameters sets when $`K`$ is grater the range of the theory validity is smaller.
Acknowledgment
The authors are thankful to F. Weber for helpful discussions. |
warning/0001/nucl-th0001035.html | ar5iv | text | # Final state interactions in 4He(e,e′p)3H at large proton energy
## 1 Introduction
Electron-nucleus scattering experiments in which a proton is detected in coincidence with the outgoing electron have long been recognized as a powerful tool to study both nuclear and nucleon dynamics (see, e.g., ref.). According to the Plane Wave Impulse Approximation (PWIA), which is expected to be valid at large momentum transfer, the nuclear $`(e,e^{}p)`$ cross section reduces to the incoherent sum of the cross sections off individual nucleons, whose distribution in momentum $`𝐤`$ and removal energy $`E`$ is dictated by the spectral function $`P(k,E)`$, the final state interactions (FSI) between the knocked out particle and the recoiling spectator system being negligible. As a consequence, the PWIA cross section of the process in which an electron of initial energy $`E_i`$ is scattered into the solid angle $`\mathrm{\Omega }_e`$ with energy $`E_f=E_i\omega `$, while a proton of kinetic energy $`T_p`$ is ejected into the solid angle $`\mathrm{\Omega }_p`$, takes the simple factorized form
$$\frac{d\sigma }{d\omega d\mathrm{\Omega }_ed\mathrm{\Omega }_pdT_p}=p(T_p+m)\stackrel{~}{\sigma }_{ep}P(p_m,E_m),$$
(1)
where $`m`$ denotes the nucleon mass, while the missing momentum $`𝐩_m`$ and missing energy $`E_m`$ are defined as
$$𝐩_m=𝐩𝐪$$
(2)
and
$$E_m=\omega T_pT_R.$$
(3)
In the above equations, $`𝐪`$ is the momentum transfer and $`T_R=p_R^2/M_{A1}`$, with $`𝐩_R=𝐩_m`$, is the kinetic energy of the recoiling spectator system of mass $`M_{A1}`$. The cross section $`\stackrel{~}{\sigma }_{ep}`$ of eq.(1) describes electron scattering off a bound nucleon of momentum $`𝐩_m`$ and removal energy $`E_m`$ .
In presence of nonnegligible FSI, the PWIA picture breaks down, and the missing momentum and energy cannot be readily interpreted as the initial momentum and removal energy of the outgoing nucleon. Therefore, a quantitative understanding of FSI is needed in order to extract the information on the nucleon spectral function from the measured $`(e,e^{}p)`$ cross section. A wealth of highly accurate theoretical calculations of FSI effects in $`(e,e^{}p)`$ reactions have been carried out within the Distorted Wave Impulse Approximation (DWIA), in which the interaction between the knocked out nucleon and the spectator system is described in terms of a complex optical potential (see, e.g., ref.). Using the results of these calculations it has been possible to obtain the spectral functions describing the single-particle states, predicted by the nuclear shell model, from the analysis of the available low missing energy data .
It has to be emphasized, however, that FSI should not only be regarded as a noise, to be removed from the measured cross sections. In fact, in many instances FSI produce a signal that carries relevant information on both the target structure and the dynamics of the scattering process. For example, it has been shown that FSI effets, which obviously depend upon the distribution in space of the spectator particles, are very sensitive to the presence of local fluctuations of the nuclear density produced by nucleon-nucleon (NN) correlations .
The analysis of FSI in $`(e,e^{}p)`$ processes may also provide information on NN scattering in the nuclear medium. At moderate proton energies ($``$ 100 MeV) nuclear structure effects, such as Pauli blocking and dispersive corrections, lead to significant changes in the NN scattering amplitude . While these effects are expected to become negligible for proton energies in the few GeV range, different effects, arising from nucleon structure, may become important in this kinematical regime, corresponding to high $`Q^2`$ ($`Q^2=q^2\omega ^2`$).
Perturbative Quantum Chromo-Dynamics (QCD) predicts that elastic scattering on a nucleon at high momentum transfer can only occur if the nucleon is found in the Fock state having the lowest number of constituents, so that the momentum can be most effectively shared among them. This state, being very compact, interacts weakly with the spectator particles and evolves to the standard proton configuration with a characteristic timescale that increases with the momentum transfer. According to this picture a proton, after absorbing a large momentum $`q`$, e.g. in an electron scattering process, can travel through the spectator system experiencing very little attenuation, i.e. exhibits color transparency (CT) . In the limit $`Q^2\mathrm{}`$ FSI effects in $`(e,e^{}p)`$ are expected to become vanishingly small.
The possible signatures of the occurrence of CT in coincidence $`(e,e^{}p)`$ and $`(p,2p)`$ processes have been recently studied within a theoretical many-body approach suitable for the calculation of semi-inclusive cross sections, involving a sum over the states of the undetected spectator system . The treatment of FSI of refs. is based on a generalization of Glauber theory of high energy proton scattering .
In this paper we extend the approach of refs. to the case of fully exclusive reactions, in which the final state of the recoiling nucleus is specified. Our treatment of the corresponding amplitude is presented in section II, where we discuss both the many-body aspects, related to the description of the nuclear initial and final states, and the structure of the scattering operator, modeling the FSI of the knocked out proton and the transition to the CT regime. The results obtained applying our approach to the case of a <sup>4</sup>He target, in which accurate numerical calculations are feasible, are given in section III, where FSI effects on different observables are discussed. Finally, the summary and conclusions are presented in section IV.
## 2 Formalism
### 2.1 $`\mathbf{(}𝒆\mathbf{,}𝒆^{\mathbf{}}𝒑\mathbf{)}`$ amplitude at high proton energy
We will focus on $`(e,e^{}p)`$ processes in which the recoiling (A-1)-particle system is left in a bound state $`|\phi _n`$. Neglecting many-body contributions to the electromagnetic current, the nuclear matrix element associated with the transition amplitude can be written
$$M_n(𝐩,𝐪)=\mathrm{\Psi }_{n𝐩}^{()}|\underset{𝐤}{}a_{𝐤+𝐪}^{}a_𝐤|\mathrm{\Psi }_0,$$
(4)
where $`a_{𝐤+𝐪}^{}`$ ($`a_𝐤`$) denotes the usual creation (anihilation) operator and the target ground state $`|\mathrm{\Psi }_0`$ satisfies the Schrödinger equation $`H_A|\mathrm{\Psi }_0=E_0|\mathrm{\Psi }_0`$.
The terms responsible for FSI can be isolated in $`H_A`$ rewriting the nuclear hamiltonian in the form
$$H_A=\underset{i=1}{\overset{A}{}}T_i+\underset{j>i=1}{\overset{A}{}}v_{ij}=H_0+H_1,$$
(5)
with (the knocked out nucleon is labelled with index 1)
$$H_0=\underset{i=1}{\overset{A}{}}T_i+\underset{j>i=2}{\overset{A}{}}v_{ij}=H_{A1}+T_1,$$
(6)
and
$$H_1=\underset{j=2}{\overset{A}{}}v_{1j}.$$
(7)
In the above equations $`T_i`$ and $`v_{ij}`$ denote the kinetic energy of the $`i`$-th nucleon and the interaction potential between nucleons $`i`$ and $`j`$, respectively. $`H_0`$ is the PWIA hamiltonian, describing the system containing (A-1) interacting spectators and the noninteracting knocked out nucleon, whereas the terms associated with FSI are included in $`H_1`$.
The decomposition of eq.(5) can be used to write the final scattering state $`|\mathrm{\Psi }_{n𝐩}^{()}`$, in the form :
$$|\mathrm{\Psi }_{n𝐩}^{()}=\mathrm{\Omega }_𝐩^{()}|\mathrm{\Phi }_{n𝐩},$$
(8)
where $`|\mathrm{\Phi }_{n𝐩}`$ denotes the asymptotic state with no interaction between particle $`1`$ and the spectators, which is obviously an eigenstate of $`H_0`$. In coordinate space it can be written
$$\mathrm{\Phi }_{n𝐩}(R)=\sqrt{\frac{1}{V}}e^{i𝐩𝐫_1}\phi _n(\stackrel{~}{R}),$$
(9)
where $`V`$ is the normalization volume, while $`R\{𝐫_1,𝐫_2\mathrm{},𝐫_A\}`$ and $`\stackrel{~}{R}\{𝐫_2,\mathrm{},𝐫_A\}`$ specify the configurations of the full A-particle system and the (A-1)-particle spectator system, respectively.
Setting $`\mathrm{\Omega }_𝐩^{()}=1`$, which amounts to disregarding the effects of FSI, and substituting into eq.(4), one obtains the PWIA amplitude, depending upon the the missing momentum $`𝐩_m=𝐩𝐪`$ only. The operator $`\mathrm{\Omega }_𝐩^{()}`$ describes the distortion of the asymptotyc wave function produced by the rescattering of the knocked out nucleon. It can be formally written as
$$\mathrm{\Omega }_𝐩^{()}=\underset{t\mathrm{}}{lim}e^{iH_At}e^{iH_0t}=\underset{t\mathrm{}}{lim}\widehat{T}\mathrm{e}^{i_0^t𝑑t^{}\widehat{H}_1(t^{})},$$
(10)
where $`\widehat{T}`$ denotes the time ordering operator and
$$\widehat{H}_1(t)=\mathrm{e}^{iH_0t}H_1\mathrm{e}^{iH_0t}.$$
(11)
In general, the calculation of $`\mathrm{\Omega }_𝐩^{()}`$ from eq.(10) with a realistic nuclear hamiltonian involves prohibitive difficulties. However, when the kinetic energy carried by the knocked out proton is large, the structure of $`\mathrm{\Omega }_𝐩^{()}`$ can be strongly simplified using a generalization of the approximation scheme originally developed by Glauber to decribe proton-nucleus scattering . The basic assumptions underlying this scheme are that i) the fast struck nucleon moves along a straight trajectory, being undeflected by rescattering processes (eikonal approximation) and ii) the spectator system can be seen as a collection of fixed scattering centers (frozen approximation).
Implementation of the eikonal and frozen approximations in the definition of the scattering operator $`\mathrm{\Omega }_𝐩^{()}`$, eq.(10), leads to the following coordinate space expression:
$`\mathrm{\Omega }_𝐩^{()}(R)`$ $`=`$ $`R|\mathrm{\Omega }_𝐩^{()}|R=P_z{\displaystyle \underset{j=2}{\overset{A}{}}}\left[1\mathrm{\Gamma }_p(1,j)\right]`$ (12)
$`=`$ $`P_z\left[1{\displaystyle \underset{j=2}{\overset{A}{}}}\mathrm{\Gamma }_p(1,j)+{\displaystyle \underset{k>j=2}{\overset{A}{}}}\mathrm{\Gamma }_p(1,j)\mathrm{\Gamma }_p(1,k)\mathrm{}\right],`$
where the positive $`z`$-axis is chosen along the eikonal trajectory and the $`z`$-ordering operator $`P_z`$ prevents the occurrence of backward scattering of the fast struck proton. The structure of the two-body operator $`\mathrm{\Gamma }_p(1,j)`$, describing the dynamics of the scattering process, will be discussed in the next section.
Inserting $`\mathrm{\Omega }_𝐩^{()}`$ of eq.(12) into the definition of $`|\mathrm{\Psi }_{n𝐩}^{()}`$ of eq.(8) one gets the following expression for the matrix element of eq.(4):
$$M_n(𝐩,𝐪)=𝑑R\left[\phi _n(\stackrel{~}{R})\mathrm{\Omega }_𝐩^{()}(R)\right]^{}\mathrm{e}^{i(𝐩𝐪)𝐫_1}\mathrm{\Psi }_0(R).$$
(13)
The calculation of the above amplitude can be simplified introducing a further approximation, whose validity rests on the same assumptions made to justify the use of the frozen approximation. Within this scheme , one replaces the many-body scattering operator $`\mathrm{\Omega }_𝐩^{()}(R)`$ with a one-body operator, depending on the position of the knocked out nucleon only, that can be obtained averaging $`\mathrm{\Omega }_𝐩^{()}(R)`$ over the positions of the spectator particles in the target ground state according to the following definition:
$$\overline{\mathrm{\Omega }}_𝐩^{()}(𝐫)=\frac{1}{\rho _A(𝐫)}𝑑R|\mathrm{\Psi }_0(R)|^2\mathrm{\Omega }_𝐩^{()}(R)\frac{1}{A}\underset{i=1}{\overset{A}{}}\delta (𝐫𝐫_i),$$
(14)
where $`\rho _A(𝐫)`$ is the target density normalized to unity.
Substitution of $`\mathrm{\Omega }_𝐩^{()}(R)`$ with $`\overline{\mathrm{\Omega }}_𝐩^{()}(𝐫)`$ in eq.(13) allows one to rewrite the amplitude in the form:
$$M_n(𝐩,𝐪)=d^3r\mathrm{e}^{i(𝐩𝐪)𝐫}\psi _{n𝐩}(𝐫),$$
(15)
the distorted overlap $`\psi _{n𝐩}(𝐫)`$ being defined as:
$$\psi _{n𝐩}(𝐫)=\left[\overline{\mathrm{\Omega }}_𝐩^{()}(𝐫)\right]^{}\chi _n(𝐫),$$
(16)
with
$$\chi _n(𝐫_1)=𝑑\stackrel{~}{R}\phi _n^{}(\stackrel{~}{R})\mathrm{\Psi }_0(R).$$
(17)
Note that, within the nuclear shell model picture, the quantity defined in eq.(17) can be interpreted as the wave function associated with the single particle state initially occupied by the knocked out nucleon .
The overlap relevant to the case of proton knock out from a <sup>4</sup>He target leading to a recoiling <sup>3</sup>H can be written:
$$\chi _0(𝐗)=d^3Yd^3Z\mathrm{\Psi }_3^{}(𝐘,𝐙)\mathrm{\Psi }_4(𝐗,𝐘,𝐙),$$
(18)
with $`𝐘=𝐫_2𝐫_3`$, $`𝐙=(2/3)𝐫_4(𝐫_2+𝐫_3)/3`$, $`𝐗=𝐫_1(𝐫_2+𝐫_3+𝐫_4)/3`$, whereas $`\mathrm{\Psi }_3(𝐘,𝐙)`$ and $`\mathrm{\Psi }_4(𝐗,𝐘,𝐙)`$ denote the ground state wave functions of <sup>3</sup>H and <sup>4</sup>He, respectively. The function $`\chi _0(𝐗)`$ has been evaluated by Schiavilla et al with highly realistic wave functions, obtained using the Variational Monte Carlo approach and nuclear hamiltonians including two- and three-nucleon interactions . We have used the results of ref. to calculate the amplitude of eq.(15) with the averaged scattering operator given by eq.(14), whose definition in the <sup>4</sup>He center of mass frame reads
$$\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)=\frac{d^3Yd^3Z|\mathrm{\Psi }_4(𝐗,𝐘,𝐙)|^2\mathrm{\Omega }_𝐩^{()}(𝐗,𝐘,𝐙)}{d^3Yd^3Z|\mathrm{\Psi }_4(𝐗,𝐘,𝐙)|^2}.$$
(19)
The integrations involved in the calculation of $`\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)`$ have been carried out using Monte Carlo configurations sampled from the probability distribution associated with the <sup>4</sup>He ground state wave function of ref..
### 2.2 Scattering operator
Within standard nonrelativistic nuclear many-body theory, i.e. treating the nucleons as pointlike structureless particles, the operator $`\mathrm{\Gamma }_p(1,j)`$ appearing in eq.(12) is a function of the particle positions $`𝐫_1`$ and $`𝐫_j`$ only. Choosing the $`z`$ axis along the direction of the eikonal trajectory (i.e. the direction of the momentum of the struck proton, specified by the unit vector $`𝐩/|𝐩|`$, the dependence of $`\mathrm{\Gamma }_p(1,j)`$ upon $`z_1`$ and $`z_j`$ can be singled out writing
$$\mathrm{\Gamma }_p(1,j)=\theta (z_jz_1)\gamma _p(|𝐛_1𝐛_j|),$$
(20)
where the step function preserves causality while $`\gamma _p(b)`$ is a function of the projection of the interparticle distance in the impact parameter plane (the $`xy`$ plane) which contains all the information on the dynamics of the scattering process. The function $`\gamma _p(b)`$ can be simply related to the coordinate space $`t`$-matrix associated with the proton-nucleon (pN) potential $`v_{ij}`$, and written in terms of the measured NN scattering amplitude at incident momentum $`p`$, $`f_p(k_t)`$, as
$$\gamma _p(b)=\frac{i}{2}\frac{d^2k_t}{(2\pi )^2}\mathrm{e}^{i𝐤_t𝐛}f_p(k_t).$$
(21)
At large $`p`$, the experimental $`f_p(k_t)`$ is usually parametrized in the form
$$f_p(k_t)=i\sigma _{pN}^{tot}(1i\alpha _{pN})\mathrm{e}^{\frac{1}{2}\frac{k_t^2}{B}},$$
(22)
where $`\sigma _{pN}^{tot}`$ and $`\alpha _{pN}`$ denote the total cross section and the ratio between the real and the imaginary part of the amplitude, respectively, while $`B`$ is related to the range of the interaction. In the case of zero-range interaction, $`B=0`$ and the impact parameter dependence of $`\gamma _p(b)`$ reduces to a two-dimensional $`\delta `$-function.
To include CT, the internal structure of the proton must be explicitely taken into account. According to the CT scenario, in the $`(e,e^{}p)`$ reaction at large $`Q^2`$ the electromagnetic interaction produces a compact three-quark state $`|E`$, which can be seen as a superposition of many hadronic states $`|\alpha `$, $`|\beta \mathrm{}`$. This state then propagates through the nuclear medium undergoing rescattering processes that eventually lead to the emergence of the detected proton. The rescattering processes can be either diagonal, when the hadronic state $`|\alpha `$ does not change, or off diagonal, when a transition to a different state $`|\beta `$ is induced. The transparency effect, i.e. the disappearance of nuclear absorption, folows from the cancelation between the contributions of diagonal and off diagonal processes at asymptotically high $`Q^2`$.
From the above discussion , it follows that, in order to describe the transition of FSI effects to the CT regime, one has to introduce a scattering operator acting in the space of the hadronic states. Its matrix element between states $`|\beta `$ and $`|\alpha `$, of mass $`m_\beta `$ and $`m_\alpha `$, respectively, can be defined as
$$\beta |\mathrm{\Gamma }_p(1,j)|\alpha =\theta (z_jz_1)\mathrm{e}^{ik_{\alpha \beta }z_j}\gamma _p^{\alpha \beta }(|𝐛_1𝐛_j|),$$
(23)
where
$$k_{\alpha \beta }=\frac{m_\alpha ^2m_\beta ^2}{2E_p},$$
(24)
$`E_p=T_p+m`$ being the energy of the detected proton in the laboratory frame. The onset oc CT is driven by the oscillating factors exp($`ik_{\alpha \beta }z_j`$), taking into account the longitudinal momentum transfer associated with each transition $`\alpha `$ \+ N $`\beta `$ \+ N. In analogy to eq.(21), $`\gamma _p^{\alpha \beta }(b)`$ is written in terms of the amplitude of the process $`\alpha +N\beta +N`$:
$$\gamma _p^{\alpha \beta }(b)=\frac{i}{2}\frac{d^2k_t}{(2\pi )^2}\mathrm{e}^{i𝐤_t𝐛}f_p^{\alpha \beta }(k_t),$$
(25)
with
$$f_p^{\alpha \beta }(k_t)=i\beta |\widehat{\sigma }|\alpha (1i\alpha _{\alpha \beta })\mathrm{e}^{\frac{1}{2}\frac{k_t^2}{B_{\alpha \beta }}}.$$
(26)
In the above equation, $`\alpha _{\alpha \beta }`$ and $`B_{\alpha \beta }`$ are the generalization of the parameters $`\alpha _{pN}`$ and $`B`$ of eq.(22), while the operator $`\widehat{\sigma }`$ describes the hadronic cross section. Unfortunately, $`\alpha _{\alpha \beta }`$ and $`B_{\alpha \beta }`$ are not known experimentally. In our numerical calculations we have made the assumption that the interactions responsible for off diagonal rescatterings have zero range, i.e. that $`B_{\alpha \beta }=0`$ for any $`\alpha \beta `$. The values of $`\alpha _{\alpha \beta }`$ have been varied within a reasonable range to gauge the sensitivity of our approach to these parameters. The results will be discussed in the next section.
Following ref., $`\beta |\widehat{\sigma }|\alpha `$ has been evaluated in configuration space, using
$$\beta |\widehat{\sigma }|\alpha =𝑑\xi d^2\rho \psi _\beta ^{}(\xi ,\rho )\sigma (\rho )\psi _\alpha (\xi ,\rho ),$$
(27)
where $`\psi _\alpha `$ and $`\psi _\beta `$ are harmonic oscillator wave functions describing a quark-diquark system with longitudinal and transverse coordinates $`\xi `$ and $`\rho `$, respectively. The quark-diquark oscillation frequency has been chosen to be $`\omega _0=`$ 0.35 GeV , yielding a realistic mass spectrum of the proton excited states, while $`\sigma (\rho )`$ has been parametrized in the form
$$\sigma (\rho )=\sigma _0\left[1\mathrm{e}^{\left(\frac{\rho }{\rho _0}\right)^2}\right],$$
(28)
with $`\sigma _0=2\sigma _{pN}`$ and $`\rho _0`$ adjusted in such a way as to reproduce the experimental pN total cross section.
A scattering operator suitable to describe the onset of CT, denoted $`\mathrm{\Omega }_{CT}(R)`$, can be constructed using eq.(12) and the two-body scattering operators $`\mathrm{\Gamma }_p(1,j)`$ whose matrix elements are defined by eq.(23):
$`\mathrm{\Omega }_𝐩^{CT}(R)`$ $`=`$ $`{\displaystyle \frac{P_zp|_{j=2}^A\left[1\mathrm{\Gamma }_p(1,j)\right]|E}{p|E}}`$ (29)
$`=`$ $`1P_z{\displaystyle \underset{j=2}{\overset{A}{}}}{\displaystyle \underset{\alpha }{}}p|\mathrm{\Gamma }_p(1,j)|\alpha {\displaystyle \frac{\alpha |E}{p|E}}`$
$`+`$ $`P_z{\displaystyle \underset{k>j=2}{\overset{A}{}}}{\displaystyle \underset{\alpha \beta }{}}p|\mathrm{\Gamma }_p(1,k)|\beta \beta |\mathrm{\Gamma }_p(1,j)|\alpha {\displaystyle \frac{\alpha |E}{p|E}}+\mathrm{},`$
where $`|p`$ is the state describing the detected proton. For the compact state $`|E`$ produced at the electromagnetic vertex, we have used the same configuration space wave function employed in refs.
$$\rho |E\mathrm{e}^{C\rho ^2Q^2},$$
(30)
with $`C=1`$. It has to be pointed out that, as shown in ref., the missing momentum distribution is not sensitive to the choice of $`C`$ as long as $`C`$ 0.05.
## 3 Results
Using the functions $`\chi _0(𝐗)`$ and $`\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)`$ defined by eqs.(18) and (19), respectively, the missing momentum distribution associated with the $`{}_{}{}^{4}He(e,e^{}p)^3H`$ process, denoted $`W_𝐩(𝐩_m)`$, can be readily obtained from
$$W_𝐩(𝐩_m)=\left|d^3X\mathrm{e}^{i𝐩_m𝐗}\chi _0(𝐗)\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)\right|^2.$$
(31)
The results discussed in the present paper have been obtained using the overlap $`\chi _0(𝐗)`$ computed in ref. using the Argonne v14 two-nucleon interaction and the Urbana VII three-body potential. The scattering operator $`\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)`$ has been calculated carrying out the integrations involved in eq.(19) with the Monte Carlo method, using a configuration set sampled from the probability distribution associated with the <sup>4</sup>He wave function of ref..
In figs. 1 and 2 we show $`W_𝐩(𝐩_m)`$ evaluated at the top of the quasi free peak, i.e. at $`\omega =Q^2/2m`$, for parallel ($`p_{m,}=|𝐩_m\times 𝐪|=0`$) and perpendicular ($`p_{m,z}=|𝐩_m𝐪|=0`$) kinematics, respectively. Each figure has four panels, corresponding to different values of $`Q^2`$ ranging from 2 (GeV/c)<sup>2</sup> to 20 (GeV/c)<sup>2</sup>. The dotted line shows the PWIA (i.e. $`\overline{\mathrm{\Omega }}_𝐩^{()}(𝐗)1`$) result, whereas the dashed and solid curves correspond to the calculations including FSI effects with and without CT, respectively. The configuration set employed in our calculations allows for an accurate determination of the missing momentum distribution over a large momentum range. However, in the region where $`W_𝐩`$ becomes very small ($`10^4`$ fm<sup>3</sup>), the statistical uncertainty of the Monte Carlo calculation becomes sizeable. Although our results indicate that the first two-three excited states saturate the contribution of the off-diagonal rescatterings at missing momentum less than 300 MeV/c, we have included six intermediate states in all numerical calculations.
Fig. 1 shows that, while within PWIA $`W_𝐩(p_{m,z})=W_𝐩(p_{m,z})`$, FSI produce a forward bakward asymmetry, whose origin has to be ascribed to the effect of the real part of the NN scattering amplitude and to the fact that the cancellation between the contributions of diagonal and off-diagonal rescattering processes depends upon the value of the missing momentum .
The main features of the distorted missing momentum distribution are the quenching in the region of the maximum, corresponding to $`p_{m,z}`$ 0, and the enhancement of the tail at negative $`p_{m,z}`$. As expected, inclusion of CT reduces the effect of FSI.
A more complicated structure is observed in the case of perpendicular kinematics, shown in fig.2. The mimimum displayed by the PWIA missing momentum distribution, almost completely washed out by FSI, reappears at lower values of $`p_{m,}`$ when the effect of CT is included. Note that at large $`p_{m,}`$ ($`p_{m,}2`$ fm$`{}_{}{}^{}1`$) the momentum distribution is dominated by FSI and that the inclusion of CT results in a sizable suppression.
The complex pattern of quenching and enhancement of the PWIA $`W_𝐩`$ has to be ascribed to the combined effects of FSI and strong NN correlations in the initial state. A similar behavior has been found in ref., where the semi-inclusive process $`{}_{}{}^{4}He(e,e^{}p)X`$ has been analyzed using the somewhat simplified Jastrow model to describe NN correlations. More recently, the $`{}_{}{}^{4}He(e,e^{}p)X`$ reaction has been studied using a four-body wave function including noncentral correlations induced by the tensor component of the NN interaction . The distorted momentum distributions of refs. and exhibit the same pattern of FSI effects on the S-wave contribution, while the D-wave is only weakly distorted. The main difference between the two approaches cancellation between the effects of central and tensor correlations at large missing momenta, leading to a suppression of the distortion by a factor of $``$ 2.
The mimimum displayed by the PWIA missing momentum distribution, almost completely washed out by FSI, reappears at lower values of $`p_{m,}`$ when the effect of CT is included. Figs. 1 and 2 show that at $`Q^2`$ = 2 (GeV/c)<sup>2</sup> and in absence of CT the missing momentum distribution at $`p_{m,}=p_{m,z}=`$ 0 gets quenched by 18 $`\%`$ on account of FSI.
The effects of FSI can be best observed in the ratio
$$T_p(𝐩_m)=\frac{W_𝐩(𝐩_m)}{\left|d^3X\mathrm{e}^{i𝐩_m𝐗}\chi _0(𝐗)\right|^2},$$
(32)
shown in figs. 3 and 4. It clearly follows from the definition that within PWIA $`T_p(𝐩_m)1`$. When FSI are included, $`T_p(𝐩_m)`$ is a function of the missing momentum and can take values both above and below than unity, reflecting the fact that the distorted missing momentum distribution can be larger or smaller than the PWIA momentum distribution. Hence, in spite of the analogy between eq.(32) and the definition of the nuclear transparency, the quantity 1 - $`T_p(𝐩_m)`$ cannot be simply interpreted as the nuclear absorption experienced by a proton carrying momentum $`𝐩`$.
The calculated $`T_p`$ corresponding to parallel and perpendicular kinematics are presented in figs. 3 and 4, respectively. The dotted curve shows the results obtained when no off-daigonal rescatterings are included, i.e. in absence of CT. The solid, dashed and long-dashed curves correspond to calculations in which CT effects have been taken into account by properly including off-diagonal rescatterings and using three different sets of parameters $`(\alpha _1,\alpha _2)`$ to describe the real part of the scattering amplitudes associated with diagonal ($`\alpha +N\alpha +N`$, with $`\alpha p`$) and off-diagonal ($`\alpha +N\beta +N`$, with $`\alpha \beta `$) processes. These amplitudes have been parametrized according to:
$`Ref(\alpha +N\alpha +N)`$ $`=`$ $`\alpha _1Ref(p+Np+N)`$ (33)
$`=`$ $`{\displaystyle \frac{\alpha _1}{3}}\left[\alpha _{pp}\sigma _{pp}+2\alpha _{pn}\sigma _{pn}\right],`$
and
$$Ref(\alpha +N\beta +N)=\alpha _2Imf(\alpha +N\beta +N).$$
(34)
The results of our calculations turn out to be insensitive to the value of $`\alpha _1`$, so we set $`\alpha _1=1`$ and show the effect of varying $`\alpha _2`$ only. The solid, dashed and long dashed curves of Fig. 3 correspond to the sets $`(\alpha _1,\alpha _2)`$ = (1,0),(1,0.5) and (1,-0.5). In parallel kinematics the dependence upon $`\alpha _2`$ appears to be sizable, particularly at the largest value of $`Q^2`$. On the other hand, Fig. 4 shows that in perpendicular kinematics, where the effect of CT at large missing momentum is large, the theoretical error bar associated with the uncertainty in $`\alpha _2`$ is rather small.
In fig. 5, we present the longitudinal forward-backward asymmetry of the missing momentum distribution defined as
$$A_z(x,y)=\frac{N_+N_{}}{N_++N_{}},$$
(35)
where
$$N_\pm =_{\pm x}^{\pm y}𝑑p_{m_z}W_𝐩\left(p_{m_z},p_{}=0\right).$$
(36)
We have evaluated $`A_z`$ from the above equations for four kinematical windows: $`(x,y)`$ = (0,0.3), (0,0.4), (0.05,0.3) and (0.1,0.4) GeV/c. To illustrate the contribution of the off-diagonal rescattering processes to $`A_z`$, the results obtained setting these contributions to zero are shown by the dotted line. As in fig. 3 and 4, the solid, dashed and long-dashed lines correspond to different choices of the parameter $`\alpha _2`$. The results of fig. 5 show that even at moderate $`Q^22.5`$ (GeV/c)<sup>2</sup>, i.e. in the region relevant to the $`(e,e^{}p)`$ experimental program at the Thomas Jefferson National Accelerator Facility (TJNAF), off-diagonal rescatterings are responsible for more than 60 $`\%`$ of the calculated asymmetry for all of considered kinematical windows. In a previous study of the asymmetry in semi-inclusive $`(e,e^{}p)`$ processes we have found a contribution of 15-20 $`\%`$ and 10-15 $`\%`$ in the case of $`{}_{}{}^{16}O`$ and $`{}_{}{}^{40}Ca`$, respectively . The results of the present calculation confirm our conclusion that CT effects are larger in light nuclei. At $`Q^2le20`$ (GeV/c)<sup>2</sup>, our calculations show a weak dependence of $`A_z`$ upon the value of the parameter $`\alpha _2`$. Hence, the uncertainty associated with the chioce of $`\alpha _2`$ does not prevent one from extracting an unambiguous signature of the onset of CT from the asymmetry.
## 4 Summary and conclusions
We have carried out a calculation of the missing momentum distribution of the process <sup>4</sup>He$`(e,e^{}p)^3`$H in quasielastic kinematics, i.e. at $`\omega Q^2/2m`$, in the range $`2Q^220`$ (GeV/c)<sup>2</sup>.
The PWIA momentum distribution, evaluated using highly realistic many-body wave functions and the Monte Carlo method, has been corrected to include the effects of FSI, treated within a coupled-channel multiple scattering approach suitable to describe the possible onset of CT. It has to be emphasized that our approach allows for a consistent treatment of short range NN correlations in both the initial and final state.
Inclusion of FSI leads to the appearance of a complex pattern of distorsions of the PWIA momentum distribution, both in parallel and perpendicular kinematics. Sizable CT effects are observed in perpendicular kinematics at $`p_m1.5`$ fm<sup>-1</sup> over the whole $`Q^2`$ range. The CT signal turns out to be much larger than the theoretical uncertainty associated with the parametrization of the amplitudes for off-diagonal rescattering. The calculated forward-backward asymmetry also shows a significant CT effect already at $`Q^2`$ 2.5 (GeV/c)<sup>2</sup>.
In conclusion, our results seem to indicate that the experimental study of the exclusive channels in $`(e,e^{}p)`$ processes off few-nucleon system may give a clue to the issue of the possible manifestation of CT in the domain of moderate $`Q^2`$ ($``$ 3 (GeV/c)<sup>2</sup>), covered by the existing electron scattering facilities.
It is a pleasure to thank Rocco Schiavilla and Robert B. Wiringa for providing Monte Carlo configurations of the <sup>4</sup>He ground state and the results of their calculation of the $`^3H|^4He`$ overlap. One of the authors (AAU) gratefully acknowledges the hospitality provided by the Sezione Sanità of the Italian National Institute for Nuclear Research (INFN) and the Institut für Kernphysik, Forschungszentrum Jülich, where part of the work described in this paper has been carried out. |
warning/0001/gr-qc0001035.html | ar5iv | text | # ON UPPER LIMITS FOR GRAVITATIONAL RADIATION
## 1 Introduction
After the initial experiments with room temperature resonant detectors, the new generation of cryogenic gravitational wave (GW) antennas entered long term data taking operation in 1990 (EXPLORER ), in 1991 (ALLEGRO ), in 1993 (NIOBE ), in 1994 (NAUTILUS ) and in 1997 (AURIGA ).
Searches for coincident events between detectors have been performed. Between EXPLORER and NAUTILUS and between EXPLORER and NIOBE in the years 1995 and 1996 . Between ALLEGRO and EXPLORER with data recorded in 1991 . In both cases no significative coincidence excesses were found and an upper limit to GW bursts was calculated .
However, the upper limit determination has been done under the hidden hypothesis that the signal-to-noise ratio (SNR) is very large. According to theoretical estimations the signals expected from cosmic GW sources are extremely feeble, so small that extremely sensitive detectors are needed. In fact, according to present knowledge, the detectors available today have not yet reached the sensitivity to detect even a few events per year.
Thus it is important to study the problem of the upper limit determination in the cases the SNRs of the observed $`events`$ are not large. In order to do this we have to discuss our definition of $`event`$.
The raw data from a resonant GW detector are filtered with a filter matched to short bursts . We describe now in more detail the procedure used for the GW detectors of the Rome group, EXPLORER and NAUTILUS.
After the filtering of the raw-data, $`events`$ are extracted as follows. Be $`x(t)`$ the filtered output of the electromechanical transducer which converts the mechanical vibrations of the bar in electrical signals. This quantity is normalized, using the detector calibration, such that its square gives the energy innovation $`E_f`$ of the oscillation for each sample, expressed in kelvin units. In absence of signals, for well behaved noise due only to the thermal motion of the bar and to the electronic noise of the amplifier, the distribution of $`x(t)`$ is normal with zero mean. The variance (average value of the square of $`x(t)`$) is called $`effectivetemperature`$ and is indicated with $`T_{eff}`$. The distribution of $`x(t)`$ is
$$f(x)=\frac{1}{\sqrt{2\pi T_{eff}}}e^{\frac{x^2}{2T_{eff}}}$$
(1)
For extracting $`events`$ (in absence of signals the events are just due to noise) we set a threshold in terms of a critical ratio defined by
$$CR=\frac{|x|<|x|>}{\sigma (|x|)}=\frac{\sqrt{SNR_f}\sqrt{\frac{2}{\pi }}}{\sqrt{1\frac{2}{\pi }}}$$
(2)
where $`\sigma (|x|)`$ is the standard deviation of $`|x|`$ and we put
$$SNR_f=\frac{E_f}{T_{eff}}$$
(3)
The threshold is set at a value CR such to obtain, in presence of thermal and electronic noise alone, a number of events which can be easily exchanged among the other groups who participate to the data exchange. For about one hundred $`events`$ per day the threshold corresponds to an energy $`E_t=19.5T_{eff}`$.
We calculate now the theoretical probability to detect a signal with a given SNR, in presence of a well behaved Gaussian noise. We put $`y=(s+x)^2`$ where $`s\sqrt{SNR}`$ is the signal we look for and $`x`$ is the gaussian noise. We obtain easily
$$probability(SNR)=_{SNR_t}^{\mathrm{}}\frac{1}{\sqrt{2\pi y}}e^{\frac{(SNR+y)}{2}}cosh(\sqrt{ySNR})𝑑y$$
(4)
We put $`SNR_t=19.5`$ for the present EXPLORER and NAUTILUS detectors.
## 2 Upper limit determination
We consider M detectors and search for M-fold coincidences over a total period of time $`t_m`$ during which all detectors are in operation. Be $`\overline{n}`$ the average number of accidental coincidences (due to chance) and $`n_c`$ the number of coincidences which are found within a given time window.
For events which have a Poissonian distribution in time the expected average number of M-fold accidental coincidences is given by
$$\overline{n}=Mw^{M1}\underset{k}{\overset{1,M}{}}n_k$$
(5)
where $`n_k`$ is the event density of the $`k^{th}`$ detector.
The accidental coincidence distribution can be estimated experimentally by proper shifting the event occurrence times of each detector. In the case of Poissonian distribution the average number of the M-fold accidental coincidences coincides with that given by eq. 5. The comparison between $`n_c`$ and $`\overline{n}`$ allows to reach some conclusion about the detection of GW or to establish an upper limit to their existence.
In paper and in the previous paper the upper limit has been estimated as follows. It has been found that, for various energy levels of the observed events, the number $`n_c`$ was smaller than or did not exceeded significantly $`\overline{n}`$. Such numbers $`n_c`$, one for each energy level, were used for calculating the upper limit. A Poissonian distribution of the number of the observed events was considered together with the hypothesis of an isotropic distribution in the sky of the GW sources. The value of $`h`$ (adimensional perturbation of the metric tensor) was then derived $`from`$ the energy levels, using the detector cross-section for gravitational waves.
This procedure can be objected on two points:
a)The most important point is that, as shown in , for SNR small and up to values of a few dozens, the energy of an event is $`not`$ the energy of the GW absorbed by the detector. This means that we cannot deduce the value of $`h`$ directly from the energy levels of the observed events;
b)In addition, the efficiency of detection, again for SNR values up to one or two dozens, is rather smaller than unity, and this changes the upper limit, particularly at small SNR.
We introduce a new procedure for estimating the upper limit, which circumvents the difficulties indicated in the above two points.
The problem to determine the upper limit has been discussed in several papers. In particular in paper , as indicated by the PDG, and, more recently, in paper . According to the upper limit can be calculated using the relative belief updating ratio
$$R(n_{GW},n_c,\overline{n})=e^{n_{GW}}(1+\frac{n_{GW}}{\overline{n}})^{n_c}$$
(6)
referring to a given period $`t_m`$ of data taking. This function is proportional to the likelihood and it allows to infer the probability to have $`n_{GW}`$ signals for given priors (using the Bayes’s theorem). It has already been used in High Energy Physics .
We calculate the upper limit by solving the equation
$$R(n_{GW},n_c,\overline{n})=0.05$$
(7)
We remark that 5% does not represent a probability but it is an useful way to put a limit independently on the priors<sup>1</sup><sup>1</sup>1 To avoid confusion we shall continue to use the words $`upperlimit`$, although it would be more appropriate to call it $`standardsensitivitybound`$ ..
Eq. 7 has a very interesting solution. Putting $`n_c=0`$ we find $`n_{GW}=2.99`$, independent on the value of the background $`\overline{n}`$. If we use the calculations of ref. we find that, for $`n_c=0`$ and $`\overline{n}=0`$, the upper limit is 3.09 (almost identical to the previous one) $`but`$ it decreases for increasing $`\overline{n}`$. The reason for this different behavior is due to the non- Bayesan character of the calculations made in , as we discuss in the following.
Suppose we have $`n_c=0`$ and $`\overline{n}0`$. This certainly means that the number of accidentals, whose average value can be determined with any desired accuracy, has undergone a fluctuation. For larger $`\overline{n}`$ values, smaller is the ($`apriori`$) probability that such fluctuation occur. Thus one could reason that it is less likely that a number $`n_{GW}`$ be associated to a large value of $`\overline{n}`$, since the observation gave $`n_c=0`$.
According to the Bayesan approach instead, as discussed in , one cannot ignore the fact that the observation $`n_c=0`$ had already being made at the time the estimation of the upper limit is considered. The Bayesan approach requires that, given $`n_c=0`$ and $`\overline{n}0`$, one evaluate the $`chance`$ that a number $`n_{GW}`$ of signals exist. This $`chance`$ of a possible signal is referred to the observation already made and, rather obviously, it cannot depend on the previous fluctuation of the background, since the presence of a signal cannot be related to the background due to the detector. Mathematically, it is easy to demonstrate, using the results obtained in , that due to the Poissonian character of the number of accidentals this $`relativechance`$ (for $`n_c=0`$) is indeed independent on $`\overline{n}`$.
It can be seen, comparing the results of with those of , that the Bayesan upper limits are for all values of $`n_c`$ and $`\overline{n}`$ (except $`n_c=\overline{n}=0`$), greater than those obtained with the non-Bayesan procedure. In our opinion the Bayesan approach has to be preferred, and so we do in this paper.
If we have $`n_c0`$ then we apply eq. 6. It is interesting to show the result for the case $`n_c=\overline{n}0`$ for the standard sensitivity bound of 5%. The result is given in fig.1
We note that for $`n_c=\overline{n}`$ and $`n_{GW}<<\overline{n}`$ eq. 6 can be approximated with
$$n_{GW}\sqrt{6\overline{n}}$$
(8)
From the result shown in fig.1 it appears evident that the lowest upper limit is obtained for $`n_c\overline{n}0`$. In order to obtain $`\overline{n}0`$ one can raise the threshold used for determining the events. However in doing this one diminish the efficiency of detection, as shown in eq.4. Whether the procedure to raise the threshold is convenient or not, it depends on the numerical effects of the two competing operations. Certainly for large GW signals, when the detection efficiency is always unity, it is much better to have a threshold that gives $`\overline{n}=0`$. For smaller signals one has to consider specific cases. However it can be seen that in the most interesting cases it is better to raise the threshold until we get $`\overline{n}0`$. This will be shown in the section where we reconsider the upper limit obtained with ALLEGRO and EXPLORER in 1991 .
In the estimation of the upper limit we consider the efficiency of detection, which we indicate with $`ϵ_k(SNR)`$ where $`k`$ refers to the $`k^{th}`$ detector. For EXPLORER and NAUTILUS the theoretical efficiency is obtained from eq. 4.
We must relate the $`h`$ values of the GW to the energy $`E`$ absorbed by the detectors. We have to consider that the absorbed energy depends on the direction of the impinging GW and on its polarization. For taking care of the various polarization we use the average value dividing the cross section by a factor of two. We then have
$$h=1.1310^{17}\sqrt{E}$$
(9)
with the energy $`E`$ expressed in kelvin unit. This formula is valid only if the GW arrives perpendicularly to the detector axis ($`\theta =90^o`$). For a given direction we calculate the absorbed energy using the $`sin(\theta )^4`$ dependency. We also consider that for an isotropic distribution of sources the number of possible GW impinging directions is proportional to $`sin(\theta )^2`$.
The procedure for calculating the upper limit is accomplished thru the following points:
a) consider various values of $`h`$;
b) assume an isotropic distribution of the GW sources;
c) for each direction $`\theta `$ and for each $`h`$ calculate the absorbed energy $`E(\theta )`$ by means of eq. 9 and the $`sin^4(\theta )`$ dependency;
d) for each detector calculate the SNR for the adsorbed energy by taking into consideration the noise $`T_{eff,k}`$:
$$SNR_k(\theta )=\frac{E(\theta )}{T_{eff,k}},k=1,..,M$$
(10)
e) using the individual efficiencies $`ϵ_k(SNR_k(\theta ))`$ consider the total efficiency $`ϵ_t(\theta )=_k^{1,M}ϵ_k(SNR_k(\theta ))`$;
f) integrate $`ϵ_t(\theta )`$ over $`\theta `$ with the weight $`sin^2(\theta )`$, because of the assumed isotropic distribution of the sources;
g) from eq.6, given $`n_c`$ and $`\overline{n}`$, we obtain $`n_{GW}`$. We then divide $`n_{GW}`$ by the result of point f) and obtain for each value of $`h`$ the upper limit during the measuring time $`t_m`$.
We remark that in this case we have not used the energy of the observed events, as done instead previously .
The total efficiency is calculated with the following eq. 11.
$$ϵ_{tot}(h)=\frac{_0^{\frac{\pi }{2}}_k^{1,M}ϵ_k(SNR_k(\theta ))sin^2(\theta )d\theta }{\frac{\pi }{4}}$$
(11)
For more clarity we show in table 1 some of the steps needed for our calculation, using two parallel detectors and $`n_c=0`$. We use the efficiency given by eq. 4, valid for a well behaved noise<sup>2</sup><sup>2</sup>2The real data often show a non gaussian behaviour. In this case the efficiency differs from the theoretical one given by eq.4, but one can easily make use of the efficiency experimentally measured..
## 3 Ricalculation of the upper limit with the data of ALLEGRO and EXPLORER in 1991
In a previous paper the upper limit for GW bursts was calculated, using the data recorded by ALLEGRO and EXPLORER in 1991. We wish now to recalculate the upper limit according the considerations discussed in this paper.
In 1991 the EXPLORER data filtering was done differently from that described in this Introduction. For both ALLEGRO and EXPLORER the output of the electromechanical transducer was sent to lock-ins referred to the frequencies of the resonant modes. Then the outputs of the lock-ins (in phase and in quadrature) were filtered searching for delta-like signals and combined for obtaining the energy innovation, which we still indicate with $`E_f`$. In this case the probability to have an event (above threshold $`SNR_t`$) due to a signal with given SNR is obtained (see ref. ) with the following equation:
$$probability(SNR)=_{SNR_t}^{\mathrm{}}e^{(SNR+y)}I_o(2\sqrt{ySNR})𝑑y$$
(12)
Here $`y=\frac{E_f}{T_{eff}}`$, $`I_o`$ is the modified Bessel function of order zero, and the noise temperature $`T_{eff}`$ is the average value of the energy innovation $`E_f`$.
We recall that in a time period of 123 days 70 coincidences were found with a background of 59.3. For extracting the events the ALLEGRO threshold was $`SNR_t=11.5`$ with a noise temperature $`T_{eff}8mK`$. For EXPLORER the threshold was $`SNR_t=10`$ also with $`T_{eff}8mK`$. Applying eq.6 we find an upper limit of $`n_{GW}=37`$ over the 123 days.
According to the previous considerations we can raise the event threshold, say for EXPLORER, in order to reduce the number of accidentals. For instance, for a threshold $`SNR_t=24`$ we get $`n_c=1`$ and $`\overline{n}=0.74`$, obtaining, from eq. 6, the value $`n_{GW}=4.8`$.
Thus the procedure for calculating the upper limit with the Bayesan approach when we have data at various thresholds, including cases with $`n_c`$ and $`\overline{n}`$ different from zero, is the following.
Start with $`n_c`$ and $`\overline{n}`$ for various thresholds and use eq.6 for obtaining $`n_{GW}`$ at each threshold. Calculate the upper limit for various values of $`h`$ as shown in the previous section. For each $`h`$ take as upper limit the smallest value among those obtained by varying the threshold. Clearly at large $`h`$ values, when we get $`n_c=0`$, the upper limit is, for the entire period of time, $`n_{GW}=2.99`$.
The result is shown in fig.2 together with that obtained previously in . It turns out that the two upper limits are similar.
The reason for this is due to the fact that in applying the previous algorithm we started from an energy level higher than the largest energy of the detected (accidental) coincidences, thus obtaining, at this level ($`n_c=\overline{n}=0`$) an upper limit of 3.09 very close to the value 2.99 obtained with the Bayesan approach. The similarity of the results at lower $`h`$ values is accidental. In the previous algorithm the increase at lower $`h`$ is due only to the increase of the number $`\overline{n}`$ of accidentals. In the present algorithm the increase is due to the smaller efficiency of detection and to the increase in $`n_{GW}`$ which roughly goes with $`\sqrt{\overline{n}}`$ (eq.8).
In spite of the similar numerical results, we believe that the procedure proposed here which does not extract the value of $`h`$ from the energy levels of the accidental coincidences and it uses the Bayesan approach is methodologically more correct.
## 4 Discussion
The best upper limit which can be obtained with an array of $`M`$ identical parallel detectors in $`M^{pl}`$ coincidence cannot go below the value 2.99, because this is the upper limit when one finds $`zero`$ coincidences independently on the background.
The basic advantage in using many detectors comes from the fact that with many detectors it is easier to obtain $`\overline{n}0`$, and thus (in absence of GW) $`n_c=0`$. Because of the Poisson distributions, the average number of accidental coincidences for M detectors in a time window $`\pm w`$ is given by eq.5. On the time scale of 1 second (w=1 s) it turns out that $`n_k<<1`$. By increasing the number of detectors one obtains smaller values of $`\overline{n}`$, thus approaching the requirement to have $`n_c=0`$ and then the lowest possible upper limit.
This is certainly true at large $`h`$ values, where the detection efficiency for all detectors is unity. The result, as shown in fig.2, is a plateau. Instead it might be convenient at low $`h`$ values to use the two most sensitive detectors, in order to have the largest possible efficiency of detection. The overall upper limit is then obtained by taking the smallest ones among the values of the various upper limit determinations.
The above procedure can be easily adjusted to the more general case of any distribution of the GW sources, and of non-parallel detectors.
## 5 Acknowledgements
We have benefited from useful discussions with P.Bonifazi, G. D’ Agostini and F. Ronga. |
warning/0001/hep-th0001061.html | ar5iv | text | # Biconformal Matter Actions
## 1 Introduction
Recently, we developed a new gauge theory of the conformal group, which solved many of the problems typically associated with scale invariance . In particular, this new class of biconformal geometries has been shown to resolve the problem of writing scale-invariant vacuum gravitational actions in arbitrary dimension without the use of compensating fields . In the cited work, we wrote the most general linear vacuum action and completely solved the resulting field equations subject only to a minimal torsion assumption. We found that all such solutions were foliated by equivalent $`n`$-dimensional Ricci-flat Riemannian spacetimes.
Reference left an open question: how are matter fields coupled to biconformal gravity? A priori, it is not at all obvious that any action for biconformal matter permits the same embedded $`n`$-dim Riemannian structure that occurs for the vacuum case since biconformal fields are $`2n`$-dimensional. Indeed, in the case of standard $`n`$-dim conformal gauging ( -), we generally require compensating fields to recover the Einstein equation with matter (see, for example, -). To answer this question for biconformal space, in the present work we extend the results of by introducing a set of Klein-Gordon-type fields $`\varphi ^m`$ of conformal weight $`m`$ into the theory. Using the Killing metric intrinsic to biconformal space, we write the natural kinetic term in the biconformally covariant derivatives of $`\varphi ^m`$ and find the resulting gravitationally coupled field equations. Then, for the case of one scalar field $`\varphi `$ of conformal weight zero, we completely solve the field equations, under the assumption of vanishing torsion. We find that, as before, the solutions are foliated by equivalent $`n`$-dim Riemannian spacetime submanifolds whose curvatures now satisfy the usual Einstein equations with scalar matter. The field $`\varphi ,`$ which a priori depended on all $`2n`$ biconformal coordinates, is completely determined by the $`n`$ coordinates of the submanifolds and satisfies the submanifold massless Klein-Gordon equation $`\eta ^{cd}D_cD_d\varphi =0`$.
Thus, the new gauging establishes a clear connection between conformal gauge theory and general relativity with scalar matter, without the use of compensating fields.
The structure of the paper is as follows. In the next section, we extend the applicability of the biconformal dual introduced in and establish how to write the usual kinetic action for scalar fields of arbitrary conformal weight using differential forms. Then, in Sec.(3) we find the field equations resulting from this action coupled to the linear gravity action introduced in . Interestingly, these equations together with previous results show that pairs of scalar matter fields of conjugate conformal weight provide a source for the Weyl vector. Next, restricting to the case of a zero-weight scalar field, Secs.(4) and (5) give the solutions for the curvature and connection, respectively, of the background geometry in the case of vanishing torsion. Finally, in Sec.(6), we examine the field equations and other constraints on the matter field.
## 2 The biconformal dual and inner product
For full detail on the new conformal gauging we refer to . We will use the same notation as in .
The Minkowski metric is written as $`\eta _{ab}=`$ $`diag(1\mathrm{}1,1)`$, where $`a,b,\mathrm{}=1,\mathrm{},n`$. We denote the connection components (gauge fields) associated with the Lorentz, translation, co-translation, and dilation generators of the conformal group $`O(n,2)`$, $`n>2`$, as the spin-connection $`\omega _b^a`$, the solder-form $`\omega ^a`$, the co-solder-form $`\omega _a`$, and the Weyl vector $`\omega _0^0`$, respectively. The corresponding $`O(n,2)`$ curvatures $`𝛀_{\stackrel{~}{b}}^{\stackrel{~}{a}}`$ ($`\stackrel{~}{a},\stackrel{~}{b},\mathrm{}=0,1,\mathrm{},n`$) are referred to as the curvature $`𝛀_b^a`$, torsion $`𝛀^a𝛀_0^a`$, co-torsion $`𝛀_a𝛀_a^0`$, and dilation $`𝛀_0^0`$, respectively, and are defined by the biconformal structure equations,
$`𝛀_b^a`$ $`=`$ $`𝐝\omega _b^a\omega _b^c\omega _c^a\mathrm{\Delta }_{cb}^{ad}\omega _d\omega ^c`$ (1)
$`𝛀^a`$ $`=`$ $`𝐝\omega ^a\omega ^b\omega _b^a\omega _0^0\omega ^a`$ (2)
$`𝛀_a`$ $`=`$ $`𝐝\omega _a\omega _a^b\omega _b\omega _a\omega _0^0`$ (3)
$`𝛀_0^0`$ $`=`$ $`𝐝\omega _0^0\omega ^a\omega _a,`$ (4)
where $`\mathrm{\Delta }_{cd}^{ab}\delta _c^a\delta _d^b\eta ^{ab}\eta _{cd}`$. In all cases differential forms are bold and the wedge product is assumed between adjacent forms. The position of any lower-case Latin index corresponds to the associated conformal weight: each upper index contributes $`+1`$ to the weight, while each lower index contributes $`1.`$
Biconformal space is the $`2n`$-dimensional base space of the $`O(n,2)`$ principal bundle with homothetic fiber, first constructed in . Each biconformal curvature may be expanded in the $`(\omega ^a,\omega _b)`$ basis as
$$𝛀_{\stackrel{~}{b}}^{\stackrel{~}{a}}=\frac{1}{2}\mathrm{\Omega }_{\stackrel{~}{b}cd}^{\stackrel{~}{a}}\omega ^{cd}+\mathrm{\Omega }_{\stackrel{~}{b}d}^{\stackrel{~}{a}c}\omega _c\omega ^d+\frac{1}{2}\mathrm{\Omega }_{\stackrel{~}{b}}^{\stackrel{~}{a}cd}\omega _{cd},$$
(5)
where we use the convention of writing
$$\omega ^{ab\mathrm{}c}\omega ^a\omega ^b\mathrm{}\omega ^c=\omega ^a\omega ^b\mathrm{}\omega ^c.$$
The three terms of eq.(5) will be called the spacetime-, cross-, and momentum-term, respectively, of the corresponding curvature.
Any $`r`$-form $`𝐔`$ ($`r2n`$) defined on the cotangent bundle to biconformal space can be uniquely decomposed into a sum of $`(p,q)`$-forms,
$$𝐔=\underset{p=0}{\overset{r}{}}𝐔_{p,rp},$$
each of which is of the form
$$𝐔_{p,q}=\frac{1}{p!q!}U_{a_1\mathrm{}a_p}^{b_1\mathrm{}b_q}\omega ^{a_1\mathrm{}a_p}\omega _{b_1\mathrm{}b_q}(p,qn,p+q=r)$$
(6)
and has conformal weight $`pq`$. For example, a $`1`$-form can be written as
$$𝐔=U_a\omega ^a+U^a\omega _aU_A\omega ^A,$$
(7)
where capital Latin indices denote both upper and lower lower-case Latin indices.
Biconformal space possesses a natural metric,
$$K^{AB}=\left(\begin{array}{cc}0& \delta _b^a\\ \delta _b^a& 0\end{array}\right)$$
which is obtained when the non-degenerate Killing form of the conformal group $`O(n,2)`$ is restricted to the biconformal base space. The Killing metric defines a natural inner product between two $`1`$-forms $`𝐔`$ and $`𝐕`$:
$$𝐔,𝐕\frac{1}{2}K^{AB}U_AV_B=\frac{1}{2}\left(U^aV_a+V^aU_a\right)$$
(8)
Notice that, because the metric is essentially $`\delta _b^a`$, whenever we sum an upper with a lower index we have implicitly used the Killing metric.
In , we demonstrated that when the indices of the $`2n`$-dimensional Levi-Civita symbol are sorted by weight, it may be written as the product of two $`n`$-dimensional Levi-Civita symbols of opposite weights:
$$\epsilon _{a_1\mathrm{}a_n}^{b_1\mathrm{}b_n}=\epsilon _{a_1\mathrm{}a_n}^{}\epsilon ^{b_1\mathrm{}b_n},$$
where the mixed index positioning indicates the scaling weight of the indices, and not any use of the metric. The Levi-Civita tensor is normalized such that traces are given by
$$\epsilon _{a_1\mathrm{}a_pc_{p+1}\mathrm{}c_n}\epsilon ^{b_1\mathrm{}b_pc_{p+1}\mathrm{}c_n}=p!(np)!\delta _{a_1\mathrm{}a_p}^{b_1\mathrm{}b_p},$$
where the antisymmetric $`\delta `$-symbol is defined as
$$\delta _{a_1\mathrm{}a_p}^{b_1\mathrm{}b_p}\delta _{a_1}^{[b_1}\mathrm{}\delta _{a_p}^{b_p]}.$$
Using the Levi-Civita tensor, the scale-invariant volume form of biconformal space is given by:
$$𝚽=\epsilon _{a_1\mathrm{}a_n}^{b_1\mathrm{}b_n}\omega ^{a_1\mathrm{}a_n}\omega _{b_1\mathrm{}b_n}.$$
Notice that despite the mixed index positions, $`\epsilon _{a_1\mathrm{}a_n}^{b_1\mathrm{}b_n}`$ is totally antisymmetric on all $`2n`$ indices.
We now define the biconformal dual of a general $`r`$-form. Let $`𝐔_{p,q}`$ be an arbitrary $`(p,q)`$-form as denoted in (6). Then the dual of $`𝐔_{p,q}`$ is an $`(nq,np)`$-form, also of weight $`pq,`$ defined as
$${}_{}{}^{}𝐔_{p,q}^{}\frac{\tau (p,q)}{p!q!(np)!(nq)!}U_{a_1\mathrm{}a_p}^{b_1\mathrm{}b_q}\epsilon _{b_1\mathrm{}b_n}^{a_1\mathrm{}a_n}\omega ^{b_{q+1}\mathrm{}b_n}\omega _{a_{p+1}\mathrm{}a_n},$$
where
$$\tau (p,q)=\{\begin{array}{c}1ifpq\\ (1)^{n(p+q)}ifp<q\end{array}$$
This choice of $`\tau (p,q)`$ guarantees that for any $`r`$-form $`𝐔,`$
$${}_{}{}^{}𝐔=(1)^{r(nr)}𝐔$$
and that for two arbitrary $`(p,q)`$\- and $`(q,p)`$-forms, $`𝐔_{p,q}`$ and $`𝐕_{q,p}`$, respectively,
$$𝐔_{p,q}^{}𝐕_{q,p}=𝐕_{q,p}^{}𝐔_{p,q},$$
where we again assume wedge products between forms. It is then easy to show that for any two $`r`$-forms $`𝐔`$ and $`𝐕`$, the product $`𝐔^{}𝐕`$ is proportional to the volume form $`𝚽`$ and
$$𝐔^{}𝐕=\underset{p=0}{\overset{r}{}}𝐔_{p,rp}^{}𝐕_{rp,p}=\underset{p=0}{\overset{r}{}}𝐕_{p,rp}^{}𝐔_{rp,p}=𝐕^{}𝐔,$$
because $`𝐔_{p,rp}^{}𝐕_{rq,q}`$ vanishes unless $`p=q`$.
Now let $`𝐔`$ be a general $`1`$-form as in (7). Then
$${}_{}{}^{}𝐔=(1)^{n1}𝐔$$
and
$$\frac{1}{2}𝐔^{}𝐔=\frac{(1)^n}{n!^2}U_aU^a𝚽.$$
(9)
Thus, by eq.(8), the term $`𝐔^{}𝐔`$ is proportional to the inner product $`𝐔,𝐔`$:
$$𝐔^{}𝐔=\frac{2(1)^n}{n!^2}𝐔,𝐔𝚽=\frac{(1)^n}{n!^2}K^{AB}U_AU_B𝚽.$$
(10)
We are now ready to build the biconformal theory of scalar matter. Let $`\varphi ^m`$ be a set of massless Lorentz-scalar fields of conformal weight $`m𝐙`$ and $`𝐃\varphi ^m`$ be their biconformally covariant derivatives defined by
$$𝐃\varphi ^m𝐝\varphi ^m+m\omega _0^0\varphi ^m.$$
(11)
This covariant derivative is a $`1`$-form which is also of weight $`m`$ and can be expanded as
$$𝐃\varphi ^m=\omega ^AD_A\varphi ^m.$$
Since the dual operator preserves the conformal weight, $`{}_{}{}^{}𝐃\varphi ^m`$ must be of weight $`m`$, so that every term in the infinite sum
$$\underset{m𝐙}{}𝐃\varphi ^m𝐃\varphi ^m$$
is of conformal weight zero. Using (10) we see that
$$𝐃\varphi ^m𝐃\varphi ^m=\frac{(1)^n}{n!^2}K^{AB}D_A\varphi ^mD_B\varphi ^m𝚽$$
for every $`m`$. We have thus arrived at the Weyl-scalar-valued action
$$S_M=\frac{1}{2}\overline{\lambda }\underset{m}{}𝐃\varphi ^m𝐃\varphi ^m=\frac{1}{2}\overline{\lambda }\frac{(1)^n}{n!^2}\underset{m}{}K^{AB}D_A\varphi ^mD_B\varphi ^m𝚽.$$
(12)
We shall make use of the ‘dual’ form of $`S_M`$ when we vary the action with respect to the field and with respect to the connection, whereas the form of $`S_M`$ that explicitly displays the dependence on the Killing metric proves more useful in varying the base forms.
## 3 The linear scalar action
In a $`2n`$-dimensional biconformal space the most general Lorentz and scale-invariant action which is linear in the biconformal curvatures and structural invariants is
$$S_G=(\alpha 𝛀_{b_1}^{a_1}+\beta \delta _{b_1}^{a_1}𝛀_0^0+\gamma \omega ^{a_1}\omega _{b_1})\omega ^{a_2\mathrm{}a_n}\omega _{b_2\mathrm{}b_n}\epsilon ^{b_1\mathrm{}b_n}\epsilon _{a_1\mathrm{}a_n},$$
first introduced in . We will always assume non-vanishing $`\alpha ,`$ $`\beta ,`$ and $`\gamma `$. For a set of massless Lorentz scalar fields $`\varphi ^m`$ of weight $`m`$ with kinetic term $`S_M`$ given by (12), we have
$$S=S_M+S_G.$$
(13)
Variation of this action with respect to the scalar fields yields the equation
$$0=𝐃^{}𝐃\varphi ^m$$
(14)
for every $`m`$, where
$$𝐃^{}𝐃\varphi ^m𝐝^{}𝐃\varphi ^m+m\omega _0^0𝐃\varphi ^m.$$
Variation with respect to the connection one-forms gives rise to the following field equations:
$`\beta (\mathrm{\Omega }^a{}_{ba}{}^{}2\mathrm{\Omega }_{ca}^d\delta _{db}^{ca})`$ $`=`$ $`\lambda \mathrm{\Theta }_b`$ (15)
$`\beta (\mathrm{\Omega }_a^{ba}2\mathrm{\Omega }_a^{cd}\delta _{dc}^{ab})`$ $`=`$ $`\lambda \mathrm{\Theta }^b`$ (16)
$`\alpha (\mathrm{\Delta }_{eg}^{af}\mathrm{\Omega }_{ab}^b+2\mathrm{\Delta }_{eb}^{cf}\delta _{dg}^{ab}\mathrm{\Omega }_{ac}^d)`$ $`=`$ $`0`$ (17)
$`\alpha (\mathrm{\Delta }_{eb}^{gf}\mathrm{\Omega }_a^{ab}+2\mathrm{\Delta }_{ed}^{af}\delta _{ab}^{gc}\mathrm{\Omega }_c^{bd})`$ $`=`$ $`0`$ (18)
$`\alpha \mathrm{\Omega }_{bac}^a+\beta \mathrm{\Omega }_{0bc}^0`$ $`=`$ $`\lambda \mathrm{{\rm Y}}_{bc}`$ (19)
$`2(\alpha \mathrm{\Omega }_{cd}^{ec}+\beta \mathrm{\Omega }_{0d}^{0e})\delta _{eb}^{ad}+\mathrm{\Lambda }_b^a`$ $`=`$ $`\lambda \mathrm{{\rm Y}}_b^a`$ (20)
$`\alpha \mathrm{\Omega }_a^{bac}+\beta \mathrm{\Omega }_0^{0bc}`$ $`=`$ $`\lambda \mathrm{{\rm Y}}^{bc}`$ (21)
$`2(\alpha \mathrm{\Omega }_{dc}^{ce}+\beta \mathrm{\Omega }_{0d}^{0e})\delta _{eb}^{ad}+\mathrm{\Lambda }_b^a`$ $`=`$ $`\lambda \mathrm{{\rm Y}}_b^a`$ (22)
where the matter sources are given by
$`\mathrm{{\rm Y}}_{ab}`$ $``$ $`{\displaystyle \underset{m}{}}D_a\varphi ^mD_b\varphi ^m=\mathrm{{\rm Y}}_{ba}`$
$`\mathrm{{\rm Y}}_b^a`$ $``$ $`{\displaystyle \underset{m}{}}\left(D^a\varphi ^mD_b\varphi ^m+D^c\varphi ^mD_c\varphi ^m\delta _b^a\right)`$
$`\mathrm{{\rm Y}}^{ab}`$ $``$ $`{\displaystyle \underset{m}{}}D^a\varphi ^mD^b\varphi ^m=\mathrm{{\rm Y}}^{ba}`$
$`\mathrm{\Theta }_b`$ $`=`$ $`{\displaystyle \underset{m}{}}m\varphi ^mD_b\varphi ^m`$
$`\mathrm{\Theta }^b`$ $`=`$ $`{\displaystyle \underset{m}{}}m\varphi ^mD^b\varphi ^m`$
and we have defined
$`\mathrm{\Lambda }_b^a`$ $``$ $`(\alpha (n1)\beta +\gamma n^2)\delta _b^a`$
$`\lambda `$ $``$ $`{\displaystyle \frac{1}{(n1)!^2}}\overline{\lambda }.`$
Note that since the spin connection does not occur in the covariant derivative (11), $`\delta _{\omega _b^a}S_M0`$, and there is no matter contribution to eq.(17) or (18). Combining equations (20) and (22) we see that the latter can be replaced by
$$\mathrm{\Omega }_{cd}^{ac}=\mathrm{\Omega }_{dc}^{ca}.$$
(23)
We remark that the biconformal structure equations together with eqs.(15)-(18) may be used to express the torsion and co-torsion in terms of the connection, the Weyl vector, and (here) the matter fields $`\varphi ^m`$. In it was observed that constraining the torsion to vanish also forces the Weyl vector to vanish, but this conclusion no longer holds with matter present. This suggests that setting the torsion to zero is not an undue constraint as assumed in , but rather, that the Weyl vector vanishes unless there are appropriate matter fields present. Taking this view, we are free to assume $`𝛀^a=0.`$ Then, there exists a gauge in which the Weyl vector is given in terms of covariant derivatives of the fields $`\varphi ^m`$ by
$$\omega _0^0=\frac{1}{(n1)(n2)}\frac{\lambda }{\beta }\underset{m}{}m\varphi ^m𝐃\varphi ^m,$$
Notice that $`\omega _0^0=0`$ unless conjugate weights, $`+m`$ and $`m`$ are both present.
We will explore such dilational sources further elsewhere. Here, since our goal is to derive the usual form of the Einstein equations with scalar matter, it is sufficient to restrict our attention to the case $`m=0.`$ Thus, for the remainder of this paper we restrict to the case of a scalar field $`\varphi `$ of conformal weight zero. Then, the covariant derivative is simply the exterior derivative
$$𝐃\varphi 𝐝\varphi =\omega ^Ad_A\varphi ,$$
so the field equations reduce to
$`0`$ $`=`$ $`{}_{}{}^{}𝐝_{}^{}𝐝\varphi `$ (24)
$`0`$ $`=`$ $`\mathrm{\Omega }^a{}_{ba}{}^{}2\mathrm{\Omega }_{ca}^d\delta _{db}^{ca}`$ (25)
$`0`$ $`=`$ $`\mathrm{\Omega }_a^{ba}2\mathrm{\Omega }_a^{cd}\delta _{dc}^{ab}`$ (26)
together with eqs.(17)-(23), where $`\mathrm{\Theta }_b=0`$, $`\mathrm{\Theta }^b=0`$ and
$`\mathrm{{\rm Y}}_{ab}`$ $``$ $`d_a\varphi d_b\varphi =\mathrm{{\rm Y}}_{ba}`$
$`\mathrm{{\rm Y}}_b^a`$ $``$ $`d^a\varphi d_b\varphi +d^c\varphi d_c\varphi \delta _b^a`$
$`\mathrm{{\rm Y}}^{ab}`$ $``$ $`d^a\varphi d^b\varphi =\mathrm{{\rm Y}}^{ba}.`$
## 4 Solution for the curvatures
We now find the most general solution to these equations subject only to the constraint of vanishing torsion. As discussed above, this condition no longer implies a vanishing Weyl vector. The Weyl vector nonetheless vanishes because of our choice to consider only zero conformal weight matter.
Despite vanishing torsion, the general approach to solving the field equations follows that of . Starting with a general ansatz for the spin connection and Weyl vector, we first solve the torsion and co-torsion equations, eqs.(25, 26) and eqs.(17, 18). The Bianchi identity following from the vanishing torsion constraint and field equations (19)-(23) then determines the form of the curvature and dilation. In Sec.(5), we show that the vanishing torsion constraint also leads to a foliation by $`n`$-dimensional flat Riemannian manifolds. By invoking the gauge freedom on each of these manifolds, we show the existence of a second foliation by $`n`$-dimensional Riemannian spacetimes satisfying the Einstein equations with scalar matter.
To begin, we write the spin connection $`\omega _b^a`$ as
$$\omega _b^a=\alpha _b^a+\beta _b^a+\gamma _b^a$$
(27)
with $`\alpha _b^a`$ and $`\beta _b^a`$ defined by
$`𝐝\omega ^a`$ $`=`$ $`\omega ^b\alpha _b^a+\frac{1}{2}\mathrm{\Omega }^{abc}\omega _{bc}`$
$`𝐝\omega _a`$ $`=`$ $`\beta _a^b\omega _b+\frac{1}{2}\mathrm{\Omega }_{abc}\omega ^{bc}.`$
Using this ansatz in structure equations (2) and (3), $`\mathrm{\Omega }^{abc}`$ and $`\mathrm{\Omega }_{abc}`$ remain related to derivatives of the solder- and co-solder forms, whereas the other torsion and co-torsion terms are algebraic in the components of $`\alpha _b^a,\beta _b^a`$, $`\gamma _b^a`$, and $`\omega _0^0`$. Thus, the separation of the connection allows us to solve the torsion/co-torsion field equations (25), (26), (17), and (18) algebraically.
We simply state the result of this reduction here. More detail is available in . Defining
$$\sigma _b^a\alpha _b^a\beta _b^a\sigma _{bc}^a\omega ^c+\sigma _b^{ac}\omega _c$$
and setting
$$𝛀^a=0,$$
field equations (25), (26), (17), and (18) imply
$`\omega _0^0`$ $`=`$ $`0`$
$`\sigma _{bc}^a`$ $`=`$ $`0`$
$`\sigma _a^{ba}`$ $`=`$ $`0`$
with no assumption concerning the co-torsion, curvature, or dilation. From these we find
$$\omega _b^a=\alpha _b^a.$$
The co-torsion cross- and momentum terms reduce to
$`\mathrm{\Omega }_{ac}^b`$ $`=`$ $`0`$
$`\mathrm{\Omega }_a^{bc}`$ $`=`$ $`\sigma _a^{bc}\sigma _a^{cb},`$
so that the full co-torsion is
$$𝛀_a=\frac{1}{2}\mathrm{\Omega }_{acd}\omega ^{cd}+\sigma _a^{bc}\omega _{bc}$$
with
$$\sigma _a^{ba}=0.$$
The spacetime co-torsion $`\mathrm{\Omega }_{abc}`$ remains undetermined.
Next, we turn our attention to the curvature and dilation equations, eqs.(19)-(23). The vanishing torsion constraint makes it possible to obtain an algebraic condition on the curvatures from the Bianchi identity associated with eq.(2). Taking the exterior derivative of eq.(2) gives
$$\omega ^b𝛀_b^a=\omega ^a𝛀_0^0,$$
(28)
which implies for the curvature components
$`\mathrm{\Omega }_0^{0cd}`$ $`=`$ $`0`$ (29)
$`\mathrm{\Omega }_b^{acd}`$ $`=`$ $`0`$ (30)
$`\mathrm{\Omega }_{0d}^{0c}`$ $`=`$ $`\frac{1}{n1}\mathrm{\Omega }_{da}^{ac}`$ (31)
$`\mathrm{\Omega }_{cd}^{ab}`$ $`=`$ $`\frac{1}{n1}\mathrm{\Delta }_{cd}^{fa}\mathrm{\Omega }_{fe}^{eb}=\mathrm{\Delta }_{cd}^{fa}\mathrm{\Omega }_{0f}^{0b}`$ (32)
$`\mathrm{\Omega }_{0cd}^0`$ $`=`$ $`\frac{1}{n2}(\mathrm{\Omega }_{dac}^a\mathrm{\Omega }_{cad}^a)`$ (33)
Next, we impose field equations (19)-(23) onto these conditions and see that eq. (21) is satisfied by virtue of eqs.(29) and (30) if and only if
$$\mathrm{{\rm Y}}^{ab}=0,$$
so that
$`d^a\varphi `$ $`=`$ $`0`$
$`\mathrm{{\rm Y}}_b^a`$ $`=`$ $`0.`$
Imposing eqs.(20) and (22) onto (31) now yields for $`\beta (n1)\alpha `$
$$\mathrm{\Omega }_{0b}^{0a}=\frac{1}{\alpha (n1)\beta }(\lambda (d^a\varphi d_b\varphi +\frac{2}{n1}d^c\varphi d_c\varphi \delta _b^a)\frac{1}{n1}(\alpha (n1)\beta +\gamma n^2)\delta _b^a$$
Since $`d^a\varphi =0,`$ this reduces to
$$\mathrm{\Omega }_{0b}^{0a}=\frac{(\alpha (n1)\beta +\gamma n^2)}{(n1)(\alpha (n1)\beta )}\delta _b^a\chi \delta _b^a.$$
(34)
Imposing eq.(19) onto (33) yields
$$\begin{array}{c}\mathrm{\Omega }_{0bc}^0=0\\ \mathrm{\Omega }_{bac}^a=\mathrm{\Omega }_{cab}^a\\ \alpha \mathrm{\Omega }_{bac}^a=\lambda \mathrm{{\rm Y}}_{bc}\end{array}if\beta \frac{1}{2}(n2)\alpha $$
If the first of these conditions is substituted back into (28), we obtain the cyclic identity on the spacetime curvature:
$$\mathrm{\Omega }_{[bcd]}^a=0.$$
Summarizing the forms of the dilation and curvature so far we have for the generic case
$`𝛀_0^0`$ $`=`$ $`\chi \omega _a\omega ^a`$ (35)
$`𝛀_b^a`$ $`=`$ $`\frac{1}{2}\mathrm{\Omega }_{bcd}^a\omega ^{cd}+\chi \mathrm{\Delta }_{bd}^{ca}\omega _c\omega ^d,`$ (36)
with
$`\mathrm{\Omega }_{[bcd]}^a`$ $`=`$ $`0`$ (37)
$`\alpha \mathrm{\Omega }_{bac}^a`$ $`=`$ $`\lambda \mathrm{{\rm Y}}_{bc}.`$ (38)
We now use the vanishing of the Weyl vector to obtain further constraints on $`𝛀_0^0`$ and $`𝛀_b^a`$. Substituting the restricted form of the dilation (35) into structure equation (4),
$$𝛀_0^0=\mathrm{\Omega }_{0b}^{0a}\omega _a\omega ^b=\omega _a\omega ^a,$$
we see that
$$\mathrm{\Omega }_{0b}^{0a}=\delta _b^a.$$
(39)
The last equation has to be equal to the dilation crossterm as given by eq.(34), which implies $`\chi =1`$, i.e. a relationship between the constants $`\alpha `$, $`\beta `$, and $`\gamma `$:
$$\frac{\gamma n}{\alpha (n1)\beta }=1.$$
Hence, the volume ($`\gamma `$) term must necessarily be present in the action. Thus, in the generic case, where
$$\beta (n1)\alpha ;\beta \frac{1}{2}(n2)\alpha ,$$
we have a two-parameter class of allowed actions, differing only in the constant $`\frac{\lambda }{\alpha }`$. Through eq.(32), eq.(39) also implies that
$$\mathrm{\Omega }_{bd}^{ac}=\mathrm{\Delta }_{db}^{ac}.$$
We have now satisfied all of field equations (15)-(23). The curvatures take the form
$`𝛀^a`$ $`=`$ $`0`$ (40)
$`𝛀_a`$ $`=`$ $`\sigma _a^{bc}\omega _{bc}+\frac{1}{2}\mathrm{\Omega }_{abc}\omega ^{bc}`$ (41)
$`𝛀_0^0`$ $`=`$ $`\omega _a\omega ^a`$ (42)
$`𝛀_b^a`$ $`=`$ $`\frac{1}{2}\mathrm{\Omega }_{bcd}^a\omega ^{cd}\mathrm{\Delta }_{bd}^{ca}\omega _c\omega ^d`$ (43)
subject to the constraints
$`\mathrm{\Omega }_{[bcd]}^a`$ $`=`$ $`0`$
$`\alpha \mathrm{\Omega }_{bac}^a`$ $`=`$ $`\lambda \mathrm{{\rm Y}}_{bc}`$
$`\sigma _a^{ba}`$ $`=`$ $`0`$
$`d^a\varphi `$ $`=`$ $`0`$
Notice that the dilation is necessarily non-degenerate, but may not be closed.
In the next section, we find further constraints on the curvatures arising from the structure equations.
## 5 Solution for the connection
While eqs.(40)-(43) for the curvatures satisfy all of the field equations, they do not fully incorporate the form of the biconformal structure equations as embodied in the Bianchi identities. Therefore, in this section, we turn to the consequences of the form of the curvatures on the connection.
So far, we have established that in the as yet unspecified original gauge the Weyl vector vanishes and the spin connection is fully determined by the solder form $`\omega ^a`$:
$`\omega _0^0`$ $`=`$ $`0`$
$`\omega _b^a`$ $`=`$ $`\alpha _b^a=\alpha _{bc}^a\omega ^c+\alpha _b^{ac}\omega _c.`$
Substituting the reduced curvatures into eqs.(1)-(3) (eq.(4) is identically satisfied by (42)), the structure equations now take the form
$`𝐝\alpha _b^a`$ $`=`$ $`\alpha _b^c\alpha _c^a+\frac{1}{2}\mathrm{\Omega }_{bcd}^a\omega ^{cd}`$ (44)
$`𝐝\omega ^a`$ $`=`$ $`\omega ^b\alpha _b^a`$ (45)
$`𝐝\omega _a`$ $`=`$ $`\alpha _a^b\omega _b+\sigma _a^{bc}\omega _{bc}+\frac{1}{2}\mathrm{\Omega }_{abc}\omega ^{bc}.`$ (46)
We observe that eq.(45) is in involution. By the Frobenius theorem, we can consistently set $`\omega ^a`$ to zero and obtain a foliation by submanifolds, where the spin connection and the co-solder form reduce to
$`𝐟_a`$ $``$ $`\omega _a|_{\omega ^a=0}`$
$`\widehat{\alpha }_b^a`$ $``$ $`\alpha _b^a|_{\omega ^a=0}=\alpha _b^{ac}𝐟_c`$
$`\widehat{\sigma }_a^{bc}`$ $``$ $`\sigma _a^{bc}|_{\omega ^a=0}.`$
Then each submanifold is described by the reduced structure equations
$`𝐝\widehat{\alpha }_b^a`$ $`=`$ $`\widehat{\alpha }_b^c\widehat{\alpha }_c^a`$
$`\mathrm{𝐝𝐟}_a`$ $`=`$ $`\widehat{\alpha }_a^b𝐟_b+\widehat{\sigma }_a^{bc}𝐟_{bc}.`$
Since the spin-connection is involute, there exists a Lorentz gauge transformation such that $`\widehat{\alpha }_b^a=0`$ on each submanifold, i.e. $`\alpha _b^{ac}=0`$. With this gauge choice the system reduces to simply
$$\mathrm{𝐝𝐟}_a=\widehat{\sigma }_a^{bc}𝐟_{bc}.$$
(47)
This can be solved in the usual way giving $`\widehat{\sigma }_a^{bc}`$ in terms of $`𝐟_a`$ and $`\mathrm{𝐝𝐟}_a`$. Since this solution has the same form on each leaf of the foliation, the expression for $`\sigma _a^{bc}`$ remains valid when it is extended back to the full space, i.e. $`\sigma _a^{bc}`$ depends on the $`2n`$ biconformal coordinates only through its dependence on $`𝐟_a`$.
The existence of an $`\widehat{\alpha }_b^a=0`$ gauge depends only on vanishing torsion, which leads to an involution of the solder form $`\omega ^a`$ and the resulting Bianchi identity, eq.(28). Therefore, the results of Sec.(4) remain valid in the $`\widehat{\alpha }_b^a=0`$ gauge.
Returning to the full biconformal space, we now have a gauge
$$\omega ^a𝐞^a$$
such that the spin connection is
$$\alpha _b^a=\alpha _{bc}^a𝐞^c,$$
while the co-solder form may be written in terms of $`𝐟_a`$ and an additional term linear in the solder form,
$$\omega _a=𝐟_a+h_{ab}𝐞^b,$$
(48)
Notice that while $`𝐟_a`$ depends on all $`2n`$ coordinates of this extension, it remains independent of the $`1`$-forms $`𝐞^a`$. This means that $`\mathrm{𝐝𝐟}_a`$ remains at least linear in $`𝐟_a,`$ and is consequently involute. We can therefore turn the problem around, setting $`𝐟_a=0`$ to obtain a second foliation of the biconformal space. We can define $`𝐡_a`$ in terms of this involution, setting
$$𝐡_a\omega _a|_{𝐟_a=0}=h_{ab}𝐞^b,$$
with $`h_{ab}`$ arbitrary. The structure equations for the $`𝐟_a=0`$ geometry are
$`𝐝\alpha _b^a`$ $`=`$ $`\alpha _b^c\alpha _c^a+\frac{1}{2}\mathrm{\Omega }_{bcd}^a𝐞^{cd}`$ (49)
$`\mathrm{𝐝𝐞}^a`$ $`=`$ $`𝐞^b\alpha _b^a`$ (50)
$`\mathrm{𝐝𝐡}_a`$ $`=`$ $`\alpha _{ac}^b𝐞^c𝐡_b+\sigma _a^{bc}𝐡_{bc}+\frac{1}{2}\mathrm{\Omega }_{abc}𝐞^{bc}`$ (51)
Eq.(51) determines $`𝐡_a`$ once the spacetime co-torsion, $`\mathrm{\Omega }_{abc},`$ is given, with little consequence for the rest of the geometry. We focus our attention on the first two equations. Since they are unchanged from their full biconformal form, the curvature
$$𝐑_b^a\frac{1}{2}\mathrm{\Omega }_{bcd}^a𝐞^{cd}$$
and connection $`\alpha _b^a`$ (and of course $`𝐞^a`$, by the first involution) are fully determined on the $`n`$-dimensional $`𝐟_a=0`$ submanifold. Thus, $`\alpha _b^a`$ is the usual spin connection compatible with $`𝐞^a`$, while $`𝐑_b^a`$ is its curvature. If we let
$$R_{ab}\mathrm{\Omega }_{acb}^c,$$
then the Bianchi identity following from (49),
$$\mathrm{𝐃𝐑}_b^a=0,$$
implies that the tensor
$$G_{ab}R_{ab}\frac{1}{2}\eta _{ab}R$$
is divergence-free.
Now that $`\mathrm{\Omega }_{acb}^c`$ is seen to be the Ricci tensor of an underlying $`n`$-dim submanifold, it follows from the remaining condition (38),
$`R_{ab}`$ $`=`$ $`\kappa \mathrm{{\rm Y}}_{ab}`$
$`\kappa `$ $``$ $`\frac{\lambda }{\alpha }`$
that these submanifolds satisfy the Einstein equations,
$$G_{ab}=\kappa T_{ab},$$
(52)
with the divergence-free stress-energy tensor given by derivatives of the matter field:
$$T_{ab}d_a\varphi d_b\varphi \frac{1}{2}\eta _{ab}\eta ^{ce}d_c\varphi d_e\varphi .$$
(53)
Even though the co-torsion has a nonvanishing spacetime projection, the curvature is the one computed from the solder form $`\omega ^a`$ alone. This is our principal result, establishing a direct connection between general relativity with scalar matter and the more general structure of biconformal gauge theory with scalar matter.
## 6 Constraints on the matter field
We have seen that the Bianchi identity associated with eq.(2) with vanishing torsion together with field equation (21) imply
$$d^a\varphi =0$$
or
$$𝐝\varphi =\omega ^ad_a\varphi .$$
(54)
We now show that this implies that the scalar field $`\varphi `$ depends only on the $`n`$ coordinates spanning each leaf of the $`𝐟_a=0`$ foliation and identically satisfies its own field equation.
Based on the involution for $`\omega ^a`$ there exist $`n`$ coordinates $`x^\mu `$ of weight $`+1`$ such that
$$\omega ^a=e_\mu ^a𝐝x^\mu $$
with the component matrices $`e_\mu ^a`$ necessarily invertible. From eq.(45), we immediately find that $`e_\mu ^a=e_\mu ^a(x).`$ Similarly, it can be shown from the $`𝐟_a`$ involution that there exist $`n`$ complementary coordinates $`y_\nu `$ of weight $`1`$ such that $`𝐟_a`$ takes the form
$$𝐟_a=f_a^\mu 𝐝y_\mu .$$
Since both $`𝐟_a`$ and $`𝐝y_\mu `$ completely span the co-tangent bundles of the $`\omega ^a=0`$ submanifolds, the component matrices $`f_a^\mu `$ are also necessarily invertible. Thus, $`(x^\mu ,y_\mu )`$ form a complete set of local coordinates on biconformal space. If we expand $`𝐝\varphi `$ in this coordinate basis, we obtain
$$𝐝\varphi =_\mu \varphi 𝐝x^\mu +^\mu \varphi 𝐝y_\mu ,$$
where $`(_\mu ,^\mu )`$ denote derivatives with respect to $`(x^\mu ,y_\nu )`$. Equating this general form with the derived form (54),
$$𝐝\varphi =e_\mu ^a(x)d_a\varphi 𝐝x^\mu ,$$
we see that $`\varphi =\varphi (x)`$. Hence, the scalar field is entirely defined on the $`n`$-dimensional Riemannian spacetime.
Using eq.(48) to expand the co-solder-form as
$$\omega _a=f_a^\mu 𝐝y_\mu +h_{ab}e_\mu ^b𝐝x^\mu ,$$
it is now easy to show that
$$d^a\varphi =0d^ad_a\varphi =0.$$
This result, together with the vanishing torsion constraint and eqs.(45)-(46), imply that the field equation for $`\varphi `$, eq.(24), is identically satisfied.
However, the field is constrained by the fact that the stress-energy tensor (53) is, by eq.(52), proportional to the divergence-free Einstein tensor and hence must be itself divergence-free. Since
$`\eta ^{be}D_eT_{ab}`$ $`=`$ $`\eta ^{be}D_e\left(D_a\varphi D_b\varphi \frac{1}{2}\eta _{ab}\eta ^{cf}D_c\varphi D_f\varphi \right)`$
$`=`$ $`\eta ^{cd}(D_cD_d\varphi )D_a\varphi `$
the stress-energy tensor is divergence-free if and only if
$$\eta ^{cd}D_cD_d\varphi =0.$$
This establishes that the scalar field $`\varphi `$ satisfies the free Klein-Gordon (wave) equation on the $`n`$-dimensional spacetime.
## 7 Conclusions
We have developed aspects of the theory of scalar matter in biconformal space. Using the existence of an inner product of $`1`$-forms and a dual operator, we constructed an action for a scalar matter field $`\varphi ^m`$ coupled to gravity and found the field equations. We solved them for the case of a scalar field of conformal weight zero in a torsion-free biconformal geometry. As in the vacuum case, the generic solutions are foliated by equivalent $`n`$-dimensional Riemannian spacetime manifolds. The curvature of each submanifold satisfies the usual Einstein equations with scalar matter. The scalar field is entirely defined on the submanifold and satisfies the $`n`$-dimensional massless Klein-Gordon equation. |
warning/0001/hep-ph0001300.html | ar5iv | text | # On Electroweak Baryogenesis in Gauge Mediated Models with Messenger-Matter Mixing
## 1 Introduction
Among various puzzles of Nature the problem of baryogenesis is one of the most intriguing. Since pioneering works , necessary conditions for the generation of the baryon asymmetry are well-known. These are baryon number violation, C- and CP-nonconservation and departure from thermal equilibrium at a certain stage of the evolution of the Universe. Search for reliable mechanisms of baryogenesis is limited to theories where these conditions are met.
One of the most appealing mechanisms is electroweak baryogenesis . It was proposed that this mechanism works during the electroweak phase transition. Unlike other mechanisms, electroweak baryogenesis was thought to be inherent in the Standard Model (SM) and seemed to require no additional fields. However, the baryon asymmetry tends to be washed out by anomalous electroweak processes if the latter come into thermal equilibrium after the phase transition completes. To prevent this wash-out, the electroweak phase transition has to be of the (strong enough) first order; this imposes a constraint on the expectation value of Higgs field $`\upsilon (T)`$ (with $`\upsilon (0)245`$ GeV) at the critical temperature $`T_c`$ ,
$$\frac{\upsilon (T_c)}{T_c}1.$$
(1)
In the case of SM this requirement implies the bound on the Higgs boson mass ,
$$m_h<40\mathrm{GeV},$$
(2)
which is inconsistent with present experimental limits . Hence electroweak baryogenesis fails in the Standard Model. <sup>1</sup><sup>1</sup>1There are recent attempts (see, e.g., Ref. ) to revive electroweak baryogenesis in SM by invoking the dynamics of preheating.
Similar situation takes place in the Minimal Supersymmetric Standard Model (MSSM). Although the corresponding bounds on the parameters of MSSM are significantly weaker, almost entire parameter space consistent with successful electroweak baryogenesis is excluded by existing experimental data. Still, electroweak baryogenesis works in MSSM with light Higgs, light $`\stackrel{~}{t}_R`$, heavy $`\stackrel{~}{t}_L`$ , small $`\stackrel{~}{t}_R\stackrel{~}{t}_L`$ mixing, and charginos with fairly degenerate masses (for brief reviews, see Refs. ). This spectrum was discussed in the framework of models with gravity mediation of supersymmetry breaking , and it was observed that the price for this “light stop” solution is non-universal boundary conditions for soft terms at GUT scale.
Recently, the phenomenology of Gauge Mediated Supersymmetry Breaking (GMSB) models attracted considerable attention (for a review, see Ref. ). In these models, supersymmetry breaking occurs in a separate sector. Supersymmetry breaking terms in the visible sector are generated due to special fields (messengers) charged under SM gauge group. The soft masses of superpartners are determined by their quantum numbers and are proportional to the corresponding gauge coupling constants. Consequently, all squarks are very heavy and electroweak baryogenesis seems not to work. So, it has been argued that the favorite mechanism to produce the baryon asymmetry in GMSB models is the Affleck–Dine baryogenesis .
This letter addresses the possibility of electroweak baryogenesis in GMSB with messenger-matter mixing. This mixing is natural since messengers carry the same quantum numbers as ordinary fields in MSSM . The constraints on mixing parameters imposed by the absence of lepton flavor violating processes and FCNC are not particularly strong . Moreover, it was pointed out that some scalar masses may be significantly reduced at large values of mixing terms without any contradiction to experimental limits . We will see that the small mass of stop required for successful electroweak baryogenesis may be explained in this way. We will show also that all other conditions on the spectrum may be satisfied.
## 2 Light stop window for electroweak baryogenesis
In any model of baryogenesis two main questions arise: i) What is the mechanism of the generation of the baryon asymmetry? ii) What protects the baryon asymmetry from being washed out by sphaleron processes?
In the framework of “light stop” solution, the answers to these questions are as follows.
i) The relevant source of CP-violation is the phase $`\varphi _{CP}`$ of $`\mu `$-term in Higgs superpotential. Baryogenesis is fueled by CP-asymmetry in chargino flux through expanding bubble wall. The realistic amount of baryon asymmetry $`n_B/s`$ is produced by sphaleron processes provided that chargino and neutralino masses are not much larger than the critical temperature $`T_c`$100 GeV. The resulting asymmetry is determined by the ratios of masses to the critical temperature, and degenerate $`\mu `$ and wino mass $`M_2`$ are favorable. Though the general tendencies are clear, the actual calculated values of $`n_B/s`$ depend on the approximations made. In what follows we make use of the constraints on the values of CP-phase and the level of degeneracy presented in Ref. . The smallest allowed value of $`\varphi _{CP}`$ is $`|\varphi _{CP}|=10^4`$; at this value one requires $`\mu =M_2`$. In the opposite case of large CP-violating phase, $`|\mathrm{sin}\varphi _{CP}|=1`$, the allowed region of masses is $`\mu =ϵM_2`$, $`0.4ϵ2.5`$ at $`\mu `$,$`M_2`$$`T_c`$ and $`0.5ϵ2`$ at $`\mu `$,$`M_2`$$`3T_c`$. The favorable region of the mass of CP-odd Higgs boson is $`m_A300`$ GeV, otherwise $`n_B/s`$ is suppressed by $`m_A^2`$ . For heavier CP-odd Higgs, the values of $`\mu `$ and $`M_2`$ are to be more degenerate in order that the baryon asymmetry at given $`\varphi _{CP}`$ be the same as at small $`m_A`$.
ii) As mentioned above, electroweak phase transition has to be of the first order, so that the constraint (1) is satisfied. Light stop helps in strengthening the phase transition. The largest baryon asymmetry is obtained in the context of MSSM with right stop plasma mass vanishing at the critical temperature, when the condition (1) becomes
$$1\frac{\upsilon (T_c)}{T_c}=\left(\frac{\upsilon (T_c)}{T_c}\right)_{SM}+\frac{2m_t^3\left(1\frac{\stackrel{~}{A}_t^2}{m_Q^2}\right)^{3/2}}{\pi \upsilon m_h^2},$$
(3)
where $`\stackrel{~}{A}_t=|A_t\mu \mathrm{cot}\beta |`$, and $`A_t`$ is the stop trilinear soft term, $`m_Q`$ is the left stop mass and $`m_t`$ is the on-shell running top quark mass in the $`\overline{MS}`$ scheme. Hereafter $`m_h`$ denotes the mass of the lightest Higgs boson. We consider the case $`m_AT_c`$, which is relevant to GMSB models; in this case CP-odd Higgs does not affect the electroweak phase transition. The first term on the right hand side of Eq. (3) is the ordinary Standard Model contribution,
$$\left(\frac{\upsilon (T)}{T}\right)_{SM}\left(\frac{40\mathrm{GeV}}{m_h}\right)^2.$$
The second term in Eq. (3) is the contribution from light right stop. Current limits on the lightest Higgs mass , $`m_h100`$ GeV, combined with the inequality (3) impose a constraint on left-right mixing
$$\stackrel{~}{A}_t/m_Q0.5.$$
(4)
At larger $`m_h`$ left-right mixing has to be smaller. At $`\stackrel{~}{A}_t=0`$ Eq. (3) gives the upper bound on the lightest Higgs boson mass in the theory with successful electroweak baryogenesis, $`m_h115`$ GeV . It is worth noting that the result (3) has been obtained by making use of improved one-loop effective potential. Higher order corrections make the phase transition slightly stronger (for a brief review and references, see Ref. ) and will be neglected in what follows.
Zero right stop plasma mass at the critical temperature implies
$$m_{\stackrel{~}{t}_R}^{2(eff)}=m_U^2+\mathrm{\Pi }_R(T_c)0,$$
(5)
where $`m_U`$ is the low energy value of the right stop soft mass term and $`\mathrm{\Pi }_R(T_c)`$ is the finite temperature contribution to the effective squared mass which is of order $`T_c^2`$ . Hence one needs negative $`m_U^2(100`$ GeV)<sup>2</sup>, which in principle may result in the existence of charge- and color-breaking (local) vacuum. A conservative requirement is that the physical vacuum has lower energy than the color-breaking minimum. The latter condition at, as an example, $`\stackrel{~}{A}_t=0`$ yields approximately
$$|m_U|<m_{crit}=\left(\frac{m_h^2\upsilon ^2\alpha _3\pi }{3}\right)^{1/4},$$
(6)
that is $`|m_U|95`$ GeV for $`m_h=100`$ GeV. One can relax this constraint by considering the physical vacuum as a metastable but long-living minimum. However, the inequality
$$m_U^2+\mathrm{\Pi }_R(T_c)>0$$
(7)
has to be satisfied in any case, otherwise the Universe would be driven to a charge- and color-breaking minimum at $`T>T_c`$ (for a discussion see Ref. ).
Let us collect the requirements which are imposed on the theory with “light stop” solution:
1. Right-left mixing in stop sector is small, $`\stackrel{~}{A}_t<0.5m_Q`$ at $`m_h>100`$ GeV, and the heavier the lightest Higgs boson, the smaller the mixing.
2. At $`m_h>115`$ GeV there is no window for electroweak baryogenesis. Consequently, the mass of the lightest Higgs boson belongs to the interval 100 GeV$`<m_h<`$115 GeV, where the lower bound comes from experiment.
3. The favorable interval of the mass of CP-odd Higgs boson is $`150\mathrm{GeV}m_A300\mathrm{GeV}`$. At smaller $`m_A`$ the phase transition is weaker, while at larger $`m_A`$ the baryon asymmetry is suppressed by $`m_A^2`$.
4. The soft mass squared of right stop is negative and of order $`T_c^2`$. Inequalities (6), (7) (or similar ones at non-zero $`\stackrel{~}{A}_t`$) have to be satisfied.
## 3 Gauge Mediated Models with messenger-matter mixing
Let us recall the spectrum of superpartners in GMSB models . In simple versions, messengers belong to complete vector-like GUT (e.g., $`SU(5)`$) multiplets. The induced soft terms depend crucially on $`n=n_5+3n_{10}`$, with $`n_5`$ and $`n_{10}`$ being the numbers of $`5+\overline{5}`$ and $`10+\overline{10}`$ messenger generations. Indeed, the spectrum of superpartners in these models at the messenger scale $`M_m`$ is
$$M_i=nc_i\frac{\alpha _i}{4\pi }\mathrm{\Lambda }f_1(\frac{\mathrm{\Lambda }}{M_m}),$$
(8)
$$m^2=2n\mathrm{\Lambda }^2\underset{i=1}{\overset{3}{}}C_ic_i\left(\frac{\alpha _i}{4\pi }\right)^2f_2(\frac{\mathrm{\Lambda }}{M_m}),$$
(9)
where $`M_i`$ denote gaugino masses and $`m`$ are soft masses of the scalar superpartners of fermions of the Standard Model. Here $`\alpha _1=\alpha /\mathrm{cos}^2\theta _W`$, $`\alpha _2=\alpha /\mathrm{sin}^2\theta _W`$, $`\alpha _3`$ are gauge coupling constants of electroweak and strong interactions; $`c_1=5/3`$, $`c_2=c_3=1`$; $`C_3=4/3`$ for color triplets (zero for singlets), $`C_2=3/4`$ for weak doublets (zero for singlets), $`C_1=\left(\frac{Y}{2}\right)^2`$, where $`Y=2(Q_{em}T_3)`$ is the weak hypercharge, and $`\mathrm{\Lambda }<M_m`$ is a dimensional parameter related to the scale of supersymmetry breaking. The functions $`f_1`$ and $`f_2`$ weakly depend on their argument and are close to 1 in the most part of the domain of definition. We consider soft mixing term in the Higgs sector $`B_\mu `$ (or $`\mathrm{tan}\beta `$) as a free parameter. The value of $`\mu `$ is determined by the relation
$$\mu ^2=\frac{1}{2}M_Z^2+\frac{m_{h_D}^2m_{h_U}^2\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1},$$
(10)
where $`m_{h_U}`$ and $`m_{h_D}`$ are soft masses of up- and down-Higgs bosons. Note that gaugino masses grow as $`n`$, whereas scalar ones grow as $`\sqrt{n}`$. It is this behavior that enables one to obtain the degeneracy $`\mu M_2`$ favorable for successful generation of baryon asymmetry.
Electroweak symmetry breaking leads to additional contributions to the soft mass spectrum. D-terms drive the masses up while mixing between left and right scalars increases the splitting between their masses. Taking into account these corrections and the renormalization group evolution from $`M_m`$ to $`M_{SUSY}`$ one calculates the low energy spectrum. In fact, reasonable estimates for the low energy soft masses are obtained by plugging $`\alpha _i=\alpha _i(M_{SUSY})`$ into Eqs. (8) and (9). As concerns the trilinear soft terms, their values at $`M_m`$ are additionally suppressed by coupling constants in comparison with Eqs. (8), (9) and become significant at low energies only due to the renormalization group evolution. The largest low energy trilinear coupling is $`A_t`$, which is numerically $`A_tM_2`$.
We consider messengers which are odd under R-parity. Then the components of the fundamental messengers have the same quantum numbers as left leptons and right down-quarks, while components of antisymmetric messengers have quantum numbers of right leptons, left quarks and right up-quarks. Therefore, symmetries allow for mixing between messengers and matter fields . In the case of one fundamental messenger multiplet, the mixing terms in the superpotential are
$$𝒲_{mm}^{(5)}=H_DL_mY_i^{(5)}E_i+H_DD_mX_i^{(5)}Q_i,$$
(11)
where $`i=1,2,3`$ counts matter generations and subscript $`m`$ refers to messenger fields. We use the standard notations for MSSM superfields ($`E_i`$ are right leptons and $`Q_i`$ are left up-quarks), $`Y_i^{(5)}`$ and $`X_i^{(5)}`$ are mixing parameters, $`L_m`$ and $`D_m`$ are messenger superfields with quantum numbers of left leptons and right down-quarks, respectively.
When messengers are integrated out, mass matrices of scalar fields acquire negative contributions,
$$\delta m_{ij}^210^2\mathrm{\Lambda }^2Y_i^{}Y_j,$$
that emerge from one-loop diagrams with messenger fields running in loops (here $`Y_i`$ generically stands for either $`Y_i^{(5)}`$ or $`X_i^{(5)}`$). In principle, these terms give rise to flavor violating processes ($`\mu e\gamma `$, $`bs\gamma `$, etc.); the corresponding limits on $`Y_i`$ are derived in Ref. . It is important for our discussion that the mixing terms also reduce some scalar masses. Indeed, let us consider colorless sector and neglect left-right mixing. Then the eigenvalues of right slepton mass matrix are
$$\{m_R^2,m_R^2,m_R^2\underset{i=1}{\overset{3}{}}\delta m_{ii}^2\},$$
where $`m_R`$ is the soft mass of right slepton. The squark masses are modified in a similar way. This hints towards the possibility to have the right stop soft mass term $`m_{\stackrel{~}{t}_R}^2`$ negative, as required by the “light stop” scenario of electroweak baryogenesis. However, fundamental messengers are not suitable for this purpose, as they do not mix with right stop. Hence, messenger fields belonging to antisymmetric representation are to be involved.
## 4 An example of GMSB model with “light stop” window
As a concrete example of a model where all criteria of “light stop” solution are satisfied, we consider GMSB model with $`n_{10}`$ antisymmetric messengers and non-zero coupling $`Y_3`$ just between right stop $`U_3`$ and corresponding messengers $`Q_m`$,
$$𝒲_{mm}^{(10)}=Y_3H_UU_3Q_m.$$
(12)
When messengers are integrated out, no additional mixing appears in squark matrix, so there is no problem with FCNC <sup>2</sup><sup>2</sup>2In what follows we will require rather large value of $`Y_3`$. A model with three large mixing terms of the same order of magnitude in right up-squark sector is ruled out by the absence of FCNC. On the other hand, if we assume the hierarchy between messenger-matter couplings similar to the hierarchy between SM fermion Yukawa couplings, we obtain a “natural” model with large messenger–stop mixing and suppressed FCNC.. As a result of mixing (12), the right stop mass squared, $`m_{\stackrel{~}{t}_R}^2`$, and up-Higgs mass squared, $`m_{h_U}^2`$, get negative contributions ,
$$\delta m_{\stackrel{~}{t}_R}^2n\frac{\mathrm{\Lambda }^2}{8\pi ^2}|Y_3|^2f_3\left(\frac{\mathrm{\Lambda }}{M_m}\right),\delta m_{h_U}^2n\frac{3\mathrm{\Lambda }^2}{16\pi ^2}|Y_3|^2f_3\left(\frac{\mathrm{\Lambda }}{M_m}\right)=\frac{3}{2}\delta m_{\stackrel{~}{t}_R}^2,$$
(13)
where $`f_3(\mathrm{\Lambda }/M_m)=(1/6)(\mathrm{\Lambda }/M_m)^2`$ at $`\mathrm{\Lambda }/M_m`$ not very close to 1. These terms come from one-loop diagrams involving the Yukawa interaction (12), with messengers running in loops.
For given $`\mathrm{\Lambda }`$ and $`n=3n_{10}`$ it is possible to tune $`Y_3`$ and obtain a negative value of $`m_{\stackrel{~}{t}_R}^2`$ of order of $`m_{crit}^2`$. In particular, Eq. (5) is obeyed at
$$|Y_3|^28\pi ^2\frac{m_{\stackrel{~}{t}_R}^2+m_{crit}^2}{3n_{10}\mathrm{\Lambda }^2}f_3^1\left(\frac{\mathrm{\Lambda }}{M_m}\right).$$
(14)
where $`m_{\stackrel{~}{t}_R}^2`$ as given by Eq. (9) is typically much larger than $`|m_{crit}|^2`$. The value of the messenger-stop Yukawa coupling turns out to be fairly large. As an example, at $`\mathrm{\Lambda }/M_m=0.5`$ one has $`Y_34\sqrt{2}\alpha _30.6`$. In this range of $`Y_3`$, Eq. (6) is also straightforward to fulfill. Hence, the fourth requirement of section 2 may be satisfied by tuning $`Y_3`$.
We turn to the other three requirements. The lower bound on the lightest Higgs boson mass is experimental and has to be satisfied regardless of baryogenesis. This is not the only constraint on the parameter space coming from collider experiments: as there is no evidence for light right sleptons, small values $`\mathrm{\Lambda }\sqrt{n}<30`$ TeV are forbidden. At the level of our accuracy we approximate the one-loop enhanced mass of the lightest Higgs boson as follows,
$$m_h^2=M_z^2\mathrm{cos}^22\beta +\frac{3\sqrt{2}}{2\pi ^2}G_Fm_t^4\left(\mathrm{log}\left(\frac{m_Q^2}{m_t^2}\right)+\frac{\stackrel{~}{A}_t^2}{m_Q^2}\right).$$
(15)
Since one of the requirements of the “light stop” solution is small left-right mixing in stop sector (see Eq. (4)), equation (15) implies the lower bound on $`\mathrm{tan}\beta `$ depending on the key GMSB parameter, $`\mathrm{\Lambda }\sqrt{n}`$, that determmines the value of $`m_Q`$ through Eq. (9). One obtains $`\mathrm{tan}\beta 3.5`$ at $`\mathrm{\Lambda }\sqrt{n}=30`$ TeV, $`\mathrm{tan}\beta 2.5`$ at $`\mathrm{\Lambda }\sqrt{n}=50`$ TeV and $`\mathrm{tan}\beta 1.5`$ at $`\mathrm{\Lambda }\sqrt{n}=100`$ TeV. Within our approximation, the largest value of $`m_h`$ consistent with “light stop” solution, $`m_h<115`$ GeV, is achieved at very large $`\mathrm{\Lambda }\sqrt{n}`$ and hence $`m_Q`$: $`\mathrm{\Lambda }\sqrt{n}400`$ TeV, $`m_Q4`$ TeV. The reason is that, as we will see below, successful baryogenesis requires rather small values of $`\mathrm{tan}\beta `$. It worth noting that similar constraints are imposed on the parameters of SUGRA models with “light stop” solution.
Let us proceed with other requirements. Before discussing concrete cases of small and large CP-violating phase $`\varphi _{CP}`$, let us make a few general remarks.
First, mixing (12) drives $`m_{h_U}^2`$ deep into the region of negative values,
$$m_{h_U}^2=m_{h_D}^2\frac{3}{8\pi ^2}m_Q^2\left(\mathrm{log}\frac{M_m}{m_Q}+\frac{3}{2}\right)+\frac{3}{2}\delta m_{\stackrel{~}{t}_R}^2,$$
(16)
that leads to very large $`\mu `$ because of Eq. (10). This implies that the CP-phase has to be relatively large, since $`m_A\mu `$, and the baryon asymmetry is suppressed by the mass of heavy CP-odd Higgs boson. Moreover, since an approximate degeneracy $`\mu M_2`$ is required, wino is also heavy. This makes a potential problem, because in GMSB models with complete messenger multiplets one has $`A_tM_2`$, and there appears large left-right mixing. The cure is to make $`\stackrel{~}{A}_t=|A_t\mu \mathrm{cot}\beta |`$ small, $`\stackrel{~}{A}_t<0.5m_Q`$, by choosing suitable sign of $`\mu `$ and tuning $`\mathrm{tan}\beta `$.
Second, the degeneracy between $`\mu `$ and $`M_2`$ will be achieved by a suitable choice of the number of messengers, $`n_{10}`$. Indeed, these two mass parameters scale differently with the number of messenger fields. Namely, from Eq. (10) one finds that $`\mu `$ is determined by scalar soft masses, so it grows as $`\sqrt{n}`$ (see Eq. (9)), while $`M_2`$ grows as $`n`$ (see Eq. (8)).
One can estimate viable values of parameters of this model by the following simple chain of reasonings. Squarks are the heaviest scalars in this theory, because the hierarchy of soft terms are governed by corresponding gauge couplings. Since $`\mathrm{\Lambda }\sqrt{n}30`$ TeV, the model has heavy strong sector, $`m_Q^2m_Z^2`$. The “light stop” solution requires $`|m_Q^2+\delta m_{\stackrel{~}{t}_R}^2|M_z^2`$, so $`|\delta m_{\stackrel{~}{t}_R}^2|`$ is also large, and $`m_{h_U}^2\frac{3}{2}\delta m_{\stackrel{~}{t}_R}^2`$ (see Eqs. (16), (13)). Then the relation (10) may be approximated as follows,
$$\mu ^2=\frac{3}{2}m_Q^2\frac{\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1}.$$
(17)
The “light stop” solution requires almost degenerate spectrum, $`\mu =ϵM_2`$ with $`1/2<ϵ<2`$. In view of the relation $`A_tM_2`$, Eq. (4) implies
$$\left|\frac{1}{ϵ}\frac{1}{\mathrm{tan}\beta }\right|\frac{1}{2}\frac{m_Q}{\mu },$$
(18)
Making use of Eq. (17) we obtain from the inequality (18)
$$\left|\frac{1}{ϵ}\frac{1}{\mathrm{tan}\beta }\right|\frac{1}{\sqrt{6}}\sqrt{1\frac{1}{\mathrm{tan}^2\beta }}.$$
(19)
This determines the viable values of $`\mathrm{tan}\beta `$ for given degeneracy $`ϵ`$. At $`ϵ=1`$ ($`M_2=\mu `$) one has $`\mathrm{tan}\beta 1.5`$, at $`ϵ=1.5`$ one has $`\mathrm{tan}\beta 4`$ while at $`ϵ=2`$ one obtains $`\mathrm{tan}\beta 10`$.
By making use of Eq. (17) and Eqs. (8), (9) we find the number of messengers required for succesfull electroweak baryogenesis at given $`ϵ`$ and $`\mathrm{tan}\beta `$,
$$n_{10}\frac{4}{3ϵ^2}\frac{\alpha _3^2}{\alpha _2^2}\frac{\mathrm{tan}^2\beta }{\mathrm{tan}^2\beta 1}.$$
(20)
The estimates are: $`n_{10}24`$ at $`ϵ=1`$ and $`\mathrm{tan}\beta 1.5`$, $`n_{10}6`$ at $`ϵ=1.5`$ and $`\mathrm{tan}\beta 4`$, $`n_{10}3`$ at $`ϵ=2`$ and large $`\mathrm{tan}\beta 10`$.
Until now we discussed mostly the bounds coming from the survival of baryon asymmetry after the electroweak phase transition. Additional bounds come from the calculations of the baryon asymmetry along the lines of Ref. . Let us present the results for the cases of small and large $`\varphi _{CP}`$. They depend on the value of $`\mathrm{\Lambda }`$ which determines the mass of the CP-odd Higgs boson $`m_A`$, and hence the suppression factor in the baryon asymmetry.
The smallest CP-phase corresponds to the case of complete degeneracy between $`M_2`$ and $`\mu `$, i.e., $`\mu /M_2ϵ=1`$. In this case $`\mathrm{tan}\beta 1.5`$, so one has to take $`\mathrm{\Lambda }\sqrt{n}100`$ TeV in order to obtain $`m_h>100`$ GeV. As discussed above, the number of messengers is large, $`n_{10}24`$. This results in the large value of Higgsino mass $`\mu 2200`$ GeV and large $`m_A\mu /\mathrm{sin}\beta `$, that makes the phase transition stronger, but reduces the amount of baryon asymmetry at given $`\varphi _{CP}`$. The realistic value of $`n_B/s`$ is obtained at $`\mathrm{\Lambda }\sqrt{n}=100`$ TeV with $`|\mathrm{sin}\varphi _{CP}|0.008`$. At $`\mathrm{\Lambda }\sqrt{n}=200`$ TeV we find $`\mathrm{tan}\beta 1.5`$, $`n_{10}24`$, $`\mu 4400`$ GeV and $`|\mathrm{sin}\varphi _{CP}|0.03`$.
In models with large CP-violation phase ($`|\mathrm{sin}\varphi _{CP}|=1`$), fine tuning between $`M_2`$ and $`\mu `$ may be somewhat relaxed, and the constraints on $`n_{10}`$ and $`\mathrm{tan}\beta `$ are, of course, weaker. As an example, let us consider the case $`\mu =1.5M_2`$. At $`\mathrm{\Lambda }\sqrt{n}=50`$ TeV one obtains $`n_{10}6`$, $`\mu `$ 800 GeV and $`\mathrm{tan}\beta 4`$, while at $`\mathrm{\Lambda }\sqrt{n}=100`$ TeV one gets $`n_{10}6`$, $`\mu `$ 1600 GeV and $`\mathrm{tan}\beta 4`$.
Let us note in passing that the case $`\mu <M_2`$ is not favorable for the “light stop” solution. Indeed, it follows from Eq. (19) that the viable region of $`ϵ\mu /M_2<1`$ is very narrow and $`\mathrm{tan}\beta `$ is smaller than $`1.5`$ at $`ϵ`$ within this interval. Also, small values of $`\mathrm{\Lambda }\sqrt{n}`$ are disfavored. As an example, at $`\mathrm{\Lambda }\sqrt{n}=30`$ TeV one finds $`\mu `$ 450 GeV, and the absence of right slepton <sup>3</sup><sup>3</sup>3In the case of large $`\mu `$ which is of relevance throughout our discussion the masses squared of right sleptons acquire negative contributions due to electroweak symmetry breaking, which are proportional to $`\mu \mathrm{tan}\beta `$ . with mass smaller than 45 GeV requires $`\mathrm{tan}\beta 3`$. This is in contradiction with the bound coming from the lightest Higgs boson mass, $`\mathrm{tan}\beta 3.5`$.
One concludes that all requirements for the “light stop” solution are satisfied in this model, hence electroweak baryogenesis is capable of generating sufficient amount of baryon asymmetry in GMSB models. The region of viable parameters is fairly narrow due to strong restrictions coming from the search for the lightest Higgs boson. Indeed, $`\mathrm{tan}\beta `$ in models with “light stop” solution tends to be small, while existing limits on the lightest Higgs boson mass favor high $`\mathrm{tan}\beta `$. At small CP-phase the suitable region stretches along $`\mathrm{tan}\beta 1.5`$ and the smallest CP-phase, $`\varphi _{CP}10^2`$, is possible at $`\mathrm{\Lambda }\sqrt{n}`$100 TeV. At $`\varphi _{CP}1`$ this region becomes wider and extends to $`\mathrm{tan}\beta 4`$. The suitable region of parameter space is the largest at the smallest possible $`m_h`$. A drawback of the model is a large number of messenger fields, so the gauge coupling constants become large below the scale of possible Grand Unification.
Electroweak baryogenesis would work better in a model where left squarks are heavy, whereas right squarks are light (or $`\mu `$ is relatively small) and $`A_t`$ is small. In the framework of GMSB these features appear in models with additional large mixing terms or extended messenger content. Let us briefly discuss these possibilities.
One can introduce additional mixing (e.g., between d-squark and corresponding messengers), which provides large negative contribution to mass squared of down-Higgs. Then larger values of $`\mathrm{tan}\beta `$ may become viable. Likewise, the value of $`\mu `$ may be smaller, especially at low $`\mathrm{tan}\beta `$, which would extend the region of suitable values of CP-phase.
One can also consider messengers which do not form complete GUT multiplets. <sup>4</sup><sup>4</sup>4The same observations apply to models where messengers carrying different quantum numbers are characterized by different scales $`\mathrm{\Lambda }`$. Let us denote the effective numbers of messengers generating gluino mass and wino mass as $`n_c`$ and $`n_w`$, respectively. Then the relevant quantity will be $`n_w^2/n_c`$ instead of $`n_{10}`$. Therefore, it will be possible to reduce significantly the number of messenger fields and relax the problem of unification of gauge coupling constants in a theory with “light stop” window.
There is yet another possibility related to messengers belonging to incomplete GUT multiplets. We outline it by making use of the simplest example of only lepton-like messengers without any messenger-matter mixing. In this model, right stop acquires negative mass squared due to renormalization group evolution, in analogy to the up-Higgs in MSSM. There is a wide region of parameter space, where electroweak symmetry breaking occurs but left squarks remain heavy. In this model the favorable hierarchy
$$m_Qm_{\stackrel{~}{t}_R},A_t$$
appears naturally. The reason for the hierarchy between $`m_Q`$ and $`A_t`$ is that the largest contribution to the low energy value of $`A_t`$ comes from two-loop diagrams involving gluino; in this model gluino is light, that reduces $`A_t`$ significantly. As a consequence, a fairly wide region of $`\mathrm{tan}\beta `$ will be consistent with the “light stop” solution in this model.
## 5 Conclusions
We have demonstrated that electroweak baryogenesis may produce acceptable amount of baryon asymmetry in the framework of GMSB models with messenger-matter mixing. The required MSSM spectrum is similar to the so-called “light stop window” of SUGRA models . We have plainly mapped that “window” onto the space of parameters of GMSB models. We have found that severe constraints are to be imposed on the GMSB parameters. At least one of the mixing terms has to be quite large in order to make right stop lighter than top. In the explicit example presented in this letter, the minimal possible value of CP-phase is reached at a large number of antisymmetric messenger generations, $`n_{10}24`$. With maximal CP-violation ($`\varphi _{CP}1`$) one has $`n_{10}6`$. Another property of GMSB models with electroweak baryogenesis is $`\mathrm{tan}\beta 4`$. We have briefly outlined also extensions of this model which have better properties with respect to electroweak baryogenesis. Let us note finally, that additional sources of CP-violation, which may originate from messenger-matter mixing, may enhance the electroweak production of baryon asymmetry.
## 6 Acknowledgments
The author is indebted to V. Rubakov and S. Dubovsky for useful discussions. This work was supported in part under Russian Foundation for Basic Research grant 99-02-18410 and by the Russian Academy of Science, JRP grant # 37. |
warning/0001/astro-ph0001219.html | ar5iv | text | # Disc galaxy evolution models in a hierarchical formation scenario: structure and dynamics
## 1 Introduction
The understanding of galaxy formation and evolution provides a crucial link between a large body of astronomical observations and cosmological theories. According to the inflationary cold dark matter (CDM) theory, large-scale cosmic structures arise from primordial density fluctuations that evolve gravitationally through a hierarchical process of mass aggregation (accretion and merging). In order to explain galaxy formation, it is also necessary to take into account the hydrodynamical and dissipative processes related with the baryon matter, as well as the poorly understood processes of star formation (SF) and its feedback. In this paper we use galaxy evolution models, which are able to predict the local and global properties of disc galaxies, and integrate them into the cosmological framework.
In the last decade, several galaxy formation and evolution models based on semi-analytical and analytical approaches were developed (e.g. White & Frenk 1991; Lacey & Silk 1993; Kauffmann, White & Guiderdoni 1993; Cole et al. 1994; Baugh, Cole & Frenk 1996; Somerville & Primack 1998; Fall & Efstathiou 1980; Dalcanton, Spergel & Summer 1997; Mo, Mao & White 1998, hereafter MMW98; van den Bosch 1998). The models we present in this paper make use of some of the ingredients of these approaches. However, as opposed to most of the previous models, we calculate the internal structure, dynamics, hydrodynamics and SF process that define the evolution of a given galaxy (dark+luminous). That is why our approach is based on models of disc galaxy evolution (Firmani, Hernández & Gallagher 1996). We link the initial and boundary conditions of these models to the cosmological setting through the hierarchical mass aggregation histories (MAHs) and the structure of the dark matter (DM) haloes. We relax the simplifying hypothesis of a universal DM halo density profile and model the evolution of disc galaxies in a wide collection of DM structures. Thus, we can explore whether the initial conditions given by the CDM models are consistent with the disc galaxy properties and their correlations in the local universe. We have carried out a sequence of semi-numerical calculations, which include (i) the generation of the hierarchical MAHs from the primordial density fluctuation field, (ii) the gravitational collapse and virialization of the DM haloes, (iii) the formation of discs in centrifugal equilibrium within the evolving haloes, (iv) the gravitational drag of the collapsing gas upon the DM halo, and (v) the evolution of a galaxy disc, including SF and hydrodynamics. Some details were already disscused in previous papers (Firmani et al. 1996; Avila-Reese, Firmani & Hernández 1998, hereafter AFH98; see also Avila-Reese 1998).
Our scenario suggests that discs form inside-out within a growing DM halo with a rate of gas accretion proportional to the rate of cosmological mass aggregation (Gunn 1982; 1987; Ryden & Gunn 1987; Avila-Reese & Firmani 1997; AFH98). We assume that the primordial angular momentum is acquired by the protogalaxy through tidal torques during its linear gravitational regime. This scenario will be called the “extended collapse scenario” as opposed to the “merging scenario” where the main properties of galaxies are a consequence of disc merging. Because the discs are dynamically frail objects (e.g. Tóth & Ostriker 1992), mergers could not have been relevant in establishing the main properties of disc galaxies, which constitute nearly the 80 per cent of normal present-day galaxies. Compelling new observational evidence suggests that a large fraction of baryon matter is in the form of small cold gas clumps which are falling onto galaxies and whose nature probably is cosmological (Blitz et al. 1998; López-Corredoira, Beckman, & Casuso 1999).
The structural and dynamical properties of disc galaxies as well as the correlations amongst them should be self-consistently explained by theories of galaxy formation and evolution. Among the local properties to be explained are the nature of the exponential surface brightness or density distribution in discs, the factors that determine the observed wide range of surface brightnesses (SB), the shapes of the rotation curves (typically flat) and their correlations with the SB. Among the global structural and dynamical correlations of disc galaxies, the Tully-Fisher relation (TFR) and the rotation velocity-radius relation are of great interest. Understanding of the origin of these observational relations —particularly the TFR and its scatter— is an outstanding problem of extragalactic astronomy. We would like to know whether the slope and zero-point of the TFR are produced by evolutionary and/or SF processes (e.g. Silk 1997; Elizondo et al. 1998), or if they are direct imprints from a fundamental scaling law of cosmic structures (e.g. Faber 1982; Frenk et al. 1988; Cole et al. 1994; MMW98; AFH98; Navarro & Steinmetz 1998). Likewise we are interested in understanding why the TFRs for high and low SB galaxies are the same, and why the scatter of this relation is small.
The purpose of this paper is to use an unified scenario of galaxy formation and evolution in the cosmological context in order to explain the items mentioned above. In a forthcoming paper, results concerning the luminosity properties, bulge formation, and the correlations which define the Hubble sequence will be discussed (Avila-Reese & Firmani 2000). Since our aim is to study general trends, we shall only use one representative CDM model. It has the following parameters: the matter density: $`\mathrm{\Omega }_m=0.35,`$ the vacuum density: $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65,`$ the Hubble constant in unities of 100 kms<sup>-1</sup>Mpc<sup>-1</sup>: h$`=0.65`$, the amplitude of perturbations on a 8 Mpch<sup>-1</sup> scale: $`\sigma _8=1`$. Hereafter, we shall refer to this model as the $`\mathrm{\Lambda }`$CDM<sub>0.35</sub> model.
In $`\mathrm{\S }2`$ we briefly describe the methods used to generate the MAHs and to calculate the virialization of the DM haloes. Results concerning the structural and dynamical properties of the DM haloes are presented. The method used to build discs within DM haloes and to follow their evolution, including SF and hydrodynamics, is presented in $`\mathrm{\S }3`$. In $`\mathrm{\S }`$4, local galaxy properties such as the surface density, the rotation curve, and the decomposition of the rotation curve are presented. Section 5 is devoted to the TFR, its scatter, and the correlations between gas fraction vs. SB and $`V_{\mathrm{max}}`$ vs. disc size. We interpret why low and high SB galaxies lie on the same location in the TFR, and why the residuals of the TFR are nearly independent of the SB. In §6 we explore how the introduction of a shallow core in the halo influences on the models. A summary and our conclusions are given in $`\mathrm{\S }7`$.
## 2 Dark matter haloes
### 2.1 The method
According to the hierarchical clustering scenario, the DM haloes are the environment in which luminous galaxies form and evolve. Therefore, some features of luminous galaxies depend on the evolution and properties of their surrounding dark haloes. In this subsection we present the method we used to calculate the formation, evolution and structure of the DM haloes.
We calculate the gravitational collapse and virialization of DM haloes beginning from a primordial density fluctuation field. This was done in AFH98 and the reader is referred to it for details (see also Avila-Reese 1998). We assume a Gaussian fluctuation field where all the statistical properties depend only on the power spectrum of fluctuations given by the chosen cosmological model. The conditional probability for Gaussian random fields is used to calculate the mass distribution of haloes at time $`t_{i+1}`$ that will be contained in a halo of mass $`M_i`$ at a later time $`t_i`$ (Bower 1991; Bond et al. 1991). Starting from a present-day mass and its cumulative density contrast, we apply recurrently this mass distribution through Monte Carlo simulations in order to construct different realizations of the MAHs of “main progenitors” (Lacey & Cole 1993). Main progenitor in our case is the most massive subunit of the distribution, at each time, and not the larger mass between the chosen mass $`M_{i+1}`$ and its complement $`M_i^{}`$, where $`M_i^{}=M_iM_{i+1}`$.
To calculate the virialized profile of the dark halo that emerges from a MAH, we use a method that expands upon the secondary infall model (e.g. Gunn 1977; Zaroubi & Hoffman 1993) by allowing non-radial orbits and arbitrary initial conditions (MAH in our case). Under the assumption of spherical symmetry, an initial density fluctuation —given by the MAH— can be described as a set of concentric mass shells that sequentially attain their first maximum expansion radius $`r_0`$. After this maximum, the shell at $`r_0`$ which contains mass $`m_0(r_0)`$, separates from the expansion of the universe, falls toward the center due the gravitational field of the internal mass, and evolves through a radially oscillatory movement. Each oscilation defines an apapsis radius $`r_a`$ which is function of $`r_0`$. The gravitational field at radius $`r`$ is given by $`m_T(r)`$, the sum of the mass shells with $`r_ar`$ that permanently oscillate inside $`r`$, $`m_P(r)`$, and of the mass shells with $`r_ar`$ that only momentarily fall inside $`r`$, $`m_M(r)`$. The estimate of $`m_M(r)`$ is obtained from the probability
$$P(r)_0^r\frac{d\eta }{v\left(\eta \right)}$$
to find a shell inside a radius $`r`$($`<r_a`$), where $`v\left(r\right)`$ is the radial velocity of the shell. This velocity is:
$$v^2\left(r\right)=2\left[EG_0^r\frac{m_T\left(\eta \right)}{\eta ^2}𝑑\eta \frac{j^2}{2r^2}\right],$$
where $`G`$ is the gravitational constant, $`E`$ is the total energy of the shell, and $`j`$ is the typical angular momentum per unit mass of a shell mass element due to the thermal motion which will be taken constant in time. The condition $`v=0`$ defines the apapsis $`r_a`$ and a value $`r_p`$ which determines the maximum penetration of the shell towards the center. We express $`j`$ through the parameter $`e_0\left(\frac{r_p}{r_a}\right)_0`$ and, although it is defined for the first $`r_a`$ and $`r_p`$ of every shell, through $`e_0`$ we are parametrically taking into account the thermal energy that could be produced by the mergers of substructures and tidal forces at all times. This parameter is calibrated to the results of N-body cosmological simulations ($`e_00.10.3`$). The sequential aggregation of new shells, combined with their motion toward the center, introduces new contributions to the gravitational field which acts on the underlying shells. This non-conservative spherical gravitational field changes $`E`$, and consequently $`r_a`$, $`r_p`$, and $`e`$ of each shell. The contraction of $`r_a`$ leads to an asymptotic value which is identified as the current virialized radius $`r_v.`$ The change of $`E`$ may be estimated by assuming an adiabatic invariant for the orbital motion. A simple iterative numerical method allows us to calculate the solution, i.e. the structure profile $`m_T(r,t)`$ at each time step.
The present-day mass used to initiate the Monte Carlo simulations will be called the nominal mass $`M_{\mathrm{nom}}`$. Our results show that the outer mass shells within $`M_{\mathrm{nom}}`$ are not virialized. The mass shells that have been virialized are roughly within the virial radius, $`r_v`$, where the mean over-density drops below the critical value, $`\mathrm{\Delta }_c`$, given by the spherical collapse model; for the cosmological model used here $`\mathrm{\Delta }_c(z=0)=334`$ (e.g. Bryan & Norman 1998). Analysis of numerical simulations show that at radii smaller than $`r_v`$, mater is indeed close to virial equilibrium (Cole & Lacey 1996; Eke, Navarro & Frenk 1998). At radii between $`r_v`$ and 2$`r_v`$, matter is still falling onto the halo, while, at larger radii, matter is expanding with the universe. The mass contained within $`r_v`$ is the virial mass $`M_v`$ which, depending upon the MAH, is 0.7-0.9 times $`M_{\mathrm{nom}}`$ (see also Kull 1999).
### 2.2 Structure of the haloes: diversity
In AFH98 it was shown that, due to the statistical nature of the MAHs, a collection of different DM virialized configurations are produced, the most typical configurations being reasonably well described by the density profile proposed by Navarro, Frenk & White (1996, 1997; hereafter NFW). Some features of the galaxies are related to the dispersion about the virialized structures. Therefore, its inclusion in the galaxy models is important. In the present work, we apply a more accurate statistical treatment of the results than in AFH98. Instead of initially selecting some relevant special cases, we generate, for a given mass and through Monte Carlo simulations, a catalog of objects with different MAHs (and eventually different spin parameters $`\lambda `$ taken from a log-normal distribution; see $`\mathrm{\S }3`$). Although most of the results of the statistical analysis of this catalog agree with AFH98, we felt it was important to perform this new procedure in order to obtain an accurate estimate of the dispersion in the mass-velocity relation.
In Figs. 1a and 2a, a sample of twenty MAHs ($`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$) and the circular velocity profiles of the haloes formed from them, are plotted as functions of the collapse redshift, $`z_c`$. The average MAH and two representative cases, a fast early collapse (L: low accretion rate at $`z0`$) and an extended collapse (H: high accretion rate at $`z0`$), are shown in Fig. 1b. The circular velocity profiles corresponding to these particular MAHs are plotted in Fig. 2b. These two MAHs were chosen in such a way that roughly 95 per cent of all the trajectories lie between of them. For a given mass, the average MAH is calculated as the mean of a large sample of trajectories.
We find that a DM halo formed from an early fast MAH is more concentrated than a halo produced by a gentle and extended MAH. Nevertheless, as seen in Fig. 1a, the MAHs are diverse and difficult to classify in uniparametrical, or even biparametrical, sequences. That is why we have preferred to generate a a catalog of MAHs for each given $`M_{\mathrm{nom}}`$ and to study the statistics of the DM haloes calculated from these MAHs a posteriori.
Some of the MAH trajectories show pronounced jumps, which indicates the occurrence of major mergers. Halo major merging may imply disc major merging that probably destroys the discs. The fraction of MAHs that show evidence of at least one major merger is between 20 and 30 per cent. When we use the results from our catalogs to calculate the TFR scatter, we keep all the MAHs since we want to obtain an upper limit of this scatter.
### 2.3 Comparison with N-body simulations and the $`M_vV_{\mathrm{max}}`$ relation
The density profiles of a large fraction of the calculated haloes are roughly described by the NFW profile. However, as it is seen in Fig. 2b, the haloes present a diversity of structures. In Avila-Reese et al. (1999) the density profiles of the haloes obtained with our method were compared with those of thousands of haloes from a cosmological N-body simulation. The statistical agreement found was rather good, in particular for the isolated haloes. If we compare, for example, the statistical distribution of the outer density profile slope ($`\beta `$) for our haloes with those from the N-body simulations, we find they are very similar (see Fig. 4 in Avila- Reese et al. 1999). Haloes in the N-body simulations present a slightly broader distribution of $`\beta `$ than our haloes.
It is also important to mention that the MAHs of our haloes for a given mass agree rather well with the mass evolution measured for the haloes in a cosmological N-body simulation (see Fig. 6 in Gottlober, Klypin & Kravtsov 1999).
As in previous works, we find that the mass $`M_v`$ and the halo maximum circular velocity $`V_{\mathrm{max}}`$ obey a $`M_vV_{\mathrm{max}}^\alpha `$ relation (in the $`M_v4\times 10^{10}M_{}4\times 10^{12}M_{}`$ range and for the $`\mathrm{\Lambda }`$CDM<sub>0.35</sub> model, $`\alpha 3.2`$ and 3.3 from our results and from results of N-body simulations, respectively). This relation is imprinted by the power spectrum of fluctuations and the MAH of the protohalo (AFH98; Avila-Reese et al. 1999).
Owing to the statistical nature of the calculated MAHs, a scatter in the $`M_vV_{\mathrm{max}}`$ relation is expected. A first attempt to estimate such a scatter was done by Eiseinstein & Loeb (1996) who used Monte Carlo simulations to generate the MAHs, and the simple spherical-collapse model to calculate the halo circular velocities. AFH98 have estimated this scatter making use of their semi-numerical method for calculating the virialization process of the DM haloes. For the SCDM model both methods lead to similar results. Here, for the $`\mathrm{\Lambda }`$CDM<sub>0.35</sub> model and in the range of $`M_v4\times 10^{10}M_{}4\times 10^{12}M_{}`$, we obtain fractional standard deviations in the velocity, $`\sigma _V/<V_{\mathrm{max}}>`$, from $`0.10`$ to $`0.07`$, respectively (see Table 1). The deviations estimated from the cosmological N-body simulations are in rough agreement with these values, being only slightly larger than our results. (see Fig. 10 in Avila-Reese et al. 1999). The deviations in velocity may be translated into logarithmic standard deviations in mass: $`\mathrm{\Delta }`$log$`M_v=\alpha `$log$`(1+\sigma _V/<V_{\mathrm{max}}>)`$ (multiplying $`\mathrm{\Delta }`$ log$`M_v`$ by 2.5, the standard deviation in mass can be expressed magnitudes). In column (4) of Table 1, the mass scatter of the $`M_vV_{\mathrm{max}}`$ relation — expressed in magnitudes — is given for several masses.
## 3 Disc build up and galaxy evolution
During or after DM reaches virial equilibrium, the baryon matter dissipates energy radiatively and falls to the bottom of the gravitational potential well. If the DM halo has some angular momentum, then a disc in centrifugal equilibrium forms at the centre of the halo. In this section we describe the methods used to calculate the formation and evolution of discs. The main structural properties of the disc will depend on the formation history of the halo, its structure, and on the amount of infalling gas and its angular momentum.
Analytical models of discs in centrifugal equilibrium surrounded by DM haloes have been used to study several galaxy and galaxy population features (e.g. Fall & Efstathiou 1980; van der Kruit 1987; Dalcanton et al. 1997; MMW98; van den Bosch 1988). The results obtained with these models encourage us to study in more detail the extended collapse scenario. Our contributions are: First, we were able to build up the disc sequentially within an evolving DM halo. Second, we calculated the disc thickness and the multiple galaxy components in a 3-D gravitational potential. Third, we calculated the SF, disc hydrodynamics and ensuing galaxy evolution.
The outline of our evolution models is:
(1) We consider that baryon matter has the same distribution of mass and angular momentum as DM till the accreted spherical (baryon+DM) mass shell virializes. We assume that each mass shell has a solid body rotation, in agreement with the Zel’dovich approximation, and that the rotation axis of the shells are aligned. We did not assume that the whole halo is a solid rotator (e.g. Dalcanton et al. 1997).
(2) Once a given mass shell virializes, a fraction $`f_d`$ of its mass is transferred to a disc in centrifugal equilibrium. Since for galaxy haloes the time-scale of gas cooling is generally smaller than the dynamical time-scale (c.f., Silk 1977; Rees & Ostriker 1977; White & Rees 1978; Ryden & Gunn 1987) we assume the gas falls from the maximum expansion radius of the current shell to the halo centre in a time equal to the shell virialization time. The radial distribution of the infalling gas is calculated by equating its specific angular momentum (the same of the DM component) to that of its final circular orbit in centrifugal equilibrium. The specific angular momentum $`j_{sh}`$ acquired by each collapsing mass shell during the linear regime is estimated under the assumption of a constant spin parameter $`\lambda `$ ($`\frac{J\left|E\right|^{1/2}}{GM^{5/2}}`$):
$$j_{sh}(t)=\frac{dJ(t)}{dM(t)}=\frac{GM(t)^{5/2}\lambda }{\left|E(t)\right|^{1/2}}\left(\frac{5}{2}\frac{1}{M(t)}+\frac{d\left|E(t)\right|}{2dM(t)}\right),$$
(1)
where $`J`$, $`M`$, and $`E`$ are the current total angular momentum, mass, and energy of the halo at time $`t`$.
(3) The gravitational drag on the dark halo produced by the collapse of each baryon mass shell is calculated through the adiabatic invariant formalism (e.g. Flores et al. 1993).
(4) We consider that the local SF and internal hydrodynamics of the disc are regulated by a balance between the kinetic energy injected by SNe and gas accretion, and the energy dissipated by the turbulent interstellar medium (Firmani et al. 1996, 1997). The star formation is turned on at radius $`r`$ when the local Toomre gravitational instability parameter for the gas disc, $`Q_g(r)\frac{v_g(r)\kappa (r)}{\pi G\mathrm{\Sigma }_g(r)}`$, falls below a given threshold; $`\kappa (r)`$, $`v_g(r)`$, and $`\mathrm{\Sigma }_g(r)`$ are the epicyclic frequency, the gas velocity dispersion, and the gas surface density at radius $`r`$, respectively. Thus, the SF is controlled by a feedback mechanism such that, when a gas disc column is overheated by the SF activity, SF is inhibited and the disc column dissipates the excess energy to lower $`v_g`$ back to the value determined from the Toomre criterion threshold. Numerical simulations (Sellwood & Carlberg 1984; Carlberg 1985; Gunn 1987) and observational estimates (e.g. Skillman 1987; Kennicutt 1989) suggest thresholds of the order of 2 instead of 1, as was analytically obtained for a thin disc (Toomre 1964). This difference is attributed to collective phenomena which are difficult to account for in the analytical studies. In our models, the value of this threshold controls the thicknesses of the gas and stellar discs; the SF rate actually is rather insensitive to the value of $`Q_g`$. When a value of 2 is used, for a model of the Galaxy, we obtain gas and stellar disc thicknesses compatible with those of the solar neighborhood. Thus, we fix $`Q_g=2`$. The gas loss from stars is also included. The gravitational dynamics of the evolving star and gas discs, and the DM halo are treated in detail. A Salpeter initial mass function is assumed. Analytical fits to simple stellar population models are used in order to calculate the luminosity in the $`B`$ band (see for more details Firmani & Tutukov 1994; Firmani et al. 1996).
Note that in our self-regulating SF mechanism the feedback happens only within the disc and not at the level of the whole halo. This is justifiable since the turbulent interstellar medium is a very dissipative system (e.g., Avila-Reese & Vázquez-Semadeni 2000). Gas and energy outflows are confined within a region close to the disc. The thick gaseous disc and the global magnetic field are efficient shields that prevent any outflow toward the halo on large scales (e.g. Mac Low, McCray & Norman 1989; Slavin & Cox 1992; Franco et al. 1995). The lack of observational evidence for significant amounts of hot gas in the haloes of disc galaxies confirms these studies and suggests that the halo-disc connection is not enough to self-regulate and drive the disc SF. Some observational evidence shows that extragalactic gas clouds reach the disc in free fall (see Blitz et al. 1998).
According to the scheme described above, time by time and at each radius, the growing disc is characterized by the infall rate of fresh gas by unit of area, $`\dot{\mathrm{\Sigma }_g}(r,t)`$, the gas and stellar disk surface density profiles, $`\mathrm{\Sigma }_g(r,t)`$ and $`\mathrm{\Sigma }_s(r,t)`$, the total rotation curve (including the growing DM halo component), $`V_r(r,t)`$, and the SF rate $`\dot{\mathrm{\Sigma }}_s(r,t)`$ determined by the energy balance in the vertical gaseous disk and by a Toomre criterion. It is interesting to point out that a consequence of our physical model of SF is that $`\dot{\mathrm{\Sigma }}_s`$ exhibits a Schmidt law with index $`2`$ at all radii and during almost all the evolution (Firmani et al. 1996; Avila-Reese & Firmani 2000). It should be mentioned that the global SF rate efficiency in our models does not depend on the mass or $`V_{\mathrm{max}}`$ of the galaxy (see Avila-Reese & Firmani 2000). The global SF rate efficiency is mainly a function of $`\mathrm{\Sigma }_g`$ (determined by $`\lambda `$) and $`\dot{\mathrm{\Sigma }}_g`$ (determined by the MAH).
Our goal is to generate a catalog of models for each given present-day nominal mass, $`M_{\mathrm{nom}}`$ (see §2). The key initial factors for a model of a given mass are (i) the MAH (it determines the structure of the halo and the rate at which the gas is accreted onto the disc), (ii) the spin parameter $`\lambda `$ (strongly influences the size and surface density of the disc), and (iii) the effective baryon fraction $`f_d`$ that is incorporated into the disc. The MAHs (see Fig. 1) and the halo evolution are calculated as was described in $`\mathrm{\S }2`$. The spin parameter $`\lambda `$ is taken to be constant in time and is chosen from a log-normal distribution through Monte Carlo simulations. The median and the dispersion of the log-normal distribution we use are 0.05 and 0.5, respectively, in agreement with results of several theoretical and numerical studies (see MMW98 and the references therein). Regarding item (iii), in a first approximation one might take the baryon-to-dark matter ratio of the protogalaxies equal to that of the universe, $`f_b=\mathrm{\Omega }_b/\mathrm{\Omega }_m`$. According to the current primordial abundance determinations of the baryon density, $`\mathrm{\Omega }_bh^20.0060.013`$ (e.g. Fukugita, Hogan & Peebles 1998), the baryon-to-dark matter fraction of the universe is $`f_b0.040.09`$ for the cosmological model used in this paper ($`\mathrm{\Omega }_m=0.35`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0.65`$, $`h=0.65`$). However, it is probable that not all the halo baryon fraction ends up in the galaxy central disc; some gas fraction may remain in the form of hot halo gas or may also be transformed into stars within some small halo substructures. Here we shall assume that the disc contains 5 per cent ($`f_d=0.05`$) of the total halo mass (see also Dalcanton et al. 1997; MMW98). We are aware that this is a free parameter in our models and when necessary, we shall explore here models for other $`f_d`$ values.
## 4 Local properties of disc galaxies
In this section we present the stellar surface density and rotation velocity profiles, and the rotation curve decompositions for our model galaxies. Most of the results in this section refer to models with the H, average, and L MAHs (see $`\mathrm{\S }2`$), and with $`\lambda =0.03,0.05,`$ and 0.08.
### 4.1 Stellar surface density
The stellar surface density (SSD) profiles of the rotationally supported discs, sequentially formed within the evolving CDM haloes, are nearly exponential over several scale lengths. Fig. 3 shows the SSD profiles for models of $`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$ with several representative values of the MAH (panel a) and the spin parameter $`\lambda `$ (panel b). It is worth remarking that, in contrast to the models of Dalcanton et al. (1997; see also Fall & Efstathiou 1980), we have not assumed the whole protogalaxy as a solid body rotator (see $`\mathrm{\S }3`$). Therefore, the exponentiality of the stellar disc profile is a prediction of our model where the crucial assumption made regarding disc formation is the time constancy of $`\lambda `$. Analysis of results of N-body simulations regarding the evolution of the spin parameter of the haloes seem to confirm this assumption (Gottlober, private communication). We find the same general trends as in Dalcanton et al. (1997) and Jimenez at al. (1998): the central SSD depends strongly upon $`\lambda `$, and less upon the mass. We also find that the SSD has some dependence upon the MAH (Fig. 3a). We conclude that the low SB galaxies may be mainly galaxies with high angular momenta and/or low masses. The value $`\lambda 0.05`$ probably introduces a natural separation between high and low SB galaxies.
The characteristic radii and SB of our disk galaxy models are realistic (see also Avila-Reese & Firmani 2000). The assumptions of detailed angular momentum conservation and rotation axis aligment of the shells were crucial for these predictions. According to numerical simulations of galaxy formation, the former assumption seems to fail (e.g., Navarro& Steinmetz 1997), posing a difficulty for the hierarchical formation scenario. Possible intermediate astrophysical processes might solve this difficulty (e.g., Weyl, Eke & Efstathiou 1998). Regarding the latter assumption, according to preliminar results of analysis of CDM cosmological N-body simulations, it seems to be reasonable, at least for the present-day isolated haloes (Avila-Reese, Klypin, Firmani, & Kravtsov, in preparation), although deeper studies are necessary, in particular for the inner regions of the halos whose collapse occured at early epochs.
### 4.2 Rotation curves
The shapes of the rotation curves depend on $`\lambda `$, on the MAH that determines the DM halo structure, and on $`f_d`$. In Fig. 4, we plot rotation curves corresponding to models of $`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$ for a variety of MAHs and spin parameters. From Fig. 4, the main tendencies appear clearly: the rotation curve shapes are more peaked as $`\lambda `$ is smaller or the MAH is more active at early epochs (the DM halo is more concentrated). Due to the correlation between SSD and $`\lambda `$, a correlation between the SSD (SB) and the rotation curve shape is natural in our models: high SSD galaxies typically present more peaked rotation curves than low SSD galaxies. Observational studies seem to find a similar trend with SB (e.g. Casertano & van Gorkom 1991; Verheijen 1997). Although less significant, observations also show a dependence of the rotation curve shape on luminosity (Persic & Salucci 1988; Casertano & van Gorkom 1991; Persic, Salucci, & Stel 1996). Our models tend to confirm this dependence (see Avila-Reese & Firmani 2000). Finally, for high values of $`f_d`$, the shapes of our rotation curves are more peaked (Fig. 5).
As it was previously pointed out by MMW98, we find that the minimum spin parameter, $`\lambda _{\mathrm{min}}`$, for which the rotation curves are still realistic (nearly flat), should be increased as $`f_d`$ increases. For instance, an eye inspection of panels (a) and (d) of Fig. 5 suggests that when $`f_d=0.05`$ the $`\lambda _{\mathrm{min}}`$ can be as low as 0.03, while $`\lambda _{\mathrm{min}}`$ probably should be increased (to $`0.04`$) when $`f_d=0.08`$. Dalcanton et al. (1997) and MMW98 proposed that models with $`\lambda <\lambda _{\mathrm{min}}`$ are gravitationally unstable. For example, according to the predicted distribution of $`\lambda `$, this means that if $`\lambda _{\mathrm{min}}0.04`$, then roughly one third of the galaxies have unstable discs (S0’s, ellipticals?), a fact that disagrees with the observations in the local universe. This problem is enhanced if some angular momentum is transfered during the gas collapse, as is seen in the N-body+hydrodynamics simulations. This difficulty could be attenuated if the DM haloes would have inner density profiles shallower than $`r^1`$ (see $`\mathrm{\S }6`$).
### 4.3 Rotation curve decompositions
In Fig. 5 we show the decompositions of the rotation curves into disc (stars+gas) and halo (contracted by the baryon collapse) components for models of $`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$ and different values of $`\lambda `$ and $`f_d`$. As is seen in Fig. 5, the maximum of the disc rotation velocity is attained at $`2.2R_s`$. The disc/halo decompositions of the observed rotation curves do not have a single solution, and an additional constraint is required in order to reproduce them. The “minimum halo” (or “maximum disc”) solution (c.f., Carignan & Freeman 1985; Sancisi & van Albada 1987) has been the commonly adopted restriction. In this case, the ratio between the disc and total maximum rotation velocities, $`V_d/V_t`$, is about $`0.85`$ (e.g. Sackett 1997). Nevertheless, as some observational works have pointed out, this solution seems to present some shortcomings or at least is not applicable to all galaxies (Persic & Salucci 1988; see Navarro 1998 for more references). From an analysis of stellar velocity dispersions Bottema (1993, 1997), concluded that $`V_d/V_t`$ is about 0.63 for high SB galaxies, and even less for low SB galaxies.
In Fig. 6, $`V_d/V_t`$ is plotted as a function of the SSD for 20 models from our catalog corresponding to $`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$ and $`f_d=0.05`$ (filled circles). The disc contribution increases as the galaxy SSD is greater. The models corresponding to high SSD galaxies ($`\mathrm{\Sigma }_{s,0}2002000M_{}/pc^2`$, where $`\mathrm{\Sigma }_{s,0}`$ is the central SSD) have on average $`V_d/V_t0.70`$. This is the approximate value of the average galaxy model with $`\lambda =0.05`$ for which the DM dominates at essentially any radius (Fig. 5). Some theoretical arguments (e.g. Athanassoula et al. 1987; Debattista & Sellwood 1998) and observational studies (Verheijen 1997; Corsini et al. 1998) suggest that in high SB galaxies the disc component should dominate to some extent at the most inner radii.
In our models, the dominance of the halo component over the disc component is mainly due to the steep inner density profile ($`r^1)`$ of the original DM halo. Therefore, the inferred rotation curve decompositions of observed galaxies seem to point out that the halo core should be shallower than $`r^1`$. More direct evidence for shallow halo cores comes from the rotation curves of dwarf and low SB galaxies (Moore 1994; Flores & Primack 1994; Burkert 1995). In $`\mathrm{\S }6`$ we shall explore models where shallow cores are artificially introduced to our DM haloes.
## 5 Structural and dynamical correlations of disc galaxies
In this section we present the $`M_sV_{\mathrm{max}}`$ (TF) and $`R_sV_{\mathrm{max}}`$ relations calculated from our catalogs of galaxy models. The scatter of the $`M_sV_{\mathrm{max}}`$ relation is studied in detail. We analise the correlation among the residuals of the above relations and we explain why low and high SB galaxies have the same TFR.
### 5.1 Infrared Tully-Fisher relations
Disc galaxies present a strong correlation between their luminosities $`L_i`$ ($`i`$ is the spectral band) and their maximum rotation velocities $`V_{\mathrm{max}}`$, commonly known as the TFR (Tully & Fisher 1977). In the infrared bands ($`i=I,H,K,\mathrm{})`$, this relation is given by:
$$L_i=A_iV_{\mathrm{max}}^{\mathrm{m}_\mathrm{i}}$$
(2)
where $`A_i`$ is related to the so called zero-point, and $`3\genfrac{}{}{0pt}{}{_<}{^{}}m_i<4`$ according to several observational studies. Since $`L_iM_s,`$ where $`M_s`$ is the disc stellar mass, eq. (2) may be interpreted as a relation between $`M_s`$ and $`V_{\mathrm{max}}`$. In Fig. 7, where $`M_s`$ is plotted as a function of $`V_{\mathrm{max}}`$, we present our results from the catalog constructed by Monte Carlo simulations with $`f_d=0.05`$ (see §2 and §3). The error bars represent the standard deviation calculated by adopting a normal distribution for the deviations of the velocity for a given mass; see also Table 1. In the range of masses considered here ($`M_s10^910^{11}M_{}),`$ the slope of the $`M_sV_{\mathrm{max}}`$ relation is approximately 3.4; this slope is slightly larger than the slope of the mass-velocity relation of the cosmological DM haloes ($`\mathrm{\S }2`$).
In Fig. 7 we have also included observational data. In the $`I`$ band we used the TFRs given by Giovanelli et al. (1997) and by Willick et al. (1995; they used the data published by Han, Mould and collaborators, see the references therein). In these studies the line widths were corrected for non-circular motions. In the $`H`$ band we used the TFRs given by Gavazzi (1993) and by Pelletier & Willner (1993). Although the $`H`$magnitudes of most of the galaxies reported in Gavazzi (1993) were obtained through aperture photometry, he used the total $`H`$magnitudes estimated with an extrapolation technique. We corrected the TFR given by Gavazzi for non-circular motions (see AFH98 for details). In the case of Pelletier & Willner, we used their TFR calculated with the total $`H`$magnitudes obtained with an infrared array and the line widths corrected for non-circular motions. Regarding the determination of $`V_{\mathrm{max}}`$ in the observational studies we plot in Fig. 7, most of the data were obtained from HI line-width measures. As Verheijen (1997) have shown, the differences in the slope of the TFR calculated with single dish and detailed synthesis data are small; if any, the slope is slightly shallower when using the detailed synthesis data.
In earlier works about the near infrared (particularly $`H`$band) TFR with aperture photometry, the slope obtained was $`4`$; however, as pointed out by Pierce & Tully (1988) and Bernstein et al. (1994), the use of aperture magnitudes results in an artificially large slope to the TFR. This is why in studies where $`CCD`$photometry is used the slope of the TFR in infrared bands resulted shallower than 4<sup>1</sup><sup>1</sup>1After the complexion of this paper a study by Tully & Pierce (1999) appeared where the authors carrefully re-evaluate observational data in order to accurately determine template TFRs in different bands. They conclude that there appears to be convergence in the infrared towards a TFR’s slope of $`3.4\pm 0.1`$ (see also Rohtberg et al. 1999) (e.g., Pierce & Tully 1988; Pelletier & Willner 1993; Bernstein et al. 1994; Verheijen 1997). The last author, for a sample of galaxies in the Ursa Major cluster obtained a $`K^{}`$band TFR with slope $`3.3`$ when he used the complete sample of 41 galaxies and with slope $`4`$ when he used only 15 galaxies out from the sample, selected to be unperturbed galaxies of late type Sb-Sd and without prominent bars.
In order to transform luminosities into stellar masses, a mass-to-luminosity ratio, $`\mathrm{{\rm Y}}`$, should be adopted. For the $`I`$band observations, we assumed $`\mathrm{{\rm Y}}_I=1.8\left(\frac{M_s}{5\times 10^{10}M_{}}\right)^{0.07}h`$ (see AFH98). This ratio is close to the one suggested by MMW98 ($`\mathrm{{\rm Y}}_I=1.7h`$ ) on the basis of the $`\mathrm{{\rm Y}}_B`$ that Bottema (1997) inferred from disc dynamics. For the $`H`$band observations $`\mathrm{{\rm Y}}_H=0.55`$ was used and $`h=0.65`$ was assumed. This mass-to-luminosity ratio is obtained from direct observational estimates in the solar neighborhood (Thronson & Greenhouse 1988; see details in AFH98).
In Fig. 7, the model results are slightly shifted to the high velocities with respect to the Giovanelli et al. (1997) and Gavazzi (1993) data, while they agree rather well with the Han-Mould (quoted by Willick et al. 1995) and Pelletier & Willner (1993) data. The assumed $`f_d`$ does not significantly change these results because our models typically shift along the main relation for different values of $`f_d`$ (see below). In order to reduce the uncertainties due to $`\mathrm{{\rm Y}}`$ we have used measured TFRs in two bands ($`I`$ and $`H`$) with $`\mathrm{{\rm Y}}`$ independently estimated. The tendency of the models to have a larger $`V_{\mathrm{max}}`$ for a given $`M_s`$ may again be showing that the CDM haloes are too cuspy. When a shallow core is introduced in the haloes (see §6), the resulting $`V_{\mathrm{max}}`$ of the galaxies are slightly smaller (empty circles in Fig. 7). Despite the uncertainties, we conclude that the agreement of the slope and the zero-point of the TFR between observations and our theoretical results is reasonable well. Note that in the theoretical calculations, both the cosmological framework and the SF process were taken into account.
Our results show that the TFR is mainly a product of the mass-velocity relation of the DM haloes and, as stated in $`\mathrm{\S }2`$, this is determined by the cosmological initial conditions (e.g. Frenk et al. 1988, Cole et al. 1994, MMW98, AFH98, Steinmetz & Navarro 1998). The slope of the TFR is linked to the shape of the power spectrum at galaxy scales and the MAH of the halo. For most of the CDM models, the power spectrum shape at galaxy scales and the MAHs are almost the same, so the slope of the TFR is expected to be a generic feature of the CDM cosmogony (Firmani & Avila-Reese 1999). Concerning the zero-point, from the point of view of the cosmological model it depends on the amplitude of the power spectrum at galaxy scales. Models with low amplitudes at galaxy scales —like the $`\mathrm{\Lambda }`$CDM or open CDM models— predict the zero-point of the TFR better than those with high amplitudes (for example, the standard CDM model; see also Firmani & Avila-Reese 1999; Jimenez & Heavens 1999).
The results presented in Fig. 7 were obtained assuming $`f_d=0.05`$. We find that the $`M_sV_{\mathrm{max}}`$ relation is rather insensitive to changes in the disc mass fraction, particularly when this fraction is such that the disc contribution to the rotation curve is relevant. In the right lower corner of Fig. 7 the solid line represents the shifts that a typical galaxy model (average MAH, $`\lambda =0.05`$) suffers when varying $`f_d`$. From right to left the dots correspond to models with $`f_d=0.03`$, 0.05 and 0.08, respectively. For the typical galaxy models, an increment in the parameter $`f_d`$ basically produces a shift along the $`M_sV_{\mathrm{max}}`$ relation. This is due to the “compensating” action of the gravitational pull exerted by the disc on the DM halo. This result also suggests that, if $`f_d`$ fluctuates or has a dependence upon mass, by effects of gas cooling or feedback, then the TFR is almost unaffected, at least in the range mentioned above. When $`f_d`$ is very small, the disc contribution to the rotation velocity is negligible, and $`V_{\mathrm{max}}`$ is determined by the DM halo alone. In this case, the $`M_sV_{\mathrm{max}}`$ relation becomes sensitive to $`f_d`$ because the stellar mass (luminosity) and the galaxy dynamics are de-coupled.
Based on the results presented above and as was pointed out in Firmani, Hernández, & Avila-Reese (1997), we find that the infrared TF relations offer a robust way to normalize the power spectrum of fluctuations at galaxy scales independently of uncertainties due to the disc mass fraction assumed. This normalization favours power spectra corresponding to $`COBE`$-normalized low density CDM models ($`\mathrm{\Lambda }`$CDM or open CDM models).
### 5.2 Scatter of the Tully-Fisher relations
The observational rms scatter of the TFR is small. In the $`I`$-band Giovanelli et al. (1997) reported a total scatter of $`0.3`$ mag for rotators of $`180`$ km/s; they found that the scatter increases from fast to slow rotators. Willick et al. (1995, 1996) and Mathewson and Ford (1994) estimated a total scatter of 0.38-0.43 mag and 0.44 mag, respectively. Bernstein et al. (1994) found a scatter of 0.23 mag. In the $`H`$-band Willick et al. (1996) found a scatter of 0.47 mag. These last authors concluded that the estimated intrinsic scatter of the infrared TFRs is not smaller than 0.3 mag. Are the theoretical models able to predict the intrinsic rms scatter inferred from the observations? In our models, there are at least two sources of the scatter in the infrared TFR. One is related to the statistical nature of the MAHs. The dispersion in the MAHs determines a dispersion in the maximum circular velocities of the virialized structures (Eisenstein & Loeb 1996; AFH98; see $`\mathrm{\S }2`$). The second source is associated with the dispersion of the spin parameter $`\lambda `$ (e.g. MMW98). As $`\lambda `$ decreases, the disc becomes more concentrated, and therefore $`V_{\mathrm{max}}`$ increases.
A first qualitative estimate of the scatter produced by the MAH and $`\lambda `$ distributions may be appreciated by the shifts from the $`M_sV_{\mathrm{max}}`$ relation shown in the right lower corner of Fig. 7. The dotted and dashed lines correspond to different MAHs and $`\lambda ^{}s`$. Dispersions in both the MAH and $`\lambda `$ tend to shift the models along the $`M_sV_{\mathrm{max}}`$ relation, in a similar way as $`f_d`$. In columns (3) and (4) of Table 2, we present the total scatters in velocity and mass of the $`M_sV_{\mathrm{max}}`$ relation obtained from our catalogs for three stellar disc masses. The scatter in mass is expressed in magnitudes (see $`\mathrm{\S }2.3`$). There is a marginal agreement with the observations. Observational and theoretical uncertainties are large, so it is still premature to draw any conclusions relating the scatter of the TFR and the cosmological initial conditions (but see Eiseinstein & Loeb 1996). It should be noted that we have maximized the scatter in our models by including both low and high SB galaxies, and taking into account major mergers in the calculation of the MAHs. The observational estimates refer mostly to high SB late-type galaxies. On the other hand, we have not included any source of scatter related to non-stationary SF or possible variations in $`f_d`$. Observational scatter due to non-stationary SF and $`f_d`$ would imply a larger scatter in our models, although Elizondo et al. (1998) have shown that, in this case, compensating effects may play an important role in keeping, or even decreasing, the scatter of the TFR.
We calculated a set of models with a constant $`\lambda =0.05`$, in order to estimate the scatter of the $`M_sV_{\mathrm{max}}`$ relation, due to variations in the MAHs. In order to estimate the contribution to the scatter of the $`M_sV_{\mathrm{max}}`$ relation due to $`\lambda `$, we calculated a set of models using the average MAH and taking the different $`\lambda `$ values from its log-normal distribution. We find that the influence of $`\lambda `$ on the scatter of the $`M_sV_{\mathrm{max}}`$ relation is smaller than that due to the MAHs; the quadratic contributions of the scatter in $`\lambda `$ and in MAH to the total scatter, are roughly 25 and 75 per cent, respectively.
### 5.3 Residuals of the $`M_sV_{\mathrm{max}}`$ and $`M_sR_s`$ relations and the TFR of low SB galaxies
In our models $`\lambda `$ strongly influences the SSD. Consequently, the small contribution of $`\lambda `$ to the scatter of the $`M_sV_{\mathrm{max}}`$ relation, implies that galaxies of different SB should have almost the same TFR. This agrees with the observations that show that the TFRs of low and high SB galaxies are the same (e.g. Zwaan et al. 1995; Tully & Verheijen 1997; Verheijen 1997).
Recently, Courteau & Rix (1998), using large catalogs of late-type, high SB galaxies, have studied the correlations among the residuals of the TF and the luminosity-radius relations. They find that the slope of the correlation among the residuals, $`V_{2.2}/R_s`$, has a mean value of $`0.19\pm 0.05`$ ($`V_{2.2}V_{\mathrm{max}}`$ is the value of the rotation velocity at $`2.2R_s`$). According to Courteau & Rix (1998), this means that the TFR scatter correlates only slightly with disc size or SB. They interpret this as a evidence for large amounts of DM in the inner parts of late-type galaxies ($`R<R_s`$). If the DM is dominant, the disc component plays almost not role in setting $`V_{\mathrm{max}}`$. This conclusion, however, might not be consistent with observations that hint that shape of the rotation curve correlates with SB for a given luminosity (Casertano & van Gorkom 1991; Verheijen 1997).
Using our models we can explain this apparent observational inconsistency. In order to interpret the observational results of Courteau & Rix (1998) one must take into account the difference in the SF histories among galaxies of different SB. From our models, in Fig. 8 we plot the deviations from the $`M_sV_{\mathrm{max}}`$ relation, $`\delta `$log$`V_{\mathrm{max}}`$, vs. the deviations from the $`M_sR_s`$ relation, $`\delta `$log$`R_s`$. The model galaxies were divided into three groups depending on their central SSD: very high SSD ($`\mathrm{\Sigma }_{s,0}>2000M_{}/pc^2`$), high SSD ($`200M_{}/pc^2<\mathrm{\Sigma }_{s,0}<2000M_{}/pc^2`$), and low SSD ($`\mathrm{\Sigma }_{s,0}<200M_{}/pc^2`$); we used black, gray, and empty circles to represent them. The slope of the correlation among the residuals changes with the SSD of the galaxy. The very high SSD models probably do not represent realistic disc galaxies because their rotation curves are too peaked. These models could be subject to instabilities that destroy the disc (see $`\mathrm{\S }\mathrm{4.0.2}`$). For the high SSD models the slope is approximately $`0.15`$ in agreement with the results of Courteau & Rix (1998). In order to compare the models with observations that include a wide range of SB’s, we estimated the residuals $`\delta `$log$`V_{\mathrm{max}}`$ and $`\delta `$log$`R_s`$ using the data given in Verheijen (1997) and Tully et al. (1996) for galaxies of the Ursa Major cluster. Fig. 8b shows these residuals where filled and empty circles represent high and low SB galaxies, respectively. In spite of the small number of objects, the tendencies in this plot are similar to those of our models. The low SB galaxies specifically demonstrate that there is no extrapolation of the trend found by Courteau & Rix (1998) to large radii (low SBs).
From Fig. 8a, it is seen that for galaxies with high SSD (small $`\delta `$log$`R_s`$) the deviates from the $`M_sV_{\mathrm{max}}`$ relation decrease toward the low velocity side as the SSD decreases. However, for low SSD galaxies (large $`\delta `$log$`R_s`$), this behavior is reversed. This result is due to the fact that not only $`V_{\mathrm{max}}`$ decreases as the SSD decreases, but so does the stellar disc mass, $`M_s`$. For this reason, for a fixed $`M_s`$ increasing the scale length $`R_s`$ (decreasing the SSD), two regimes are present in Fig. 8a. When the SSD is large ($`\mathrm{\Sigma }_{s,0}\genfrac{}{}{0pt}{}{_>}{^{}}200M_{}/pc^2`$), the disc contribution to the circular velocity becomes important and the residuals of the $`M_sV_{\mathrm{max}}`$ relation decrease with the residuals of the $`M_sR_s`$ relation ($`V_{\mathrm{max}}`$ decreases with $`R_s`$). But, when the SSD is small ($`\mathrm{\Sigma }_{s,0}\genfrac{}{}{0pt}{}{_<}{^{}}200M_{}/pc^2`$) the fraction of gas converted into stars is smaller than in discs with larger SSDs. Thus, in the last case $`M_s`$ should correspond to a galaxy with a more massive halo than in the former case. Therefore, the residuals of the $`M_sV_{\mathrm{max}}`$ relation increase with the residuals of the $`M_sR_s`$ relation ($`V_{\mathrm{max}}`$ increases with $`R_s`$). At the same time, when the SSD is large the rotation curve is dominated by the disc and it results peaked, while when the SSD is small, the rotation curve is dominated by the halo contribution, it slowly increases, and presents a broad maximum.
We conclude that the dependence of $`M_s`$ upon the SSD is responsible for the almost flat and non-monotonic slope of the correlation among the residuals of the $`M_sV_{\mathrm{max}}`$ and $`M_sR_s`$ relations (Fig. 8a,b). One expects that this slope becomes steeper and monotonic if the total (stars+gas) disc mass $`M_{\mathrm{tot}}`$ is used instead of $`M_s`$. Since $`M_{\mathrm{tot}}`$ does not depend upon the disc surface density, then as the disc surface density increases (i.e. the disc scale length $`R_{\mathrm{tot}}`$ decreases), $`V_{\mathrm{max}}`$ increases for all the SB’s. The residuals from the $`M_{\mathrm{tot}}V_{\mathrm{max}}`$ and $`M_{\mathrm{tot}}R_{\mathrm{tot}}`$ relations, for the same models from Fig. 8a are plotted in Fig. 8c. As can be seen from Fig. 8c, if we assume that the disc stellar mass (luminosity) is exclusively proportional to the total disc mass (e.g. Dalcanton et al. 1998; MMW98; AFH98; van den Bosch 1998), then we arrive to an incorrect result: the slope of the correlation among the residuals of the TF and luminosity-radius relations is steep and monotonic, or in other words, the deviations from the TFR is highly correlated with the SB. This contradicts the observational inferences of Courteau & Rix (1998), and that based upon the data from Verheijen (1997) and plotted in Fig. 8b.
In our model galaxies, the stellar mass $`M_s`$ depends upon the disc surface density, because the SF efficiency depends upon the disc surface density. The stellar mass in our models, also depends upon the gas accretion rate given by the MAH. In Fig. 9 we have plotted the fraction of gas ($`f_g=M_{gas}/(M_{gas}+M_s))`$ as a function of the central SB in the $`B`$ band, for our models and for the data compiled by de Blok & McGaugh (1997). The agreement between models and observations is rather good, suggesting that our SF histories are realistic. We conclude that the disc surface density plays an important role in determining the gas fraction, $`f_g`$, of disc galaxies.
The analysis presented in this subsection allows one to understand why high and low SB galaxies have the same TFR, even though the shapes of their rotation curves depend upon the SB. In the $`M_s`$ vs. $`V_{\mathrm{max}}`$ plot (Fig. 7), as the SB decreases, the galaxies not only shift towards lower velocities, but also shift towards a smaller stellar mass (luminosity). Thus, the deviations from the $`M_sV_{\mathrm{max}}`$ relation (TFR) are almost independent of the disc size or SB. This effect also explains why the scatter about the $`M_sV_{\mathrm{max}}`$ relation due to the dispersion of the parameter $`\lambda `$ becomes so small (see $`\mathrm{\S }5.2`$).
### 5.4 Disc sizes vs. maximum rotation velocities
The distribution of disc galaxies in the $`R_dV_{\mathrm{max}}`$ diagram is shown in Fig. 10a. The data were taken from a catalog of late-type galaxies in the $`r`$band elaborated by Courteau (1996,1997) (circles), and from a sample of galaxies of the Ursa Major cluster in the $`K^{}`$band elaborated by Verheijen (1997) (triangles). We transformed the surface brightness in the $`r`$ and $`K^{}`$bands to surface densities using $`\mathrm{{\rm Y}}_r=1.4`$ and $`\mathrm{{\rm Y}}_K^{}=0.6`$. We have tentatively divided the samples into two groups: high SB ($`\mathrm{\Sigma }_{s,0}\genfrac{}{}{0pt}{}{_>}{^{}}200M_{}/pc^2`$, filled symbols) and low SB ($`\mathrm{\Sigma }_{s,0}\genfrac{}{}{0pt}{}{_<}{^{}}200M_{}/pc^2`$, empty symbols). The threshold we used for the division in Courteau’s (1996,1997) sample was $`\mu _0=20.5`$ R-mag/arcsec<sup>2</sup> (low SB galaxies are under-represented in this sample). For the Verheijen (1997) sample the threshold we used was 18.5 $`K^{}`$mag/arcsec<sup>2</sup>. Despite the incompleteness of the samples and the differences between them, Fig. 10a indicates that the disc size correlates with the maximum rotation velocity, particularly for families of similar SB. In Fig. 10b we plot the results from our model catalogs corresponding to three masses ($`M_{\mathrm{nom}}=5\times 10^{10}M_{}`$, $`5\times 10^{11}M_{}`$, and $`5\times 10^{12}M_{}`$). As was noted in the previous subsection, the very high SB discs with $`\mathrm{\Sigma }_{s,0}>2000M_{}/pc^2`$ ($`\mu _0<16`$ K’-mag/arcsec<sup>2</sup>; black filled circles) probably are not realistic; they might be subject to gravitational instabilities (these models have $`\lambda <0.025`$). The solid lines are the linear regressions for the normal and low SSD models (gray and empty circles, respectively). These lines are reproduced in panel (a) for comparison with the observations. The agreement is reasonably good.
In the lower right corner of Fig. 10b, we show how galaxy models ($`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$) shift in the $`R_dV_{\mathrm{max}}`$ plane when the MAH, $`\lambda `$, and $`f_d`$ are varied. The largest shift is due to $`\lambda `$, which is why the $`R_dV_{\mathrm{max}}`$ relation is well defined only for discs of similar SSD (SB).
## 6 Effects of a shallow core
From the observed rotation curves of some dwarf and low SB galaxies, Kravtsov et al. (1998, hereafter KKBP98) empirically inferred an approximate self-similar density profile for the haloes of these galaxies assuming that they are completely dominated by DM:
$$\rho (r)=\frac{\rho _0}{(r/r_0)^\gamma [1+(r/r_0)^\alpha ]^{(\beta \gamma )/\alpha }},$$
(3)
with $`(\alpha ,\beta ,\gamma )=(2.0,3.0,0.2)`$, where the particular value of $`\gamma =0.2`$ is weakly motivated; it should be considered just as the evidence of a shallow core ($`\gamma 0`$, see also Burkert 1995). Kravtsov et al. concluded that these density profiles are in reasonable agreement with those obtained in their high-resolution cosmological N-body simulations. Other high-resolution simulations, however, did not confirm the numerical results of KKBP98 (Fukushige & Makino 1997; Moore et al. 1998, 1999; Jing 1999), posing a potential difficulty for the CDM models with respect to the observations. The numerical, physical or observational analysis of the inner structure of galactic DM halos is beyond the scope of the present paper. Nevertheless, through our models, we shall explore the effects of a possible shallow core on the properties of the disc galaxies.
We shall introduce artificial shallow cores to our haloes. We look for that the scaling parameters of the cores to be introduced are in rough agreement with the rotation curves of low SB galaxies, taking into account the contraction of their haloes after disc formation. Since the density profile given by eq. (3) is in rough agreement with that observed low SB rotation curves suggest, it is reasonable to use this profile for the galaxy halos. As a matter of fact, the shape of this profile is similar to the typical density profile of our models (and to the NFW profile) except for the inner regions where $`\rho (r)r^{0.2}`$. Therefore, we deform ad hoc the inner profile of our DM haloes by smoothly imposing the $`\rho (r)r^{0.2}`$ behavior from a given radius $`r_{\mathrm{core}}=\nu r_{\mathrm{max}}`$ down to the center, where $`r_{\mathrm{max}}`$ is the radius at $`V_{\mathrm{max}}`$ and $`\nu <<1`$. Because our approach is evolutionary, we need to introduce the mentioned deformation at all the times. Unfortunately, we only have information at $`z=0`$, therefore, for all the other epochs we use $`r_{\mathrm{core}}`$ with the same $`\nu `$ defined at $`z=0`$. The fittings applied by KKBP98 to low SB galaxies and to the haloes obtained in their simulations show that the two scaling parameters of the profile (3), $`r_0`$ and $`\rho _0`$, are linked. This allows us to roughly fix $`r_0`$ given other parameter as $`\rho _0`$ or $`V_{\mathrm{max}}`$ (see also Burkert 1995). With the formation of a disc —even if this is of low surface density— the original mass distribution of the dark halo changes. For instance, a model with the average MAH and $`\lambda =0.08`$ (low SB model) where the shallow core given by ec. (3) whith $`r_0`$ normalized to observations is introduced, after disc formation presents a too small $`r_0`$ for its $`V_{\mathrm{max}}`$. Now, if we increase $`\nu `$ for the initial halo, then at some value the parameters $`r_0`$ and $`V_{\mathrm{max}}`$ of the system after disc formation will agree with the parameters estimated for the observed low SB galaxies. The agreement occurs when $`r_0`$ increases roughly by a factor 1.4 with respect to the the KKBP98 estimates.
The main influence of the introduction of a shallow core in the DM haloes appears on the dynamics of the galaxy system. The same models presented in Figure 5 were again calculated but with the inclusion of a shallow core in the DM haloes. This core was introduced as described above: we deformed our DM haloes in order their density profiles fitted eq.(3) with the scaling parameters inferred by KKPB98 from rotation curves of low SB galaxies but making the core size ($`r_0`$) 1.4 times larger in order to “return back” the halo to its initial structure before disc formation. The obtained model rotation curve decompositions are plotted in Figure 11. The disc component contribution to the total rotation curve is now more significant than in the models plotted in Fig. 5, and is in better agreement with several observational and theoretical studies (see $`\mathrm{\S }4.3`$). The $`V_d/V_t`$ ratios for the catalog models of $`M_{\mathrm{nom}}=5\times 10^{11}M_{}`$ calculated with the modified DM haloes are plotted in Fig. 6 (empty circles). It is also important to note that the inclusion of a shallow core helps to obtain nearly flat rotation curves for models with small values of $`\lambda `$ and/or high values of $`f_d`$ that, otherwise, would have too steep rotation curves (see the end of $`\mathrm{\S }4.2`$).
The inclusion of a shallow core in the DM haloes does not produce considerable changes in the SSD, scale lengths, and gas fractions of the models. Some influence appears on $`V_{\mathrm{max}}`$. The $`M_sV_{\mathrm{max}}`$ relation slightly shifts to the low velocity side (see Figure 7, empty circles), being even in better agreement with the estimates from the observations than in the case of the models without core. On average, the amplitude of this relation increases $`\mathrm{\Delta }`$log $`M_s0.15`$. The shallow core refers only to the most inner regions of the halo-galaxy system. Therefore, its influence on the $`V_{\mathrm{max}}`$ is small because the $`V_{\mathrm{max}}`$ of the original halo (without gravitational contraction due to disk formation) is typically attained at radii larger than the disc size, i.e. far from the shallow core region.
Our models, which include the adiabatic contraction of the DM due to disc formation, suggest that (i) the original haloes of the present-day low SB galaxies had to have a shallow core larger by roughly a factor 1.4 than that their rotation curves suggest according to KKBP98, and (ii) the dynamics of normal (high SB) disc galaxies could be better explained if the original DM haloes had such a core. The inner structure of the galaxy dark haloes offers an important test for structure formation theories (Moore 1994; Flores & Primack 1994; Burkert 1995; Moore et al. 1999). More observational efforts are necessary in this direction.
## 7 Summary and conclusions
We have studied the formation and evolution of disc galaxies in a $`\mathrm{\Lambda }`$CDM<sub>0.35</sub> cosmology. We constructed a self-consistent model trying to avoid free parameters. Our main assumptions were (i) spherical symmetry and adiabatic invariance during the gravitational collapse of the DM, (ii) aggregation of the baryon matter to the disc in form of gas (no merger) with a rate given by the cosmological aggregation rate, (iii) detailed angular momentum conservation and adiabatic invariance during the gas collapse, and (iv) stationarity and self-regulation of SF in the disc. The obtained different density profiles of DM haloes and the dispersion agree with the results of cosmological N-body simulations (Avila-Reese et al. 1999). The most typical profiles are well described by the NFW profile. The properties of the model galaxies depend upon three initial factors: mass, MAH, and spin parameter, $`\lambda `$. The results allow us to predict and to understand several observational features of disc galaxies:
(1) Within the observational and theoretical uncertainties, we find that the slope and zero-point of the TFR in the $`I`$ and $`H`$bands may be directly determined by the cosmological initial conditions, principally the power spectrum of fluctuations of the CDM models. The $`M_sV_{\mathrm{max}}`$ relation (TFR) remains the same for different disc mass fractions $`f_d`$, in the cases where the disc makes a non-negligible gravitational contribution to the total rotation curve (when $`f_d\genfrac{}{}{0pt}{}{_>}{^{}}0.03`$ for the cosmological parameters used here). The TFR can be used to normalize the power spectrum at galaxy scales independently of uncertainties due to the assumed $`f_d`$. $`COBE`$-normalized low density CDM models are favoured by this normalization.
(2) The rms scatter in the TFR, according to our models, is produced by the scatter in the halo formation histories (MAHs) and by the dispersion of $`\lambda `$. As a result of compensating effects, the quadratic contribution of the latter is only a 25 per cent in the total rms scatter. Thus, a major contribution to the TFR scatter comes from the stochastic nature of the protogalaxy MAHs. The total scatter we obtain does not disagree with the observational data, but, owing to the observational and theoretical uncertainties, it is still premature to claim definitive conclusions.
(3) We explain why high and low SB galaxies show approximately the same TFR, and why the slope of the correlation among the residuals of the TF and luminosity-radius relations reported by Courteau & Rix (1998) is so small. We obtain similar results to those of Courteau & Rix although the shape of the rotation curves of our models correlates with the SSD and the rotation curves are not strongly dominated by the DM component. For models with a given total (star+gas) disc mass, as the SSD decreases, $`V_{\mathrm{max}}`$ decreases, but, owing to the dependence of the SF rate on the disc surface density, the stellar mass $`M_s`$ also decreases. Thus, models of different surface density lie on the same relation in the $`M_sV_{\mathrm{max}}`$ plane. Indeed, as our models and the observational data show, the disc gas fraction ($`f_g=M_g/(M_g+M_s)`$) strongly correlates with the SSD (SB).
(4) For a given $`f_d`$, the shapes of the rotation curves are steeper as $`\lambda `$ decreases or as the MAH is more active at early epochs (the DM halo is more concentrated). The SSD depends upon $`\lambda `$ and the MAH, which explains why the shape of the rotation curves depends upon the SSD, as the observations indicate. If $`f_d`$ is too high ($`\genfrac{}{}{0pt}{}{_>}{^{}}0.08`$), the rotation curves of discs with small $`\lambda `$ ($`<\lambda _{\mathrm{min}}0.04)`$ decrease too fast, and these discs are probably unstable. If $`f_d`$ is too small ($`\genfrac{}{}{0pt}{}{_<}{^{}}0.03`$), then low and high SSD galaxies will have very different TFRs, contrary to observations.
(5) The rotation curve decompositions show a dominance of the DM component down to the very central regions for most of the models with $`f_d0.05`$. This occurs because the inner density profile of the DM haloes are steep ($`\rho (r)r^1`$). The $`V_d/V_t`$ ratio increases with the SSD. On average, $`V_d/V_t0.70`$ for the high SSD models.
(6) The discs in centrifugal equilibrium that form in the centre of evolving CDM haloes with $`\lambda `$ constant in time, and with an accretion rate dictated by the hierarchical MAH, have a nearly exponential SSD distribution. The central SSD strongly depends on $`\lambda `$, and less on the mass and the MAH.
(7) We have studied the effects a shallow core in the DM haloes would produce on the galaxy properties. We found that the rotation curve decompositions of high SB galaxies agree better with the decompositions inferred from observations when the haloes have a shallow core. Because the rotation curve with a shallow core tends to be flatter, the minimum possible value of $`\lambda `$ and the fraction of unstable discs decrease. On average, $`V_d/V_t0.76`$ for the high SSD models. The introduction of shallow cores slightly shifts the models in the $`M_sV_{\mathrm{max}}`$ plane towards lower velocities, improving the agreement with the estimates inferred from the $`I`$ and $`H`$band TFRs.
The disc galaxy evolution models presented here make use of several ingredients of the CDM-based hierarchical formation scenario. We assumed that the hierarchical aggregation of mass takes place as a gentle accretion process, discarding major mergers. Although a realistic galaxy formation model should take into account both accretion and merging, in the case of disc galaxies, the former could not have played an important role. One expects that the aggregation of baryon matter was more uniform and gentle than that of the DM due to its hydrodynamical properties and due to the re-heating at high redshifts (e.g. Blanchard, Valls-Gabaud, & Mamon 1992). Thus, even if the DM haloes suffered an active merging process, luminous galaxies could have formed within large haloes in the way envisaged in the extended collapse picture. Our results showed that several structural and dynamical properties of disc galaxies and their correlations are closely related to the cosmological background, whereas other are consequence of evolutionary processes. Among the former, the disc density and velocity profiles, and the TFR are remarkable and they agree with the observations for the $`\mathrm{\Lambda }`$CDM<sub>0.35</sub> universe used here. However, we should also emphasize that the inner structure of the CDM haloes and probably the scatter of the TFR, are in conflict with observations. Both are probably associated with the statistical nature of the primordial fluctuation field (Gaussian?) and/or to the nature of the DM particles.
## Acknowledgments
We thank Julieta Fierro, Michael Richer, and Xavier Hernández for helpful comments and for critically reading the original manuscript. We also thank S. Courteau for providing the data shown in Fig. 10a in electronic form and F. Angeles for computing assistance. We are grateful to the announymous referee for comments which were very helpful to improve the quality of the paper. |
warning/0001/hep-ph0001255.html | ar5iv | text | # 1 Introduction
## 1 Introduction
During the last few years a number of hard diffractive processes were suggested for probing the short-range structure of hadrons and the behaviour of perturbative QCD (pQCD) at high energies. The description of these processes is greatly simplified if a factorization theorem is valid which allows to write the amplitude as a convolution of the wave function of the produced meson, a hard interaction block and a block related to the density of partons in the target.
The key element of the factorization proof for exclusive meson production in DIS is the selection of processes with longitudinally polarized incomming photons. This ensures that small $`1/Q`$ transverse interquark distances dominate in the $`\gamma _L^{}meson`$ transition and that extra soft interactions between consitituents moving in the same direction as the $`\gamma _L^{}`$ and those moving in the same direction as the nucleon are suppressed by a factor $`1/Q^2`$. As a result the study of exclusive vector meson production can be used for probing “small dipole” - hadron interactions at high energies.
Another appealing alternative is to use hard diffractive processes with scattering of the colliding particles at large enough $`t`$ to reach the so called sqeezed regime . Here the key question is whether indeed the sqeezing occurs, i.e. whether multiple interactions can be neglected. This is the case for $`\gamma _L^{}+pV+X`$ . However the rates in this case are low. The rates are much higher for the real photon case $`\gamma +pV+X`$.
The photoproduction of $`J/\mathrm{\Psi }`$ meson at large $`t`$ was studied in the papers . The predictions of pQCD for photoproduction of longitudinally polarized light meson at large $`t`$ were derived in refs. . Photoproduction of transversely polarized meson was discussed in a phenomenological model in .
Here we will derive pQCD predictions for all helicity amplitudes of this process $`M_{\lambda _1\lambda _2}`$, where $`\lambda _1=\pm `$ is the photon helicity and $`\lambda _2=\pm ,0`$ is the vector meson helicity. We calculate the factorizing contribution of small quark-antiquark separation and indicate the non-factorizing contributions appearing in our calculations as singularities at the end point in momentum fraction. We shall show that the chiral-odd configuration (the helicities of the quark and the antiquark are parallel) in the quark loop of Fig.1 gives a very important contribution if $`t`$ is not asymptotically large. Within usual perturbation theory a photon can split only into a chiral-even (the helicities of the quark and the antiquark are antiparallel) massless quark pair. The violation of chiral symmetry, which is well known to be a soft QCD phenomenon, generates a nonperturbative chiral-odd component of the real photon wave function. The interaction of this additional chiral contribution can however be described in pQCD since high $`t`$ quark-dipole scattering chooses a $`q\overline{q}`$ configuration with small transverse interquark distances. As a result the chiral-odd contribution can be factorized into a convolution of two nonperturbative photon and vector meson light-cone wave functions with the hard scattering amplitude. The chiral-odd wave function of the real photon is proportional to the quark condensate and its magnetic sucseptability . We shall show that the helicity amplitudes are very sensitive to the values of these fundamental parameters of the QCD vacuum. Let us stress that this new hard production mechanism for high $`t`$ photoproduction was not considered before.
We shall demonstrate that in the spin non-flip amplitude $`M_{++}`$ both the chiral-even and the chiral-odd mechanisms give contributions $`t^2`$ which are of the same sign. The contributions of these two mechanisms to the single spin-flip amplitude $`M_{+0}`$ are of opposite signs, they behave as $`t^{3/2}`$ for the chiral-even and $`t^{5/2}`$ for the chiral-odd case. For the double spin–flip amplitude $`M_+`$ the contributions from both mechanisms are of opposite sign and $`t^2`$ for the chiral-even and $`t^3`$ for the chiral-odd case. Due to large numerical coefficient in front of the chiral-odd contributions, even in the case of $`M_{+0}`$ and $`M_+`$ the chiral-odd mechanism is very important in a wide region of intermediate $`t`$ despite the fact that it is formally $`1/t`$ suppresed. This observation could explain why onset of the asymptotic regime, namely the dominance of the $`M_{+0}`$ helicity amplitude was not observed experimentally at large $`t`$.
The cross section for a reaction with rapidity gap $`\eta _0`$ (see the diagram of Fig.1) can be related to those for the photoproduction of $`V`$ off a quark and a gluon via the gluon and quark densities in a proton $`G(x,t)`$ and $`q(x,t)`$ :
$`{\displaystyle \frac{d^2\sigma (\gamma pVX)}{dtdx}}`$ $`=`$ $`{\displaystyle \underset{f}{}}\left(q(x,t)+\overline{q}(x,t)\right){\displaystyle \frac{d\sigma (\gamma qVq)}{dt}}+`$ (1.1)
$`G(x,t){\displaystyle \frac{d\sigma (\gamma GVG)}{dt}};x={\displaystyle \frac{4p_{}^2}{s}}\mathrm{cosh}^2{\displaystyle \frac{\eta _0}{2}}.`$
$`\eta _0`$ is the difference in rapidity between the struck parton and produced meson. (We consider here the case of small angle scattering so that $`t/xs1`$).
The factorization formula (1.1) is valid if the typical transverse distances between quarks in the upper part of Fig. 1 are small. In this case the contribution of the additional soft interactions which are schematically depicted in Fig. 2 will be power suppressed and therefore the jet balancing the high transverse momentum of the meson would be produced close to the gap edge.
We shall show below for each helicity amplitude of the $`\gamma qVq`$ process that there exist soft non-factorizable contributions originating from the region of large transverse distances between quark and antiquark in Fig. 1. For such soft contributions the factorization (1.1) breaks down and the rapidity $`\eta `$ of the leading jet for the system $`X`$ can lie far away from the gap edge, $`\eta >>\eta _0`$.
We shall show that in the region of intermediately large $`t`$ the relative contribution of these soft nonfactorizable interactions is not large numerically. Therefore in this $`t`$ region eq. (1.1) should be a good approximation and we expect that the leading jet of $`X`$ should be close to the gap edge. It would be very interesting to study the jet structure of the system $`X`$ experimentally since this is a clean signature for the dominance of the hard factorizable production mechanism.
We shall discuss now the helicity amplitudes of the parton subprocess $`\gamma qVq`$ in detail. The two-dimensional polarizations vectors of a real photon and transversly polarized vector meson are denoted as
$$𝐞^{(\pm )}=\frac{1}{\sqrt{2}}(1,\pm i).$$
(1.2)
There are three independent helicity amplitudes which we choose as
$$M_{++}(M_{}=M_{++}),M_{+\mathrm{\hspace{0.33em}0}}(M_{\mathrm{\hspace{0.33em}0}}=M_{+\mathrm{\hspace{0.33em}0}}),M_+(M_+=M_+).$$
Although in the following we calculate the helicity amplitudes only for the production of $`\rho ^0`$ meson, the resulting formulas are valid - after a suitable change of the coupling constant - also for $`\varphi `$ and $`\omega `$ meson production.
We shall use the results of , for the to describe the $`\rho `$ and real photon $`q\overline{q}`$ light-cone wave functions at small interquark distances.
In the present letter we present only the main steps of our calculations and the main results. Technical details of the derivation as well as a comparison with data will be presented in a later publication.
We use the following definitions
$$\gamma _5=\gamma ^5=i\gamma ^0\gamma ^1\gamma ^2\gamma ^3,\sigma ^{\mu \nu }=\frac{i}{2}[\gamma ^\mu ,\gamma ^\nu ],\epsilon ^{0123}=1$$
(1.3)
The quark charge is denoted by $`eQ_q`$. The flavour structure of the $`\rho ^0`$ vector meson can be described by the replacement $`eQ_qe/\sqrt{2}`$.
## 2 Factorization formulae for $`\gamma qVq`$ amplitude
The amplitude of the $`\gamma qVq`$ process can be written as a convolution of the hard scattering amplitude which describes the production of a quark pair $`A_{\alpha \beta }(l,pl)`$ and the amplitude of the transition of this quark pair into a meson $`\mathrm{\Psi }_{\alpha \beta }(x,x)`$
$$M=A_{\alpha \beta }(l,pl)e^{i(2lp)x}\mathrm{\Psi }_{\alpha \beta }(x,x)\frac{d^4l}{(2\pi )^4}2^4d^4x$$
(2.1)
$$\mathrm{\Psi }_{\alpha \beta }(x,x)=<\rho (p)|\overline{\mathrm{\Psi }}_\alpha (x)\mathrm{\Psi }_\beta (x)|0>,$$
(2.2)
$`p`$ and $`l`$ are the momenta of the meson and the quark respectively, see Fig.1. The momentum of the target quark (gluon) is $`p_2`$. The factor $`2^4`$ is present in the above equation since the separation between quark and antiquark is $`r=2x`$.
We introduce light cone variables
$$l_\pm =1/\sqrt{2}(l_0\pm l_3),d^4l=dl_+dl_{}d^2𝐥,\xi =2u1,l_+=up_+.$$
(2.3)
The two-dimensional transverse vectors are denoted by bold face. We denote the longitudinal momentum fraction for the quark as $`u`$, for the antiquark as $`(1u)=\overline{u}`$.
The hard scattering amplitude $`A_{\alpha \beta }`$ depends weakly on $`l_{}`$, neglecting this dependence we perform integration over $`l_{}`$ and $`x_+`$
$$M=2A_{\alpha \beta }(l,pl)e^{i(\mathrm{𝐥𝐫})}\frac{d^2𝐥d^2𝐫}{(2\pi )^2}𝑑u\mathrm{\Psi }_{\alpha \beta }(u,𝐫)$$
(2.4)
$$\mathrm{\Psi }_{\alpha \beta }(u,𝐫)=\frac{d(p_+x_{})}{(2\pi )}e^{i(p_+x_{})\xi }\mathrm{\Psi }_{\alpha \beta }(x,x).$$
(2.5)
The $`𝐱`$ values contributing essentially to the integral are small $`1/q`$, therefore $`x^2=2x_+x_{}𝐱^20`$.
We use the usual Fierz transformation
$$\delta _{\alpha \alpha ^{}}\delta _{\beta \beta ^{}}=\frac{1}{4}\mathrm{\Gamma }_{\alpha \beta }^t\mathrm{\Gamma }_{t\beta ^{}\alpha ^{}}=\frac{1}{4}\mathrm{\Gamma }_{\beta \alpha }^t\mathrm{\Gamma }_{t\alpha ^{}\beta ^{}}$$
(2.6)
$$\mathrm{\Gamma }^t=\{1,\gamma ^\mu ,\gamma ^\mu \gamma _5,\sigma ^{\mu \nu },i\gamma _5\},\mathrm{\Gamma }_t=\left(\mathrm{\Gamma }^t\right)^1,$$
to disentangle the spinor indices of $`A_{\alpha \beta }(l,pl)`$ and $`\mathrm{\Psi }_{\alpha \beta }`$.
We use
$$<\rho (p)|\overline{\mathrm{\Psi }}(x)\mathrm{\Gamma }^t\mathrm{\Psi }(x)|0>=\left[<0|\overline{\mathrm{\Psi }}(x)\mathrm{\Gamma }^t\mathrm{\Psi }(x)|\rho (p)>\right]^{}.$$
(2.7)
Expressions for the light-cone wave functions of vector mesons were derived in and and we shall adopt the notations used there. For instance (see eq.(2.8) of ),
$`<0|\overline{u}(x)\gamma _\mu u(x)|\rho (p,\lambda )>=f_\rho m_\rho [{\displaystyle \frac{e^{(\lambda )}x}{px}}p_\mu {\displaystyle \underset{0}{\overset{1}{}}}due^{i\xi (px)}\varphi _{||}(u)`$
$`+(e_\mu ^{(\lambda )}p_\mu {\displaystyle \frac{e^{(\lambda )}x}{px}}){\displaystyle \underset{0}{\overset{1}{}}}due^{i\xi (px)}g_{}^{(v)}(u)]`$ (2.8)
Here we have neglected the terms proportional to $`m_\rho ^3`$ which involve wave functions of twist-4. The polarization state of vector meson with definite helicity $`\lambda `$ is described by $`e^{(\lambda )}`$.
According to all wave functions like $`\varphi _{||}(u),g_{}^{(v)}(u)`$ which parametrize the matrix elements $`<0|\overline{u}(x)\mathrm{\Gamma }_tu(x)|\rho >`$ are normalized in the same way
$$\underset{0}{\overset{1}{}}\varphi (u)𝑑u=1.$$
(2.9)
## 3 Hard scattering amplitude.
The hard scattering amplitude in eq. (2.4) was calculated in . The result is given as integral over the transverse gluon momentum $`𝐤`$
$$A_{\alpha \beta }=is\frac{J_{\alpha \beta }^{\gamma q\overline{q}}(𝐤,𝐪)J_{qq}(𝐤,𝐪)}{𝐤^2(𝐤𝐪)^2}d^2𝐤.$$
(3.1)
This impact factor representation is well known from the pioneering works . $`s`$ is the total c.m.s energy squared of the photon-quark collision, $`t=𝐪^2`$. The impact–factors $`J_{\gamma V}`$ and $`J_{qq}`$ correspond to the upper and the lower blocks in Fig. 1. They are $`s`$–independent. For colorless exchange the impact–factors contain factors $`\delta _{ab}`$, where $`a`$ and $`b`$ are the color indices of the exchanged gluons.
Due to gauge invariance the impact–factors, which describe the coupling to colorless state vanish when the gluon momenta tend to zero:
$$J_{\gamma V}(𝐤,𝐪)0\text{ at }\{\begin{array}{cc}𝐤& \hfill 0,\\ (𝐩𝐤)& \hfill 0.\end{array}$$
(3.2)
This property garanties the infrared safety of the integral (3.1). The quark and gluon impact factors are
$$J_{qq}=\alpha _s\frac{\delta _{ab}}{N};J_{gg}=\alpha _s\delta _{ab}\frac{2N}{N^21}.$$
(3.3)
The helicity and color state of the quark or gluon target are conserved by these vertices. The relations (3.3) show that the cross section for photoproduction of vector mesons by gluons is about 5 times larger than by quarks:
$$d\sigma _{\gamma gVg}=\left(\frac{2N^2}{N^21}\right)^2d\sigma _{\gamma qVq}=\frac{81}{16}d\sigma _{\gamma qVq}.$$
(3.4)
The impact factor describing the upper part of the diagram Fig. 1 is
$$J_{\gamma q\overline{q}}=eQ_qg^2\frac{\delta _{ab}}{2N}(\left[mR\widehat{e}2u(\mathrm{𝐏𝐞})\widehat{P}\widehat{e}\right]\frac{\widehat{p}_2}{s})_{\alpha \beta }.$$
(3.5)
Here $`R`$ and the transverse vector $`P=(0,𝐏,0)`$ are:
$`𝐏`$ $`=`$ $`[{\displaystyle \frac{𝐪_1}{𝐪_1^2+m^2}}+{\displaystyle \frac{𝐤𝐪_1}{(𝐤𝐪_1)^2+m^2}}][𝐪_1𝐪_2];`$
$`R`$ $`=`$ $`[{\displaystyle \frac{1}{𝐪_1^2+m^2}}{\displaystyle \frac{1}{(𝐤𝐪_1)^2+m^2}}]+[𝐪_1𝐪_2].`$ (3.6)
According to eq. (2.5) we have to switch to the mixed representation, i.e. the momentum representation with respect to gluon $`t`$channel momenta and the coordinate representation with respect to transverse distance between quark and antiquark $`𝐫`$.
To perform the corresponding Fourier transform we have to express the quark and antiquark momenta $`𝐪_1,𝐪_2`$ through the momentum of quark $`𝐥`$ which is transverse with respect to the meson momentum $`p`$
$$𝐪_\mathrm{𝟏}=𝐥+𝐪u𝐪_\mathrm{𝟐}=𝐥+𝐪\overline{u}.$$
(3.7)
As result we obtain
$$𝐏(𝐫)=\frac{d^2𝐥}{(2\pi )^2}e^{i\mathrm{𝐥𝐫}}𝐏=\frac{d^2𝐥}{(2\pi )^2}\frac{𝐥}{𝐥^2+m^2}e^{i\mathrm{𝐥𝐫}}f^{dipole}.$$
(3.8)
$$R(𝐫)=\frac{1}{2\pi }K_0(rm)f^{dipole}.$$
(3.9)
The dipole amplitude has now appeared. It is given by the formulae
$$f^{dipole}=e^{i\mathrm{𝐪𝐫}u}\left(1e^{i\mathrm{𝐤𝐫}}\right)\left(1e^{i(𝐪𝐤)𝐫}\right)$$
(3.10)
In the massless limit
$$𝐏(𝐫)=\frac{im}{2\pi }K_1(rm)\frac{𝐫}{r}f^{dipole}|_{m0}\frac{i}{2\pi }\frac{𝐫}{r^2}f^{dipole}.$$
(3.11)
The trace calculations for those hard scattering amplitudes with Fierz structures ($`\gamma _\mu ,\gamma _\mu \gamma _5,\sigma _{\mu \nu }`$) which lead to the dominant ($`s`$) contribution is straightforward. Let us define
$$\widehat{Q}=\left(mR\widehat{e}2u(\mathrm{𝐏𝐞})\widehat{P}\widehat{e}\right)\frac{\widehat{p}_2}{s},$$
then the relevant traces which have to be calculated are
$`{\displaystyle \frac{1}{4}}Tr[\gamma _\mu \widehat{Q}]=(12u)(\mathrm{𝐏𝐞}){\displaystyle \frac{p_2^\mu }{s}}`$
$`{\displaystyle \frac{1}{4}}Tr[\gamma _\mu \gamma _5\widehat{Q}]={\displaystyle \frac{i}{s}}\epsilon _{\mu \nu \sigma \tau }P^\nu e^\sigma p_2^\tau `$
$`{\displaystyle \frac{1}{4}}Tr[\sigma _{\mu \nu }\widehat{Q}]={\displaystyle \frac{im}{s}}(p_{2\mu }e_\nu e_\mu p_{2\nu })R`$ (3.12)
## 4 Meson wave functions
Now we transform $`\rho `$ meson wave functions separately for longitudinal and transverse polarization into a form convenient for subsequent calculations.
a) $`\mathrm{\Gamma }_t=\gamma _\mu `$, longitudinal polarization
$$<0|\overline{u}(x)\gamma _\mu u(x)|\rho (p,\lambda =0)>=f_\rho p_\mu \underset{0}{\overset{1}{}}𝑑ue^{i\xi (p_+x_{})}\varphi _{||}(u)$$
(4.1)
The Fourier transform (FT) of this expression with respect to $`(p_+x_{})`$, gives
$$\frac{1}{2}p_\mu f_\rho \varphi _{||}(u)$$
(4.2)
b) $`\mathrm{\Gamma }_t=\gamma _\mu `$, transverse polarization
$`<0|\overline{u}(x)\gamma _\mu u(x)|\rho (p,\lambda =T)>=`$
$`=f_\rho m_\rho p_\mu {\displaystyle \frac{(𝐞^{(T)}𝐱)}{(p_+x_{})}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑ue^{i\xi (p_+x_{})}(\varphi _{||}(u)g_{}^{(v)}(u))`$ (4.3)
$$FT:if_\rho m_\rho p_\mu (𝐞^{(T)}𝐫)\underset{0}{\overset{1}{}}𝑑ue^{i\xi (p_+x_{})}\underset{0}{\overset{u}{}}𝑑v(\varphi _{||}(v)g_{}^{(v)}(v))$$
(4.4)
We consider only the contribution to the amplitude of the lowest Fock component of the meson wave function, i.e. the quark antiquark component, and we neglect quark masses. The twist-3 vector meson wave function $`g_{}^{(v)}(u)`$ is expressed through the twist-2 wave function $`\varphi _{||}(v)`$ with the help of a relation (WW) derived in which is similar to the one derived by Wandzura and Wilczek for $`g_2`$ structure function
$$g_{}^{(v)}(u)=g_{}^{(v)WW}(u)=\frac{1}{2}\left[\underset{0}{\overset{u}{}}\frac{dv}{\overline{v}}\varphi _{||}(v)+\underset{u}{\overset{1}{}}\frac{dv}{v}\varphi _{||}(v)\right].$$
(4.5)
$$FT:\frac{i}{2}f_\rho m_\rho (𝐞^{(T)}𝐫)\left(\frac{\overline{u}}{2}\underset{0}{\overset{u}{}}\frac{dv}{\overline{v}}\varphi _{||}(v)\frac{u}{2}\underset{u}{\overset{1}{}}\frac{dv}{v}\varphi _{||}(v)\right).$$
(4.6)
c) $`\mathrm{\Gamma }_t=\gamma _\mu \gamma _5`$, longitudinal polarization: no contribution in our approximation
d) $`\mathrm{\Gamma }_t=\gamma _\mu \gamma _5`$, transverse polarization:
This matrix element is given by eq. (2.9) of
$$<0|\overline{u}(x)\gamma _\mu \gamma _5u(x)|\rho (p,\lambda )>=\frac{1}{2}f_\rho m_\rho \epsilon _{\mu \nu \alpha \beta }e^{(\lambda )\nu }p^\alpha x^\beta \underset{0}{\overset{1}{}}𝑑ue^{i\xi (px)}g_{}^{(a)}(u).$$
(4.7)
The diffrence in sign between our definition for this matrix element and the corresponding one of is related to different sign conventions for $`\epsilon _{\mu \nu \alpha \beta }`$.
$`<0|\overline{u}(x)\gamma _\mu \gamma _5u(x)|\rho (p,\lambda =T)>=`$
$`={\displaystyle \frac{1}{2}}f_\rho m_\rho \epsilon _{\mu \nu \alpha \beta }e^{(T)\nu }p^\alpha x^\beta {\displaystyle \underset{0}{\overset{1}{}}}𝑑ue^{i\xi (p_+x_{})}g_{}^{(a)}(u)`$ (4.8)
$`FT:{\displaystyle \frac{1}{8}}f_\rho m_\rho \epsilon _{\mu \nu \alpha \beta }e^{(T)\nu }p^\alpha r^\beta g_{}^{(a)}(u)`$ (4.9)
$$g_{}^{(a)}(u)=g_{}^{(a)WW}(u)=2\overline{u}\underset{0}{\overset{u}{}}\frac{dv}{\overline{v}}\varphi _{||}(v)+2u\underset{u}{\overset{1}{}}\frac{dv}{v}\varphi _{||}(v)$$
(4.10)
e) $`\mathrm{\Gamma }_t=\sigma _{\mu \nu }`$, longitudinal polarization:
The parametrization of this matrix element is given by (2.16) of
$`<0|\overline{u}(x)\sigma _{\mu \nu }u(x)|\rho (p,\lambda )>=if_\rho ^T[(e_\mu ^{(\lambda )}p_\nu e_\nu ^{(\lambda )}p_\mu ){\displaystyle \underset{0}{\overset{1}{}}}due^{i\xi (px)}\varphi _{}(u)`$
$`+\left(p_\mu x_\nu p_\nu x_\mu \right){\displaystyle \frac{(e^{(\lambda )}x)}{(px)^2}}m_\rho ^2{\displaystyle \underset{0}{\overset{1}{}}}𝑑ue^{i\xi (px)}\left(h_{||}^{(t)}(u){\displaystyle \frac{1}{2}}\varphi _{}(u){\displaystyle \frac{1}{2}}h_3(u)\right)`$
$`+{\displaystyle \frac{1}{2}}(e_\mu ^{(\lambda )}x_\nu e_\nu ^{(\lambda )}x_\mu ){\displaystyle \frac{m_\rho ^2}{(px)}}{\displaystyle \underset{0}{\overset{1}{}}}due^{i\xi (px)}(h_3(u)\varphi _{}(u))].`$ (4.11)
$$FT:\frac{1}{2}f_\rho ^Tm_\rho \left(p_\mu r_\nu p_\nu r_\mu \right)\underset{0}{\overset{u}{}}𝑑v\left(h_{||}^{(t)}(v)\varphi _{}(v)\right).$$
(4.12)
$$h_{||}^{(t)}(u)=h_{||}^{(t)WW}(u)=\xi \left(\underset{0}{\overset{u}{}}\frac{dv}{\overline{v}}\varphi _{}(v)+\underset{u}{\overset{1}{}}\frac{dv}{v}\varphi _{}(v)\right).$$
(4.13)
With the help of this relation we find the following formulae
$$FT:\frac{1}{2}f_\rho ^Tm_\rho \left(p_\mu r_\nu p_\nu r_\mu \right)u\overline{u}\left(\underset{u}{\overset{1}{}}\frac{dv}{v}\varphi _{}(v)\underset{0}{\overset{u}{}}\frac{dv}{\overline{v}}\varphi _{}(v)\right)$$
(4.14)
f) $`\mathrm{\Gamma }_t=\sigma _{\mu \nu }`$, transverse polarization:
We have to consider two cases: the case which gives the chiral-odd contribution to the amplitude without spin-flip and the case with double spin-flip. In both cases the spins of the quarks in the loop are parallel. But in the double spin-flip case a spin-flip of orbital angular momentum by two units occures.
The wave function of a transversely polarized meson which contributes to the amplitude without spin-flip reads
$$\frac{i}{2}f_\rho ^T\left(e_\mu ^{(T)}p_\nu e_\nu ^{(T)}p_\mu \right)\varphi _{}(u).$$
(4.15)
For double spin-flip it is given by the second term in the square bracket of eq. (4)
$$\frac{i}{2}f_\rho ^Tm_\rho ^2\left(p_\mu r_\nu p_\nu r_\mu \right)(𝐞^{(T)}𝐫)\underset{0}{\overset{u}{}}𝑑v\underset{0}{\overset{v}{}}𝑑\eta \left(h_{||}^{(t)}(\eta )\frac{1}{2}\varphi _{}(\eta )\frac{1}{2}h_3(\eta )\right).$$
(4.16)
## 5 Helicity amplitudes
Now we are in a position to combine all results and calculate then the helicity amplitudes with the help of eq.(2.4).
We introduce short-hand notation
$$C=is\alpha _s^2\frac{N^21}{N^2}eQ_q$$
The contributions of the chiral-even configurations to various helicity amplitudes have the forms
$`M_{+0}^{even}=iC{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }(12u)}`$
$`{\displaystyle \frac{(\mathrm{𝐫𝐞}^{(+)})}{r^2}}f^{dipole}{\displaystyle \frac{f_\rho }{2}}\varphi _{||}(u)`$ (5.1)
$`M_{++}^{even}=C{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }f^{dipole}f_\rho m_\rho }`$
$`{\displaystyle \frac{u\overline{u}}{2}}\left({\displaystyle \underset{0}{\overset{u}{}}}{\displaystyle \frac{dv}{\overline{v}}}\varphi _{||}(v)+{\displaystyle \underset{u}{\overset{1}{}}}{\displaystyle \frac{dv}{v}}\varphi _{||}(v)\right)`$ (5.2)
$`M_+^{even}=C{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }f^{dipole}f_\rho m_\rho }`$
$`{\displaystyle \frac{(𝐫_x+i𝐫_y)^2}{2r^2}}\left(\overline{u}^2{\displaystyle \underset{0}{\overset{u}{}}}{\displaystyle \frac{dv}{\overline{v}}}\varphi _{||}(v)+u^2{\displaystyle \underset{u}{\overset{1}{}}}{\displaystyle \frac{dv}{v}}\varphi _{||}(v)\right)`$ (5.3)
Next we calculate the two-dimensional integrals over $`𝐫`$, $`𝐤`$ and the integral over $`u`$. The result for the single spin-flip amplitude can be written as
$`M_{+\mathrm{\hspace{0.17em}0}}^{even}=C{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}𝑑uf_\rho (12u)}`$
$`𝐞^{(+)}[{\displaystyle \frac{𝐪}{𝐪^2u}}+{\displaystyle \frac{𝐤𝐪u}{(𝐤𝐪u)^2}}(u\overline{u})]`$
$`=Cf_\rho \mathrm{\hspace{0.33em}2}\pi {\displaystyle \frac{(\mathrm{𝐪𝐞}^{(+)})}{|𝐪|^4}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{du}{u\overline{u}}}(12u)\varphi _{||}(u)\mathrm{ln}{\displaystyle \frac{\overline{u}}{u}}`$ (5.4)
This formula was derived first in .
For the asymptotic form of $`\varphi _{||}(u)=6u\overline{u}`$ we get
$$M_{+\mathrm{\hspace{0.17em}0}}^{even}=is\alpha _s^2\frac{N^21}{N^2}eQ_qf_\rho \mathrm{\hspace{0.33em}12}\pi \frac{(\mathrm{𝐪𝐞}^{(+)})}{|𝐪|^4}.$$
(5.5)
For the non spin-flip amplitude the integration over $`𝐫`$ and $`𝐤`$ leads to the expression
$$M_{++}^{even}=Cf_\rho m_\rho \frac{2\pi }{|𝐪|^4}\underset{0}{\overset{1}{}}\frac{du}{u\overline{u}}g_{}^{(v)WW}(u).$$
(5.6)
For the asymptotic form of the meson wave function (see eq. (4.5)) the integration over $`u`$ gives
$$M_{++}^{even}=is\alpha _s^2\frac{N^21}{N^2}eQ_qf_\rho m_\rho \frac{6\pi }{|𝐪|^4}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}1+2u_{min}\right),$$
(5.7)
where we restricted the integration region to $`[1u_{min},u_{min}]`$, $`u_{min}=m_\rho ^2/𝐪^2`$. The integration over this interval gives the contribution to $`M_{++}^{even}`$ of the region where the interquark distances $`𝐫^2\frac{1}{𝐪^2u\overline{u}}<\frac{1}{m_\rho ^2}`$ are small, and where we can use pQCD.
For the double spin-flip amplitude the two-dimensional integrations over $`𝐫`$ and $`𝐤`$ give
$`M_+^{even}=Cf_\rho m_\rho {\displaystyle \frac{4\pi }{|𝐪|^4}}`$
$`{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{du}{u^2\overline{u}^2}}\left({\displaystyle \frac{1}{4}}g_{}^{(a)WW}(u)u\overline{u}g_{}^{(v)WW}(u)\right)\left[{\displaystyle \frac{1}{2}}+(12u)\mathrm{ln}{\displaystyle \frac{\overline{u}}{u}}\right]`$ (5.8)
For the asymptotic form of wave function (see eqs. (4.5), (4.10)) this reads
$$M_+^{even}=is\alpha _s^2\frac{N^21}{N^2}eQ_qf_\rho m_\rho \frac{18\pi }{|𝐪|^4}.$$
(5.9)
The calculations of the chiral-odd contributions to helicity amplitudes are not so straightforward as in the chiral-even case.
Let us look again at eq. (2.4). This formula gives the amplitude as a convolution of the hard scattering amplitude $`A`$ and the vector meson light-cone wave function $`\mathrm{\Psi }`$. The hard scattering amplitude is given by eqs. (3.1), (3.5), (3.3). It is well known that at high energies such hard scattering amplitudes can be further factorized into a product of a dipole scattering amplitude and a perturbative photon wave function given by an $`e\overline{u}\gamma _\mu u`$ vertex. Therefore our scattering amplitude (2.4) can be rewritten as a product of the photon wave function, the dipole scattering amplitude and the vector meson wave function
$$M^{odd}=is\alpha _s^2\frac{N^21}{N^2}\frac{1}{4\pi }\frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}d^2𝐫𝑑u\mathrm{\Psi }^\gamma (𝐫,u)f^{dipole}\mathrm{\Psi }^\rho (𝐫,u).$$
(5.10)
This property is often called ’diffractive factorization’. It is valid even in the non-perturbative region where the photon splits into a dipole of large size. There are lots of phenomenological models in the literature which describe the photon wave function in this non-perturbative region.
Fortunately in our process high $`t`$ quark-dipole scattering chooses small interquark separations $`𝐫`$. Therefore we need only the photon wave functions on the light-cone which are very similar to the light-cone wave functions of a vector meson .
To simplify the calculation we use the following trick which permits us to proceed in a similar way as in the chiral-even case. First we calculate the chiral-odd amplitude for the perturbative photon wave function with finite quark mass. For that we have to insert into eq. (2.4) the result of the trace of $`\sigma _{\mu \nu }`$ structure with the hard scattering amplitude (see the last eq. (3)) and the chiral-odd vector meson wave function (eq.(4.15)). For the no spin-flip case the result has the form
$`M_{++}^m=C{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }f^{dipole}}`$
$`K_0(mr)f_\rho ^Tm(𝐞^{(+)}𝐞^{(+)})\varphi _{}(u),`$ (5.11)
whereas for the single spin-flip transition it reads
$`M_{+\mathrm{\hspace{0.17em}0}}^m=(i)C{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }f^{dipole}}`$
$`mK_0(mr)f_\rho ^Tm_\rho u\overline{u}(𝐞^{(+)}𝐫)\left({\displaystyle \underset{u}{\overset{1}{}}}{\displaystyle \frac{dv}{v}}\varphi _{}(v){\displaystyle \underset{0}{\overset{u}{}}}{\displaystyle \frac{dv}{\overline{v}}}\varphi _{}(v)\right).`$ (5.12)
In these two cases only the quark spin configuration in which the spins of quark and antiquark are parallel to the that of the incomming photon contributes.
A comparison of (5) and (5.10) leads to the identification of the relevant chiral-odd meson and photon wave functions
$$\mathrm{\Psi }_+^{\rho odd}(u)=\frac{2\pi }{\sqrt{2}}\frac{\delta ^{ij}}{N}f_\rho ^T\varphi _{}(u),$$
(5.13)
$$\mathrm{\Psi }_{+pert}^{\gamma odd}=eQ_q\sqrt{2}\delta ^{ij}\frac{m}{2\pi }K_0(mr),$$
(5.14)
with definite colours $`i`$ and $`j`$ of quark and antiquark respectively (note that the perturbative chiral-odd photon wave function is well known. see e.g. ). The non-perturbative light-cone wave function of the photon has to have a similar form as the wave function of vector mesons (5.13)
$$\mathrm{\Psi }_{+nonpert}^{\gamma odd}(u)=\frac{2\pi }{\sqrt{2}}eQ_q\frac{\delta ^{ij}}{N}f_\gamma \varphi _{}^\gamma (u).$$
(5.15)
Except the evident factor $`eQ_q`$ the only difference between the meson and the photon wave function are the different dimensional coupling constants, $`f_\gamma `$ and $`f_\rho ^T`$. The photon coupling constant $`f_\gamma `$ is a product of the quark condensate $`<\overline{q}q>`$ and its magnetic susceptibility ,
$$f_\gamma =<\overline{q}q>\chi 70MeV,$$
(5.16)
at the normalization point $`\mu =1GeV`$. Both parameters $`<\overline{q}q>`$ and $`\chi `$ have been tested in various QCD sum rule applications. Let us emphasize in view of our final results that the value of $`\chi `$ is large. The rough sum rule estimate which takes into account only the lowest lying hadronic state, namely the $`\rho `$ meson is $`\chi =2/m_\rho ^2=3.3\text{ GeV}^2`$. A more accurate analysis gives $`\chi =4.4\text{ GeV}^2\pm 0.4`$ (the value of $`f_\gamma `$ in eq.(5.16) corresponds roughly to this value for $`<\overline{q}q>=0.017`$ GeV<sup>3</sup>). Note that the quantity $`\chi `$ being a parameter describing the QCD vacuum plays an important role in the known sum rules for the electromagnetic form factor of baryons.
Now in order to obtain the amplitude of the interest we have only to substitute
$$\mathrm{\Psi }_{+pert}^{\gamma odd}\mathrm{\Psi }_{+nonpert}^{\gamma odd}(u)$$
(5.17)
in eqs. (5) and (5) and after taking into account the trace over colour indicies obtain
$`M_{++}^{odd}={\displaystyle \frac{C}{N}}{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}\frac{d^2𝐫du}{4\pi }f^{dipole}}`$
$`2\pi ^2f_\gamma \varphi _{}^\gamma (u)f_\rho ^T(𝐞^{(+)}𝐞^{(+)})\varphi _{}(u)`$
$`={\displaystyle \frac{C}{N}}f_\rho ^Tf_\gamma {\displaystyle \frac{4\pi ^3}{|𝐪|^4}}{\displaystyle \underset{0}{\overset{1}{}}}{\displaystyle \frac{du}{u^2\overline{u}^2}}\varphi _{}^\gamma (u)\varphi _{}(u).`$ (5.18)
For the asymptotic forms of the photon and the vector meson wave functions $`\varphi _{}^\gamma (u)=\varphi _{}(u)=6u\overline{u}`$ this gives
$$M_{++}^{odd}=is\alpha _s^2\frac{N^21}{N^3}eQ_q\frac{144\pi ^3}{|𝐪|^4}f_\rho ^Tf_\gamma .$$
(5.19)
Let us compare this result with the chiral-even contribution to the non-flip amplitude (5.7). The origin of the large relative coefficient $`(2\pi )^2`$ in the chiral-odd contribution can be related to the factors of $`2\pi `$ in Eqs. (5.13), (5.15) which in turn are related to QCD sum rules calculations for which the appearence of such factors is actually typical. When a quark loop is ”broken up” by inserting a quark condensate $`<\overline{q}q>`$ there is no integration over the full loop momentum range or $`\overline{q}q`$ phase space left as the virtuality is restricted by some hadronic scale. The numerical factor is just the (two dimensional) minimal phase space volume and the actual expansion paramenter is $`(2\pi )^2<\overline{q}q>`$.
Due to the above mentioned large numerical factor and large value of $`\chi `$ the chiral-odd contributions to helicity amplitudes are large. Note that the direct source of chiral symmetry breaking related to a nonzero quark mass $`m`$ leads to negligible contributions to the helicity amplitudes since light quarks have very small current masses and the $`\gamma q\overline{q}`$ chiral-odd vertex is proportional to this mass. If one would instead consider a model where the point-like form of the $`\gamma q\overline{q}`$ vertex is used and the quark masses are of order of constituent quark masses, $`m200\text{MeV}`$, then still the chiral-odd amplitudes would be much smaller since the value of constituent quark mass is much smaller than the parameter $`(2\pi )^2f_\gamma `$.
For the single spin flip case we obtain
$`M_{+\mathrm{\hspace{0.33em}0}}^{odd}={\displaystyle \frac{C}{N}}{\displaystyle \frac{i\sqrt{2}\pi ^2}{4\pi }}{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}d^2𝐫𝑑uf_\gamma \varphi _{}^\gamma (u)f_\rho ^Tm_\rho }`$
$`\left(𝐫_x+i𝐫_y\right)f^{dipole}u\overline{u}\left({\displaystyle \underset{u}{\overset{1}{}}}{\displaystyle \frac{dv}{v}}\varphi _{}(v){\displaystyle \underset{0}{\overset{u}{}}}{\displaystyle \frac{dv}{\overline{v}}}\varphi _{}(v)\right)`$
$`={\displaystyle \frac{C}{N}}{\displaystyle \frac{\sqrt{2}\pi (2\pi )^2}{|𝐪|^5}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑uf_\gamma \varphi _{}^\gamma (u)f_\rho ^Tm_\rho {\displaystyle \frac{(12u)}{u^3\overline{u}^3}}`$
$`\left({\displaystyle \underset{u}{\overset{1}{}}}{\displaystyle \frac{dv}{v}}\varphi _{}(v){\displaystyle \underset{0}{\overset{u}{}}}{\displaystyle \frac{dv}{\overline{v}}}\varphi _{}(v)\right).`$ (5.20)
Using the asymptotic forms for $`\varphi _{}^\gamma (u)`$ and $`\varphi _{}(u)`$ this expression takes the form
$$M_{+\mathrm{\hspace{0.33em}0}}^{odd}=\frac{C}{N}f_\gamma f_\rho ^Tm_\rho \frac{72\sqrt{2}\pi ^3}{|𝐪|^5}\underset{0}{\overset{1}{}}𝑑u\frac{(12u)^2}{u\overline{u}}.$$
(5.21)
Performing the remaining integral over $`u`$ we obtain
$$M_{+\mathrm{\hspace{0.33em}0}}^{odd}=is\alpha _s^2\frac{N^21}{N^3}eQ_qf_\gamma f_\rho ^Tm_\rho \frac{144\sqrt{2}\pi ^3}{|𝐪|^5}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}2(12u_{min})\right).$$
(5.22)
where we have again introduced suitable integration limit.
Finally we calculate the chiral-odd part of the double spin-flip amplitude. This case differs from the cases with no spin-flip and single spin-flip in one important aspect. In the last two cases only one spin configuration in the quark loop contributes, when the sum of quark helicities is equal to the helicity of incomming photon. For double spin-flip both chiral-odd spin configurations $`(\lambda _q=+1/2,\lambda _{\overline{q}}=+1/2)`$ (case (a)) and $`(\lambda _q=1/2,\lambda _{\overline{q}}=1/2)`$ (case (b)) contribute. For (a) the helicity of the initial photon $`\lambda _\gamma =+1`$ is carried by the helicity of quarks $`S_z=\lambda _q+\lambda _{\overline{q}}=+1`$. After interaction the dipole aquires $`L_z=2`$ which results in the meson helicity $`\lambda _\rho =S_z+L_z=1`$. In the case (b) the photon with $`\lambda _\gamma =+1`$ splits into $`q\overline{q}`$ state with $`(\lambda _q=1/2,\lambda _{\overline{q}}=1/2)`$ $`S_z=1`$ and $`L_z=+2`$. A flip by 2 units of $`L_z`$ happens and the dipole has $`S_z=1`$ and $`L_z=0`$, which gives again the meson helicity $`\lambda _\rho =S_z+L_z=1`$.
The contribution related to case (a) $`M_+^{odda}`$ can be calculated in a similar way. Here only the second term in eq. (4) contributes and after introducing the chiral-odd non-perturbative photon wave function (5.15) and using the asymptotic forms of the functions $`h_{||}^{(t)}(u)=3(2u1)^2,\varphi _{}(u)`$ and $`h_3(u)=1C_2^{1/2}(2u1)`$ ($`C_2^{1/2}`$ is the Gegenbauer polynom) we obtain
$`M_+^{odda}={\displaystyle \frac{C}{N}}{\displaystyle \frac{3\pi }{4}}f_\gamma f_\rho ^Tm_\rho ^2`$
$`{\displaystyle \frac{d^2𝐤}{𝐤^2(𝐤𝐪)^2}d^2𝐫𝑑u\varphi _{}^\gamma (u)u^2\overline{u}^2f^{dipole}(𝐞^{(+)}𝐫)(𝐞^{()}𝐫)}`$ (5.23)
After performing the relevant integrations the final result takes the form
$$M_+^{odda}=is\alpha _s^2\frac{N^21}{N^3}eQ_q\frac{144\pi ^3}{|𝐪|^6}f_\gamma f_\rho ^Tm_\rho ^2\left(2\mathrm{ln}\frac{1u_{min}}{u_{min}}3(12u_{min})\right).$$
(5.24)
For case (b) we need the photon wave function with $`S_z=1`$ and $`L_z=+2`$. For a vector meson the corresponding wave function is described by the second term in the square bracket of eq. (4). Unfortunately, we do not know a comprehensive QCD analysis of the photon wave function beyond twist-2. In this situation we make the assumption that the photon wave function with $`S_z=1`$, $`L_z=+2`$ differs from the corresponding meson wave function only by the replacement $`f_\rho ^TeQ_qf_\gamma `$. The
$$M_+^{odda}=M_+^{oddb}$$
(5.25)
and
$$M_+^{odd}=2M_+^{odda},$$
(5.26)
where $`M_+^{odda}`$ is given by eq.(5.24).
## 6 Discussion
We have explicitely derived helicity amplitudes of the process $`\gamma qVq`$. They are the sums of the chiral-even and the chiral-odd contributions
$$M_{\lambda _1\lambda _2}=M_{\lambda _1\lambda _2}^{even}+M_{\lambda _1\lambda _2}^{odd}.$$
(6.1)
The formulaes for the amplitudes are given by eqs. (5.5), (5.7), (5.9), (5.19), (5.22), (5.26), (5.24), they are the main results of the present work.
At asymptotically high momentum transfer the dominant helicity amplitude is $`M_{+\mathrm{\hspace{0.33em}0}}`$. Its chiral-even part $`M_{+\mathrm{\hspace{0.33em}0}}^{even}`$ has the minimal, $`1/|𝐪|^3`$, suppression. To discuss the onset of this asymptotic regime we consider the ratios $`M_{\lambda _1\lambda _2}/M_{+\mathrm{\hspace{0.33em}0}}^{even}`$
$$\frac{M_{+\mathrm{\hspace{0.33em}0}}}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=1\frac{24\pi ^2f_\gamma f_\rho ^Tm_\rho }{Nf_\rho 𝐪^2}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}2(12u_{min})\right)$$
(6.2)
$$\frac{M_{++}}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=\frac{m_\rho }{|𝐪|\sqrt{2}}\left(\mathrm{ln}\frac{1u_{min}}{u_{min}}1+2u_{min}\right)+\frac{24\pi ^2f_\gamma f_\rho ^T}{Nf_\rho |𝐪|\sqrt{2}}$$
(6.3)
$$\frac{M_+}{M_{+\mathrm{\hspace{0.33em}0}}^{even}}=\frac{3m_\rho }{\sqrt{2}|𝐪|}\frac{48\pi ^2f_\gamma f_\rho ^Tm_\rho ^2}{Nf_\rho |𝐪|^3\sqrt{2}}\left(2\mathrm{ln}\frac{1u_{min}}{u_{min}}3(12u_{min})\right)$$
(6.4)
The chiral-odd parts of the helicity amplitudes are proportional to dimensionfull coupling constants $`f_\gamma ,f_\rho ^T`$. The scale dependence for three active flavours is, see ,
$$\frac{f_\rho ^T(Q^2)}{f_\rho ^T(\mu ^2)}=\frac{f_\gamma (Q^2)}{f_\gamma (\mu ^2)}=L^{4/27},L=\frac{\alpha _s(Q^2)}{\alpha _s(\mu ^2)}.$$
The factorization scale in the case of our process is $`Q^2t`$.
The counting of relative factors of $`1/|𝐪|`$ for the chiral-even and the chiral-odd contributions to the helicity amplitudes in eqs. (6.2), (6.3), (6.4) can be easily understood. These contributions can be represented as convolutions of the corresponding photon and vector meson wave functions with the dipole scattering amplitude (3.10). The dipole amplitude does not depend on the helicity state of the quarks. (Helicity states do not change during the dipole interaction.) The projection of $`q\overline{q}`$ angular momentum on the axies of dipole motion $`L_z`$ can change due to the interaction, $`L_zL_z^{}`$. The important observation is that in the case of high $`t`$ scattering such a change ($`L_zL_z^{}`$) is not suppressed by a factor $`1/|𝐪|`$. Therefore the power counting of various contributions to helicity amplitudes is determined entirely by the behaviour at small $`𝐫`$ of the corresponding photon and vector meson wave functions.
The wave function describing the chiral-even $`q\overline{q}`$ fluctuation of the photon is given by first order perturbation theory. In the impact parameter representation it is simply a Fourier transform of the quark propagator, see eqs. (3.8) and (3.11). Note that it is power-like divergent at small interquark separations, $`𝐫/r^2`$. The total helicity of the quarks in the chiral even configuration is zero, $`S_z=0`$, therefore the helicity of the photon is carried by the orbital angular momentum of the quarks, $`\lambda _\gamma =S_z+L_z=L_z=\pm 1`$.
In contrast in the non-perturbative chiral-odd $`q\overline{q}`$ configuration the helicity of the photon is carried by the total helicity of quarks $`\lambda _\gamma =S_z,L_z=0`$ (the other possibility $`\lambda _\gamma =+1`$, $`S_z=1`$, $`L_z=+2`$ is relevant only for $`M_+^{oddb}`$ (5.25)). The small $`𝐫`$ asymptotics of this chiral-odd wave function is constant, it is described by the dimensionful non-perturbative parameter $`f_\gamma `$, see eq. (5.15). This asymptotics should be compared with the asymptotics of the chiral-even photon wave function $`1/r|𝐪|`$.
Various vector meson wave functions have different small $`𝐫`$ behaviour. The wave functions of (scaling) twist two describe $`q\overline{q}`$ pair with $`L_z=0`$, they are constants at small $`𝐫`$, see eqs. (4.2) and (4.15) for chiral-even and chiral-odd ones respectively. The configurations with $`L_z0`$ are described by the wave functions of higher twists. In the case $`L_z=\pm 1`$ they behave as $`𝐫1/|𝐪|`$, see eqs. (4.4), (4.9) for chiral-even and (4.12) for chiral-odd wave functions. The state with $`L_z=\pm 2`$ is given by the the wave function $`r^21/q^2`$, see eq. (4.16).
Let us discuss the helicity non flip amplitude (6.3). The chiral-even part of this amplitude is given by the configuration with $`L_z=+1`$, therefore the product of the corresponding photon and meson wave functions (4.4, 4.9) is $`\frac{1}{r}m_\rho f_\rho r=const`$. The chiral-odd part in this amplitude is given by the configuration with $`L_z=0`$ and in this case the product of chiral-odd photon wave function (5.15) and vector meson wave function (4.15) $`f_\gamma f_\rho ^T=const`$. This is the explanation why both parts of (6.3) have similar $`1/|𝐪|`$ suppression.
In the chiral-odd case the leading twist $`L_z=0`$ wave function of the meson enters the amplitude, in the chiral-even part it is the wave function of higher twist $`L_z=\pm 1`$. This difference is compensated by the difference in the small $`𝐫`$ behaviour of the point-like chiral-even and the non-perturbative chiral-odd wave functions of the real photon. The same arguments apply to the spin-flip amplitudes (6.2) and (6.4).
We find that the chiral-odd contributions to the amplitudes are accompanied with astonishingly large numerical coefficients. In the case of $`M_{++}`$ the chiral-even and the chiral-odd parts of (6.1) add with the same signs. In contrast, the chiral-even and the chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ or $`M_+`$ enters with opposite sign which leads to an effective reduction of these amplitudes for intermediate $`|t|`$. Even in the case of $`M_{+\mathrm{\hspace{0.33em}0}}`$ (or $`M_+`$) where the chiral-odd part is formally $`1/q^2`$ suppressed in comparison with the corresponding chiral-even part, the chiral-odd part cannot be neglected up to very large momentum transfers. For instance in the $`|t|`$ interval $`3\text{ GeV}^2÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$ the ratio (6.2) takes values from $`0.56`$ to $`0.62`$, and the ratio (6.4) varies from $`0.05`$ to $`0.08`$. Though the chiral-even and the chiral-odd contributions to $`M_{++}`$ are of the same order with respect to $`1/|𝐪|`$ counting the chiral-odd one can be dominant up to very large $`|t|`$. According to (6.3) for the $`t`$ range $`3÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$ the chiral-even part constitutes only $`10÷\mathrm{\hspace{0.33em}20}\%`$ of the $`M_{++}`$ amplitude and for $`|t|100\text{GeV}^2`$ $`M_{++}^{even}0.72M_{++}^{odd}`$.
Let us discuss next the relative magnitude of the various helicity amplitudes. On the one hand at asymptotically large $`|t|`$ $`M_{+\mathrm{\hspace{0.33em}0}}M_{+\mathrm{\hspace{0.33em}0}}^{even}`$ will dominate. On the other hand in the intermediately large $`t`$ region there is a large compensation between chiral-even and chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$, the similar compensation takes place for $`M_+`$ amplitude. Therefore in this region the non spin-flip $`M_{++}`$ amplitude dominates strongly. According to eqs. (6.2), (6.3), $`M_{+\mathrm{\hspace{0.33em}0}}`$ will exceed $`M_{++}`$ only at $`|t|>40(\text{ GeV})^2`$. For the $`|t|`$ interval $`3\text{ GeV}^2÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$
$$\frac{M_{+\mathrm{\hspace{0.33em}0}}}{M_{++}}0.25÷\mathrm{\hspace{0.33em}0.35}$$
(6.5)
$$\frac{M_+}{M_{++}}\mathrm{\hspace{0.17em}0.02}÷\mathrm{\hspace{0.33em}\hspace{0.17em}0.04}$$
(6.6)
Both chiral-even and chiral–odd parts of the amplitudes were calculated expecting $`|t|`$ to be large, or in the leading order of $`1/|𝐪|`$ expansion. As usual in the QCD approach to any exclusive reaction, the question about the region of applicability of these results is open untill the power corrections have not studied. In our case the situation is more difficult because the factorization of the amplitude into hard and soft parts is violated for the chiral-even part of $`M_{++}`$ and the chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ and $`M_+`$. In these cases the corresponding integrals over the quark longitudinal momentum $`u`$ contain the end point logarithmic singularities. This means that in these cases the contribution of the soft regions, where the transverse separation between quark and antiquark is large, is not power suppressed. The account of higher orders would result in an appearence of Sudakov like form factor which describes the suppression due to the change of the colour direction of motion without radiation of gluons. In the hard region, where $`q\overline{q}`$ pair scatters as a colourless dipole of the small size, this Sudakov suppression doesn’t work. It will come into a game for the scattering of dipoles of large sizes and therefore will lead to an effective suppression of the soft region.
At present we simply restrict the corresponding $`u`$ integrals to the interval $`[1u_{min},u_{min}],u_{min}=m_\rho ^2/𝐪^2`$, which corresponds to the contribution of the hard region only. It is not known at the moment how to calculate in a model independent way the soft contributions to the chiral-even part of $`M_{++}`$ and the soft contributions to the chiral-odd parts of $`M_{+\mathrm{\hspace{0.33em}0}}`$ and $`M_+`$. The good news is, however, that the chiral-odd part $`M_{++}^{odd}`$ of the dominant in the intermediatly large $`|t|`$ region helicity amplitude $`M_{++}`$ is free from this end-point singularity. This chiral-odd part is numerically considerably larger than the hard contribution to its chiral-even counterpart $`M_{++}^{even}`$. Therefore we can expect that the relative uncertainty related with the uncalculated soft contributions to $`M_{++}^{even}`$ is small for $`|t|3\text{ GeV}^2÷\mathrm{\hspace{0.33em}8}\text{ GeV}^2`$. For $`M_{+\mathrm{\hspace{0.33em}0}}`$ and $`M_+`$ the uncertainties related with the corresponding nonfactorizable parts of $`M_{+\mathrm{\hspace{0.33em}0}}^{odd}`$ and $`M_+^{odd}`$ can be larger. Despite that we believe that in the intermediately large $`t`$ region our predictions (6.2) and (6.4) for the ratios of the helicity amplitudes will be true on a qualitative level.
A recently reported ZEUS analysis of the angular distribution of $`\rho ^0`$ meson decay products from the process $`\gamma p\rho ^0p`$ at $`t1÷2\text{ GeV}^2`$ have actually demonstrated the dominance of the non spin-flip amplitude in this $`t`$ range. Unfortunately, the precision of these data is low and the only conclusions which can be drawn are that the spin-flip amplitudes are small, the relative sign between $`M_{++}`$ and $`M_{+\mathrm{\hspace{0.33em}0}}`$ tends to be positive and the sign between $`M_{++}`$ $`M_+`$ tends to be negative. Note that this is in agreement with our predictions (6.2), (6.3), (6.4).
Acknowledgments
We are grateful to Lonya Frankfurt and Mark Strikman for extremely helpfull discussions during the whole duration of our studies.
We also thank Vladimir Braun and Oleg Teryaev for many clarifying discussions.
L.Sz. acknowledges support by DFG. He acknowledges also support by Saxonian Ministry SMWK during his visit to Leipzig University. D. I. acknowledges support from BMBF and by a grant from Sankt-Petersburg Center of Fundamental Natural Science and by the grant RFBR 99-02-17211. |
warning/0001/math0001082.html | ar5iv | text | # Une identité remarquable en théorie des partitions
## 1 Notations
Nous démontrons dans cet article une conjecture présentée dans un précédent travail . Il s’agit d’une identité qui se rencontre dans l’étude des polynômes “symétriques décalés” , où elle permet le développement explicite de certains “polynômes de Jack décalés”, notamment ceux assocés aux partitions lignes et colonnes.
Cette identité se formule de manière extrêmement simple dans le cadre de la théorie classique des partitions. Cependant il nous a semblé que sa preuve ne s’obtient commodément qu’en utilisant la structure (élémentaire) de $`\lambda `$-anneau de l’anneau des polynômes.
Une partition $`\lambda `$ est une suite décroissante finie d’entiers positifs. On dit que le nombre $`n`$ d’entiers non nuls est la longueur de $`\lambda `$. On note $`\lambda =(\lambda _1,\mathrm{},\lambda _n)`$ et $`n=l(\lambda )`$. On dit que $`\left|\lambda \right|=\underset{i=1}{\overset{n}{}}\lambda _i`$ est le poids de $`\lambda `$, et pour tout entier $`i1`$ que $`m_i(\lambda )=\text{card}\{j:\lambda _j=i\}`$ est la multiplicité de $`i`$ dans $`\lambda `$. On identifie $`\lambda `$ à son diagramme de Ferrers $`\{(i,j):1il(\lambda ),1j\lambda _i\}`$. On pose
$$z_\lambda =\underset{i1}{}i^{m_i(\lambda )}m_i(\lambda )!.$$
La généralisation suivante du coefficient binomial classique a été introduite dans . Soient $`\lambda `$ une partition et $`r`$ un entier $`1`$. On note $`\genfrac{}{}{0.0pt}{}{\lambda }{r}`$ le nombre de façons dont on peut choisir r points dans le diagramme de $`\lambda `$ de telle sorte que au moins un point soit choisi sur chaque ligne de $`\lambda `$.
Les coefficients binomiaux généralisés $`\genfrac{}{}{0.0pt}{}{\lambda }{r}`$ possèdent la fonction génératrice suivante
$$\underset{r1}{}\genfrac{}{}{0.0pt}{}{\lambda }{r}q^r=\underset{i=1}{\overset{l(\lambda )}{}}\left((1+q)^{\lambda _i}1\right)=\underset{i1}{}\left((1+q)^i1\right)^{m_i(\lambda )}.$$
Soient $`z`$ une indéterminée et $`n`$ un entier $`1`$. On note désormais
$$(z)_n=z(z+1)\mathrm{}(z+n1),[z]_n=z(z1)\mathrm{}(zn+1)$$
les factorielles “ascendante” et “descendante” classiques. On pose
$$\left(\genfrac{}{}{0pt}{}{z}{n}\right)=\frac{[z]_n}{n!}.$$
Soit $`X=\{X_1,X_2,X_3,\mathrm{}\}`$ une famille (infinie) d’indéterminées indépendantes. Pour tous entiers $`j,k0`$ on pose
$$P_{jk}(X)=\underset{|\mu |=j}{}\frac{{\displaystyle \genfrac{}{}{0.0pt}{}{\mu }{k}}}{z_\mu }\underset{i1}{}X_{i}^{}{}_{}{}^{m_i(\mu )}.$$
(1)
Comme on a $`\genfrac{}{}{0.0pt}{}{\mu }{k}=0`$ si $`k<l(\mu )`$, la sommation est limitée aux partitions $`\mu `$ telles que $`l(\mu )k`$. Il en résulte que $`P_{jk}(X)`$ est un polynôme de degré $`k`$. Comme on a $`\genfrac{}{}{0.0pt}{}{\mu }{k}=0`$ si $`k>\left|\mu \right|`$, on a $`P_{jk}(X)=0`$ pour tout $`k>j`$. On pose par convention $`P_{00}(X)=1`$.
On a par exemple facilement
$$P_{j1}(X)=X_j,$$
$$P_{j2}(X)=\frac{1}{2}(j1)X_j+\frac{1}{2}\underset{\begin{array}{c}j_1+j_2=j\hfill \\ j_1,j_21\hfill \end{array}}{}X_{j_1}X_{j_2}.$$
## 2 Notre résultat
Le but de cet article est de démontrer la conjecture suivante, que l’un de nous a formulée dans un précédent travail ( , Conjecture 2). Cette conjecture explicite un développement en série formelle.
###### Théorème 1.
Soient $`z,u`$ et $`X=\{X_1,X_2,X_3,\mathrm{}\}`$ des indéterminées indépendantes. Pour tous entiers $`n,r1`$ on a
$$\begin{array}{c}\underset{\left|\mu \right|=n}{}(1)^{rl(\mu )}\frac{{\displaystyle \genfrac{}{}{0.0pt}{}{\mu }{r}}}{z_\mu }\underset{i1}{}\left(z+\underset{k1}{}u^k\frac{(i)_k}{k!}X_k\right)^{m_i(\mu )}=\hfill \\ \hfill \underset{j0}{}u^j\left(\genfrac{}{}{0pt}{}{n+j1}{nr}\right)\left(\underset{k=0}{\overset{min(r,j)}{}}\left(\genfrac{}{}{0pt}{}{zj}{rk}\right)P_{jk}(X)\right).\end{array}$$
Cette conjecture est triviale pour $`r>n`$ car on a alors $`\genfrac{}{}{0.0pt}{}{\mu }{r}=0`$. Pour $`r=n`$ on obtient le résultat suivant.
###### Théorème 2.
Soient $`z,u`$ et $`X=\{X_1,X_2,X_3,\mathrm{}\}`$ des indéterminées indépendantes. Pour tout entier $`n1`$ on a
$$\begin{array}{c}\underset{\left|\mu \right|=n}{}\frac{(1)^{nl(\mu )}}{z_\mu }\underset{i1}{}\left(z+\underset{k1}{}u^k\frac{(i)_k}{k!}X_k\right)^{m_i(\mu )}=\hfill \\ \hfill \underset{j0}{}u^j\left(\underset{k=0}{\overset{min(n,j)}{}}\left(\genfrac{}{}{0pt}{}{zj}{nk}\right)P_{jk}(X)\right).\end{array}$$
Le Théorème 2 avait été auparavant conjecturé dans (Conjecture 4, page 465). Pour $`X=0`$ le Théorème 1 redonne le Théorème 1’ de (page 462).
## 3 Fonctions symétriques
Nous donnons d’abord ici les notations dont nous aurons besoin à propos de l’algèbre $`\mathrm{𝐒𝐲𝐦}`$ des fonctions symétriques, considérée d’un point de vue formel.
Soit $`A=\{a_1,a_2,a_3,\mathrm{}\}`$ un ensemble de variables, qui peut être infini (nous dirons que $`A`$ est un alphabet). On introduit les fonctions génératrices
$$\lambda _t(A)=\underset{aA}{}(1+ta),\sigma _t(A)=\underset{aA}{}\frac{1}{1ta},\mathrm{\Psi }_t(a)=\underset{aA}{}\frac{a}{1ta}$$
dont le développement définit les fonctions symétriques élémentaires $`\mathrm{\Lambda }^i(A)`$, les fonctions complètes $`S^i(A)`$ et les sommes de puissances $`\psi ^i(A)`$ :
$$\lambda _t(A)=\underset{i0}{}t^i\mathrm{\Lambda }^i(A),\sigma _t(A)=\underset{i0}{}t^iS^i(A),\mathrm{\Psi }_t(A)=\underset{i1}{}t^{i1}\psi ^i(A).$$
Lorsque l’alphabet $`A`$ est infini, chacun de ces trois ensembles de fonctions forme une base algébrique de $`\mathrm{𝐒𝐲𝐦}[A]`$, l’algèbre des fonctions symétriques sur $`A`$ (c’est-à-dire que ses éléments sont algébriquement indépendants).
On peut donc définir l’algèbre $`\mathrm{𝐒𝐲𝐦}`$ des fonctions symétriques, sans référence à l’alphabet $`A`$, comme l’algèbre sur $`𝐐`$ engendrée par les fonctions $`\mathrm{\Lambda }^i`$, $`S^i`$ ou $`\psi ^i`$.
Pour toute partition $`\mu =(\mu _i,1il(\mu ))=(i^{m_i(\mu )},i1)`$, on définit les fonctions $`\mathrm{\Lambda }^\mu `$, $`S^\mu `$ ou $`\psi ^\mu `$ en posant
$$f^\mu =\underset{i=1}{\overset{l(\mu )}{}}f^{\mu _i}=\underset{k1}{}(f^k)^{m_k(\mu )},$$
$`f^i`$ désigne respectivement $`\mathrm{\Lambda }^i`$, $`S^i`$ ou $`\psi ^i`$. Les fonctions $`\mathrm{\Lambda }^\mu `$, $`S^\mu `$, $`\psi ^\mu `$ forment une base linéaire de l’algèbre $`\mathrm{𝐒𝐲𝐦}`$.
On a la formule de Cauchy
$$\mathrm{\Lambda }^i=\underset{\left|\mu \right|=i}{}(1)^{il(\mu )}\frac{\psi ^\mu }{z_\mu }$$
ou encore
$$S^i=\underset{\left|\mu \right|=i}{}\frac{\psi ^\mu }{z_\mu }.$$
Pour toute partition $`\mu `$, on peut définir les fonctions symétriques monomiales $`\psi _\mu `$ et les fonctions de Schur $`S_\mu `$, qui forment également une base linéaire de l’algèbre $`\mathrm{𝐒𝐲𝐦}`$.
Les bases $`\mathrm{\Lambda }^\mu `$, $`S^\mu `$, $`\psi ^\mu `$, $`\psi _\mu `$ ou $`S_\mu `$ sont notées respectivement $`e_\mu `$, $`h_\mu `$, $`p_\mu `$, $`m_\mu `$ ou $`s_\mu `$ dans la littérature, notamment dans . Les notations utilisées ici sont celles de , qui sont plus adaptées aux $`\lambda `$-anneaux.
## 4 L’anneau des polynômes comme $`\lambda `$-anneau
Nous allons démontrer le Théorème 1 en utilisant le fait que l’anneau des polynômes possède une structure de $`\lambda `$-anneau.
Un $`\lambda `$-anneau est un anneau commutatif avec unité muni d’opérateurs qui vérifient certains axiomes. Nous renvoyons le lecteur à pour la théorie générale, et au chapitre 2 de pour leur application à l’analyse multivariée.
Nous n’utiliserons cette théorie que dans le cadre élémentaire suivant. Soit $`A=\{a_1,a_2,a_3,\mathrm{}\}`$ un alphabet quelconque. On considère l’anneau $`𝐑[A]`$ des polynômes en $`A`$ à coefficients réels. La structure de $`\lambda `$-anneau de $`𝐑[A]`$ consiste à définir une action de $`\mathrm{𝐒𝐲𝐦}`$ sur $`𝐑[A]`$.
### 4.1 Action de $`\mathrm{𝐒𝐲𝐦}`$
Les fonctions $`\psi ^i`$ formant un système de générateurs algébriques de $`\mathrm{𝐒𝐲𝐦}`$, écrivant tout polynôme sous la forme $`_{c,u}cu`$, avec $`c`$ constante réelle et $`u`$ un monôme en $`\{a_1,a_2,a_3\mathrm{}\}`$, on définit une action de $`\mathrm{𝐒𝐲𝐦}`$ sur $`𝐑[A]`$, notée $`[.]`$, en posant
$$\psi ^i[\underset{c,u}{}cu]=\underset{c,u}{}cu^i.$$
Pour tous polynômes $`P,Q𝐑[A]`$ on en déduit immédiatement $`\psi ^i[PQ]=\psi ^i[P]\psi ^i[Q]`$ et $`\psi ^\mu [PQ]=\psi ^\mu [P]\psi ^\mu [Q]`$.
L’action ainsi définie s’étend à tout élément de $`\mathrm{𝐒𝐲𝐦}`$. Ainsi on a
$$\lambda _t[\underset{c,u}{}cu]=\underset{c,u}{}(1+tu)^c,\sigma _t[\underset{c,u}{}cu]=\underset{c,u}{}(1tu)^c.$$
On en déduit $`S^i[P]=(1)^i\mathrm{\Lambda }^i[P]`$.
On notera le comportement différent des constantes $`c𝐑`$ et des monômes $`u`$ :
$$\begin{array}{c}\hfill \psi _i[c]=c,S^i[c]=\frac{(c)_i}{i!},\mathrm{\Lambda }^i[c]=\frac{[c]_i}{i!}\\ \hfill \psi _i[u]=u^i=S^i[u],\mathrm{\Lambda }^i[u]=0,i>1,\mathrm{\Lambda }^1[u]=u.\end{array}$$
(2)
Il est plus correct de caractériser les ”monômes” $`u`$ comme éléments de rang 1 (i.e. les $`u0,1`$ tels que $`\mathrm{\Lambda }^i[u]=0i>1`$), et les ”constantes” $`c𝐑`$ comme les éléments invariants par les $`\psi _i`$ (on dira aussi élément de type binomial).
Lorsqu’on utilise la théorie des $`\lambda `$-anneaux pour démontrer une identité algébrique, il est donc toujours nécessaire de préciser le statut de chaque élément. En particulier nous aurons à employer des indéterminées de rang 1, et d’autres de type binomial.
### 4.2 Extension aux séries formelles
On remarquera que si $`a_1,a_2,\mathrm{},a_n`$ sont des éléments de rang un, alors
$$\psi ^i[a_1+a_2+\mathrm{}+a_n]=a_1^i+a_2^i+\mathrm{}+a_n^i$$
est la valeur de la $`i`$-ème somme de puissance $`\psi ^i(a_1,a_2,\mathrm{},a_n)`$.
Dans la suite pour tout alphabet $`A=\{a_1,a_2,a_3,\mathrm{}\}`$, on notera $`A^{\mathrm{}}=_ia_i`$ la somme de ses éléments. Lorsque A est formé d’éléments de rang 1, on a ainsi pour toute fonction symétrique $`f`$,
$$f[A^{\mathrm{}}]=f(A).$$
(3)
En particulier si $`q`$ est de rang 1, on a
$$\psi ^i(1,q,q^2,q^3,\mathrm{},q^{n1})=\psi ^i[\underset{k=0}{\overset{n1}{}}q^k].$$
Il est naturel de vouloir écrire
$$\underset{k=0}{\overset{n1}{}}q^k=\frac{1q^n}{1q},$$
et d’étendre ainsi l’action de $`\mathrm{𝐒𝐲𝐦}`$ aux fonctions rationnelles. Il est également naturel de considérer un alphabet infini $`(1,q,q^2,q^3,\mathrm{})`$, de vouloir sommer la série
$$\underset{k0}{}q^k=\frac{1}{1q},$$
et d’étendre ainsi l’action de $`\mathrm{𝐒𝐲𝐦}`$ aux séries formelles à coefficients réels.
Pour cela on pose
$$\psi _i\left(\frac{cu}{dv}\right)=\frac{cu^i}{dv^i},$$
avec $`c,d`$ constantes réelles et $`u,v`$ des monômes en $`(a_1,a_2,a_3\mathrm{})`$.
L’action ainsi définie s’étend à tout élément de $`\mathrm{𝐒𝐲𝐦}`$. On munit ainsi l’anneau des séries formelles à coefficients réels d’une structure de $`\lambda `$-anneau. On a par exemple
$$\lambda _t[\frac{1}{1q}]=\underset{i1}{}(1+tq^i)=(t;q)_{\mathrm{}}$$
$$\sigma _t[\frac{1}{1q}]=\underset{i1}{}\frac{1}{1tq^i}=\frac{1}{(t;q)_{\mathrm{}}},$$
ce qui fait apparaître des quantités bien connues en $`q`$-calcul.
### 4.3 Formulaire
Les relations fondamentales suivantes sont des conséquences directes des relations (2). Certaines nous seront nécessaires. Pour tous $`P,Q`$ on a d’abord
$$\begin{array}{cc}\hfill S^i[P+Q]& =\underset{j=0}{\overset{i}{}}S^{ij}[P]S^j[Q]\hfill \\ \hfill \mathrm{\Lambda }^i[P+Q]& =\underset{j=0}{\overset{i}{}}\mathrm{\Lambda }^{ij}[P]\mathrm{\Lambda }^j[Q],\hfill \end{array}$$
(4)
ou de manière équivalente :
$$\begin{array}{c}\hfill \sigma _t[P+Q]=\sigma _t[P]\sigma _t[Q]\\ \hfill \lambda _t[P+Q]=\lambda _t[P]\lambda _t[Q].\end{array}$$
(5)
Pour tous $`P,Q`$ on a d’autre part
$$\begin{array}{cc}\hfill S^i[PQ]& =\underset{\left|\mu \right|=i}{}\frac{1}{z_\mu }\psi ^\mu [P]\psi ^\mu [Q]\hfill \\ & =\underset{\left|\mu \right|=i}{}\psi _\mu [P]S^\mu [Q]\hfill \\ & =\underset{\left|\mu \right|=i}{}S_\mu [P]S_\mu [Q],\hfill \end{array}$$
(6)
ou de manière équivalente :
$$\begin{array}{cc}\hfill \mathrm{\Lambda }^i[PQ]& =\underset{\left|\mu \right|=i}{}\frac{(1)^{il(\mu )}}{z_\mu }\psi ^\mu [P]\psi ^\mu [Q]\hfill \\ & =\underset{\left|\mu \right|=i}{}\psi _\mu [P]\mathrm{\Lambda }^\mu [Q]\hfill \\ & =\underset{\left|\mu \right|=i}{}S_\mu [P]S_\mu ^{}[Q],\hfill \end{array}$$
(7)
$`\mu ^{}`$ désigne la partition transposée de $`\mu `$.
Si $`P`$ est de rang 1 et $`Q`$ arbitraire, on a
$$\mathrm{\Lambda }^i[PQ]=P^i\mathrm{\Lambda }^i[Q].$$
Ainsi lorsque $`P`$ et $`Q`$ sont de rang 1, on a
$$\lambda _t[PQ]=1+tPQ.$$
(8)
## 5 Démonstration du Théorème 1
### 5.1 Préliminaires
###### Lemme 1.
Soit $`q^{}`$ un élément de rang 1. Si on pose $`q=q^{}1`$, on a
$$\psi ^\mu [q]=\underset{k1}{}\genfrac{}{}{0.0pt}{}{\mu }{k}q^k.$$
###### Preuve.
On a
$$\psi ^i[q]=\psi ^i[q^{}1]=(q^{})^i1=(1+q)^i1.$$
On applique la fonction génératrice des entiers $`\genfrac{}{}{0.0pt}{}{\mu }{k}`$. ∎
A l’aide des relations (2) et (4) on obtient facilement
$$\mathrm{\Lambda }^i[q]=(1)^{i1}q,S^i[q]=(1+q)^{i1}q,i1.$$
(9)
On en déduit
$$\mathrm{\Lambda }^\mu [q]=(1)^{\left|\mu \right|l(\mu )}q^{l(\mu )},S^\mu [q]=(1+q)^{\left|\mu \right|l(\mu )}q^{l(\mu )}.$$
On rappelle la définition du polynôme $`P_{jk}`$ introduit en (1) et de la fonction symétrique monomiale $`\psi _\mu `$ (somme de tous les monômes différents ayant pour exposant une permutation de $`\mu `$).
###### Lemme 2.
Soit $`A=\{a_1,a_2,a_3,\mathrm{}\}`$ un alphabet (fini ou infini) quelconque. Pour tout $`i1`$ on pose $`X_i=_{aA}a^i`$. Alors pour tous entiers $`j,k0`$ on a
$$P_{jk}(X)=(1)^k\underset{|\mu |=j,l(\mu )=k}{}\psi _\mu (A).$$
###### Preuve.
L’égalité à établir est une identité algébrique entre polynômes en les $`a_i`$. Elle est entièremnt indépendante de la structure de $`\lambda `$-anneau de l’anneau des polynômes. Pour la démontrer dans le cadre de la théorie des $`\lambda `$-anneaux, nous pouvons donc choisir le statut de chacune des indéterminées $`a_i`$.
Nous pouvons par exemple supposer que tous les éléments de l’alphabet $`A`$ sont de rang 1. Compte-tenu de (3), la relation à démontrer devient dans ce cas
$$P_{jk}(X)=(1)^k\underset{|\mu |=j,l(\mu )=k}{}\psi _\mu [A^{\mathrm{}}].$$
(10)
Compte-tenu de (3), on a aussi dans ce cas
$$\psi ^i[A^{\mathrm{}}]=\psi ^i(A)=X_i,i1,$$
d’où pour toute partition $`\mu `$,
$$\psi ^\mu [A^{\mathrm{}}]=\psi ^\mu (A)=\underset{i1}{}X_{i}^{}{}_{}{}^{m_i(\mu )}.$$
La formule de Cauchy (7) implique alors
$$\begin{array}{cc}\hfill \mathrm{\Lambda }^j[qA^{\mathrm{}}]& =\underset{\left|\mu \right|=j}{}\frac{(1)^{jl(\mu )}}{z_\mu }\psi ^\mu [q]\psi ^\mu [A^{\mathrm{}}]\hfill \\ & =\underset{\left|\mu \right|=j}{}\frac{(1)^{jl(\mu )}}{z_\mu }\left(\underset{k1}{}\genfrac{}{}{0.0pt}{}{\mu }{k}q^k\right)\underset{i1}{}X_{i}^{}{}_{}{}^{m_i(\mu )}\hfill \\ & =(1)^j\underset{k1}{}P_{jk}(X)q^k.\hfill \end{array}$$
Et d’autre part on a aussi
$$\begin{array}{cc}\hfill \mathrm{\Lambda }^j[qA^{\mathrm{}}]& =\underset{\left|\mu \right|=j}{}\psi _\mu [A^{\mathrm{}}]\mathrm{\Lambda }^\mu [q]\hfill \\ & =\underset{\left|\mu \right|=j}{}\psi _\mu [A^{\mathrm{}}](1)^{jl(\mu )}q^{l(\mu )}.\hfill \end{array}$$
On en déduit (10) par comparaison. ∎
### 5.2 Méthode
Dans toute la suite de cet article, on considère un alphabet (fini ou infini) $`A=\{a_1,a_2,a_3,\mathrm{}\}`$. Pour le moment, nous ne faisons aucune hypothèse sur le statut des éléments de $`A`$. En particulier nous ne supposons pas que les $`a_k`$ sont de rang 1. Pour tout $`i1`$ on pose
$$X_i=\psi ^i(A)=\underset{aA}{}a^i.$$
On considère quatre éléments $`q^{},z,t,u`$. On suppose que $`z`$ est de type binomial et que $`q^{}=1+q`$ est de rang 1.
Pour démontrer l’identité du Théorème 1, on va montrer l’égalité des fonctions génératrices de ses deux membres. Plus précisément on écrit chaque membre de l’identité du Théorème 1 en changeant les $`X_i`$ en $`X_i`$, et on somme sur $`n`$ et $`r`$ après avoir multiplié par $`(t)^n`$ $`(q)^r`$.
L’égalité à démontrer devient
$$\begin{array}{c}\underset{nr1}{}\underset{|\mu |=n}{}(1)^{rl(\mu )}(t)^n(q)^r\frac{{\displaystyle \genfrac{}{}{0.0pt}{}{\mu }{r}}}{z_\mu }\underset{i1}{}\left(z\underset{k1}{}u^k\frac{(i)_k}{k!}X_k\right)^{m_i(\mu )}=\hfill \\ \hfill \underset{nr1}{}\underset{j0}{}(t)^n(q)^ru^j\left(\genfrac{}{}{0pt}{}{n+j1}{nr}\right)\left(\underset{k=0}{\overset{min(r,j)}{}}\left(\genfrac{}{}{0pt}{}{zj}{rk}\right)P_{jk}(X)\right).\end{array}$$
(11)
### 5.3 Membre de droite
Compte tenu du Lemme 2, le membre de droite de (11) s’écrit, en notant $`uA=\{ua_1,ua_2,ua_3,\mathrm{}\}`$,
$$\underset{nr1}{}\underset{\nu }{}(t)^n(q)^ru^{|\nu |}\left(\genfrac{}{}{0pt}{}{n+|\nu |1}{nr}\right)\left(\genfrac{}{}{0pt}{}{z|\nu |}{rl(\nu )}\right)(1)^{l(\nu )}\psi _\nu (A)=$$
$$\underset{\nu }{}(1)^{l(\nu )}\psi _\nu (uA)\underset{nrl(\nu )}{}(t)^n(q)^r\left(\genfrac{}{}{0pt}{}{n+|\nu |1}{nr}\right)\left(\genfrac{}{}{0pt}{}{z|\nu |}{rl(\nu )}\right).$$
Mais on a la relation suivante, qui est une autre façon d’écrire la formule classique du binôme :
$$\underset{ij}{}\left(\genfrac{}{}{0pt}{}{i1}{j1}\right)t^{ij}=\frac{1}{(1t)^j}.$$
On en déduit immédiatement
$$\underset{nr}{}(t)^n\left(\genfrac{}{}{0pt}{}{n+|\nu |1}{nr}\right)=\frac{(t)^r}{(1+t)^{|\nu |+r}}.$$
Le membre de droite de (11) s’écrit donc
$$\underset{\nu }{}(1)^{l(\nu )}\psi _\nu (uA)\left(\underset{rl(\nu )}{}\frac{(qt)^r}{(1+t)^{|\nu |+r}}\left(\genfrac{}{}{0pt}{}{z|\nu |}{rl(\nu )}\right)\right).$$
Ce qui peut se reformuler
$$\underset{\nu }{}\psi _\nu (uA)\frac{(qt)^{l(\nu )}}{(1+t)^{|\nu |+l(\nu )}}\left(\underset{k0}{}\left(\frac{qt}{1+t}\right)^k\left(\genfrac{}{}{0pt}{}{z|\nu |}{k}\right)\right).$$
Finalement le membre de droite de (11) s’écrit
$$\underset{\nu }{}\psi _\nu (uA)\frac{(qt)^{l(\nu )}}{(1+t)^{|\nu |+l(\nu )}}\left(1+\frac{qt}{1+t}\right)^{z|\nu |}.$$
Soit encore en posant $`y=qt/(1+t)`$,
$$\underset{\nu }{}\psi _\nu (uA)\frac{y^{l(\nu )}}{(1+t)^{|\nu |}}(1y)^{z|\nu |}.$$
(12)
### 5.4 Membre de gauche
Comme on a $`X_k=_{aA}a^k`$, la quantité suivante, écrite au membre de gauche de (11), devient
$$\begin{array}{cc}\hfill z\underset{k1}{}u^k\frac{(i)_k}{k!}X_k& =z\underset{k1}{}u^k\frac{(i)_k}{k!}\left(\underset{aA}{}a^k\right).\hfill \end{array}$$
On introduit alors l’alphabet
$$A^{}=\{\frac{1}{1ua_1},\frac{1}{1ua_2},\frac{1}{1ua_3},\mathrm{}\}=\left\{\frac{1}{1ua},aA\right\}.$$
Nous faisons désormais l’hypothèse suivante : chaque élément $`\frac{1}{1ua}`$ est de rang 1. Sous cette hypothèse on a
$$\begin{array}{cc}\hfill \psi ^i[\underset{aA}{}\frac{1}{1ua}]& =\underset{aA}{}(1ua)^i\hfill \\ & =\underset{aA}{}\left(\underset{k0}{}\frac{(i)_k}{k!}u^ka^k\right)\hfill \\ & =\underset{aA}{}\left(1+\underset{k1}{}\frac{(i)_k}{k!}u^ka^k\right).\hfill \end{array}$$
On introduit l’élément
$$B=z\underset{aA}{}\frac{ua}{1ua}=z+\underset{aA}{}\left(1\frac{1}{1ua}\right).$$
On a ainsi
$$z\underset{k1}{}u^k\frac{(i)_k}{k!}X_k=\psi ^i[B].$$
Pour toute partition $`\mu `$, on en déduit
$$\underset{i1}{}\left(z\underset{k1}{}u^k\frac{(i)_k}{k!}X_k\right)^{m_i(\mu )}=\psi ^\mu [B].$$
Compte-tenu de cette relation, le membre de gauche de (11) s’écrit
$$\begin{array}{cc}\hfill \underset{nr1}{}\underset{|\mu |=n}{}(1)^{rl(\mu )}(t)^n(q)^r\frac{{\displaystyle \genfrac{}{}{0.0pt}{}{\mu }{r}}}{z_\mu }\psi ^\mu [B]& =\underset{n1}{}t^n\underset{\left|\mu \right|=n}{}\frac{(1)^{nl(\mu )}}{z_\mu }\left(\underset{r1}{}\genfrac{}{}{0.0pt}{}{\mu }{r}q^r\right)\psi ^\mu [B]\hfill \\ & =\underset{n1}{}t^n\underset{\left|\mu \right|=n}{}\frac{(1)^{nl(\mu )}}{z_\mu }\psi ^\mu [q]\psi ^\mu [B]\hfill \\ & =\underset{n1}{}t^n\mathrm{\Lambda }^n[qB]\hfill \\ & =\lambda _t[qB].\hfill \end{array}$$
La démonstration sera terminée en prouvant que le développement (12) est exactement la décomposition de $`\lambda _t[qB]`$ sur la base des fonctions monomiales $`\psi _\nu (uA)`$.
### 5.5 Développement de $`\lambda _t[qB]`$.
On maintient les notations précédentes en faisant le changement de variables $`uaa`$. On considère un alphabet $`A=\{a_1,a_2,a_3,\mathrm{}\}`$ et trois éléments $`q^{},z,t`$ avec les hypothèses suivantes :
\- on suppose que $`z`$ est de type binomial et que $`q^{}=1+q`$ est de rang 1,
\- on suppose que pour tout $`aA`$, l’élément $`a^{}=\frac{1}{1a}`$ est de rang 1.
On a maintenant
$$B=z\underset{aA}{}\frac{a}{1a}=z+\underset{aA}{}\left(1\frac{1}{1a}\right).$$
Nous allons démontrer le Théorème 1 sous la forme suivante.
###### Théorème 3.
En posant $`y=qt/(1+t)`$, on a
$$\lambda _t[qB]=\underset{\nu }{}\psi _\nu (A)\frac{y^{l(\nu )}}{(1+t)^{|\nu |}}(1y)^{z|\nu |}.$$
###### Preuve.
On a d’abord
$$\lambda _t[qB]=\lambda _t[qz+q\underset{aA}{}(1a^{})]$$
Comme $`z`$ est de type binomial, on a
$$\lambda _t[qz]=(\lambda _t[q])^z.$$
Et d’autre part la relation (9) implique
$$\lambda _t[q]=1+q\underset{i1}{}(1)^{i1}t^i=1+\frac{qt}{1+t}.$$
Compte-tenu de (5) on en déduit
$$\begin{array}{cc}\hfill \lambda _t[qB]& =\lambda _t[qz]\lambda _t[q\underset{aA}{}(1a^{})]\hfill \\ & =(1y)^z\lambda _t[q\underset{aA}{}(1a^{})]\hfill \\ & =(1y)^z\underset{aA}{}\lambda _t[q(1a^{})].\hfill \end{array}$$
Maintenant on a $`q(1a^{})=(q^{}1)(1a^{})=q^{}1q^{}a^{}+a^{}`$. Les éléments $`q^{}`$ et $`a^{}`$ étant de rang 1, les relations (5) et (8) impliquent
$$\lambda _t[q(1a^{})]=\frac{\lambda _t[q^{}]}{\lambda _t[1]}\frac{\lambda _t[a^{}]}{\lambda _t[q^{}a^{}]}=\frac{1+tq^{}}{1+t}\frac{1+ta^{}}{1+tq^{}a^{}}=(1y)\frac{1+ta^{}}{1+t(1+q)a^{}}.$$
Finalement on obtient
$$\begin{array}{cc}\hfill \lambda _t[qB]& =(1y)^z\underset{aA}{}(1y)\frac{1+ta}{1+t(1+q)a}\hfill \\ & =(1y)^z\underset{aA}{}(1y)\left(1\frac{qt}{1+t+qta}\right).\hfill \end{array}$$
Posons alors
$$v=\frac{1}{(1+t)(1y)}=\frac{1}{1+t+qt}.$$
La relation précédente devient
$$\lambda _t[qB]=(1y)^z\underset{aA}{}\left(1+y\frac{va}{1va}\right).$$
Maintenant on a
$$\underset{aA}{}\left(1+y\frac{va}{1va}\right)=\underset{NA}{}\underset{aN}{}y\frac{va}{1va}.$$
Nous allons utiliser la propriété suivante, qui se vérifie facilement :
$$\underset{\begin{array}{c}NA\hfill \\ \mathrm{card}N=n\hfill \end{array}}{}\underset{aN}{}y\frac{va}{1va}=y^n\underset{l(\nu )=n}{}v^{|\nu |}\psi _\nu (A).$$
Soit encore
$$\lambda _t[qB]=(1y)^z\underset{\nu }{}\frac{y^{l(\nu )}}{(1+t)^{|\nu |}(1y)^{|\nu |}}\psi _\nu (A).$$
On conclut immédiatement. ∎
## 6 Application
Les Théorèmes 1 et 2 peuvent permettre d’obtenir des identités remarquables en spécialisant les indéterminées $`X_i`$ et $`z`$.
Nous revenons seulement ici sur les conjectures de , rencontrées en étudiant les polynômes symétriques décalés . Soit $`\alpha `$ un nombre réel positif. Pour toute partition $`\lambda `$ et tout entier $`k0`$, on note
$$d_k(\lambda )=\underset{(i,j)\lambda }{}\left(j1\frac{i1}{\alpha }\right)^k.$$
On introduit la généralisation suivante de la “factorielle ascendante”:
$$(z)_\lambda =\underset{(i,j)\lambda }{}\left(z+j1\frac{i1}{\alpha }\right).$$
Pour tous entiers $`j,k0`$ on pose
$$F_{jk}(\lambda )=P_{jk}(d_1(\lambda ),d_2(\lambda ),d_3(\lambda ),\mathrm{}).$$
C’est-à-dire qu’on choisit la spécialisation suivante
$$X_k=d_k(\lambda ),k1.$$
En d’autres termes, l’alphabet $`A`$ tel que $`X_k=_{aA}a^k`$ est alors
$$A_\lambda =\left\{j1\frac{i1}{\alpha },(i,j)\lambda \right\}.$$
###### Théorème 4.
Soient $`x,y`$ deux indéterminées indépendantes. Pour toute partition $`\lambda `$ on a
$$\frac{(yx)_\lambda }{(y)_\lambda }=\underset{i0}{}\underset{j0}{}(1)^{i+j}\frac{x^i}{y^{i+j}}\left(\underset{k=0}{\overset{min(i,j)}{}}\left(\genfrac{}{}{0pt}{}{|\lambda |j}{ik}\right)F_{jk}(\lambda )\right).$$
###### Preuve.
On montre comme dans (p. 464) que
$$\frac{(yx)_\lambda }{(y)_\lambda }=\underset{\mu }{}v^{|\mu |}\frac{(1)^{|\mu |l(\mu )}}{z_\mu }\underset{i1}{}\left(\underset{p0}{}u^p\frac{(i)_p}{p!}d_p(\lambda )\right)^{m_i(\mu )},$$
avec $`v=x/y`$ et $`u=1/y`$. On écrit le Théorème 2 spécialisé avec $`X_k=d_k(\lambda )`$ et $`z=d_0(\lambda )=|\lambda |`$. ∎
Il est important de noter que la sommation a lieu sur tout $`j0`$ et pas seulement sur $`|\lambda |j0`$. Le degré en $`x`$ du membre de gauche étant clairement $`|\lambda |`$, on obtient pour tout $`i>|\lambda |`$,$`j0`$,
$$\underset{k=0}{\overset{min(i,j)}{}}\left(\genfrac{}{}{0pt}{}{|\lambda |j}{ik}\right)F_{jk}(\lambda )=0.$$
En effet c’est seulement lorsque l’alphabet $`A_\lambda `$ est infini que les indéterminées $`d_k(\lambda )`$ sont indépendantes.
Le cas où $`\lambda `$ est une partition-ligne $`(n)`$ correspond au développement en série de la formule classique de Chu-Vandermonde . |
warning/0001/cond-mat0001062.html | ar5iv | text | # The Interacting Impurity Josephson Junction: Variational Wavefunctions and Slave Boson Mean Field Theory
## I Introduction
The Ambegaokar-Baratoff formula, $`I_\mathrm{c}=\pi \mathrm{\Delta }/2eR`$, relates the critical current in a Josephson junction to the superconducting gap $`\mathrm{\Delta }`$ and the normal state junction resistance $`R`$. This result is perturbative in the electron tunneling amplitude $`t`$; the normal state conductance is proportional to $`|t|^2`$ when $`t`$ is small. In certain instances, however, it may be that the tunneling is mediated through a magnetic impurity, rather than taking place directly from one superconductor to the other. This situation has been considered by a number of authors . The principal result is that magnetic impurity-mediated tunneling results in a negative contribution to $`I_\mathrm{c}`$. As Kulik originally argued , the magnetic impurity gives rise to an effective spin-flip hopping amplitude $`t_{\mathrm{sf}}`$ between the superconductors. Since spin-flip tunneling results in a sign change of the Cooper pair singlet ($``$ to $``$) , the critical current is $`I_\mathrm{c}|t|^2|t_{\mathrm{sf}}|^2`$. When $`I_\mathrm{c}<0`$, one has a $`\pi `$-junction, for which the ground state energy is minimized when the phase difference between the superconductors is $`\delta =\pi `$. Such $`\pi `$-junctions break time-reversal symmetry ($`𝒯`$), and a ring containing a single $`\pi `$-junction will enclose trapped flux .
This analysis suffices in the perturbative limit where $`t`$ is small. If the impurity level energy is $`\epsilon _0^{}`$ and the Coulomb integral is $`U`$, the condition for a magnetic ground state (and a $`\pi `$-junction) is $`U>\epsilon _0^{}>0`$ . In the nonperturbative regime, a new energy scale arises: the bare impurity level width $`\mathrm{\Gamma }\pi \rho |t|^2`$, where $`t`$ is the electrode-impurity hopping matrix element and $`\rho `$ the electrode density of states. The $`T=0`$ phase diagram as a function of $`\epsilon _0^{}/\mathrm{\Delta }`$, $`U/\mathrm{\Delta }`$, and $`\mathrm{\Gamma }/\mathrm{\Delta }`$ was investigated by the authors within the Hartree-Fock (HF) approximation, where it was found that $`\pi `$-junction behavior occurs for $`U>\epsilon _0^{}>0`$, provided $`\mathrm{\Gamma }`$ is sufficiently small ($`\mathrm{\Gamma }<U`$). However, when $`\mathrm{\Delta }=0`$, it is known that the HF approximation is unable to describe the formation of a Kondo singlet at energy scales below $`T_\mathrm{K}^{}W\mathrm{exp}(\pi |\epsilon _0^{}|/2\mathrm{\Gamma })`$, where $`W`$ is the half-bandwidth in the electrodes. In our problem, then, one expects that the Kondo effect will be mitigated whenever $`\mathrm{\Delta }>T_\mathrm{K}^{}`$. Roughly speaking, if $`\mathrm{\Delta }>T_\mathrm{K}`$, the ground state of the system is a Kramers doublet, and breaks time reversal symmetry, whereas if $`T_\mathrm{K}^{}>\mathrm{\Delta }`$, the ground state is a hybrid singlet formed from electrons on the impurity and in the superconducting electrodes .
Recently, Clerk and Ambegaokar applied a generalization of the non-crossing approximation (NCA), a partial summation scheme, to attack this problem in the $`U\mathrm{}`$ limit . Adaptation of this method to the interacting impurity Josephson junction allows one to see a transition from $`0`$-junction to $`\pi `$-junction when the superconducting gap becomes comparable to the Kondo temperature.
In this paper, we also explore the $`U\mathrm{}`$ limit, using two approaches. The first is a variational wave function calculation, similar to that used by Varma and Yafet in the Kondo problem . We generalize this wave function in two respects. Firstly, there are two superconductors connected to the impurity. This is in contrast with the usual setup of Kondo problem where we have only metallic electrode. Secondly, the spin-$`\frac{1}{2}`$ state must be considered as well. We find, in agreement with ref. , that a first order transition occurs at $`T_\mathrm{K}^{}/\mathrm{\Delta }1`$. We also show how this transition may be precipitated by a change in the phase difference $`\delta `$, as we found in ref. .
The second calculation we describe is the slave boson mean field theory. While the this method does confirm that the singlet state gives a $`0`$-junction and the Kramers doublet a $`\pi `$-junction, within the static slave boson mean field solution the singlet is always lower in energy (or the states are degenerate). The NCA, which goes beyond the mean field level, is able to describe the transition.
## II Variational Wavefunction Approach
We start with the grand canonical Hamiltonian $`𝒦\mu N`$,
$`𝒦`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{q}{}}\{\xi _{q\alpha }^{}(\psi _{q\alpha }^{}\psi _{q\alpha }^{}+\psi _{q\alpha }^{}\psi _{q\alpha }^{})`$ (4)
$`+\mathrm{\Delta }_\alpha ^{}\left(e^{i\delta _\alpha }\psi _{q\alpha }^{}\psi _{q\alpha }^{}+e^{i\delta _\alpha }\psi _{q\alpha }^{}\psi _{q\alpha }^{}\right)`$
$`{\displaystyle \frac{1}{\sqrt{𝒩_\alpha }}}{\displaystyle \underset{\sigma =,}{}}(t_\alpha ^{}\psi _{q\alpha \sigma }^{}c_\sigma ^{}+t_\alpha ^{}c_\sigma ^{}\psi _{q\alpha \sigma }^{})\}`$
$`+\epsilon _0^{}\left(c_{}^{}c_{}^{}+c_{}^{}c_{}^{}\right)+Uc_{}^{}c_{}^{}c_{}^{}c_{}^{},`$
where $`\alpha `$ labels the electrode, $`\xi _{q\alpha }`$ is the dispersion in the $`\alpha `$ electrode relative to the chemical potential, $`\mathrm{\Delta }_\alpha `$ is the modulus of the superconducting gap and $`\delta _\alpha `$ its phase, $`𝒩_\alpha `$ is the number of unit cells in electrode $`\alpha `$, $`t_\alpha `$ is the hopping amplitude from electrode $`\alpha `$ to the impurity, and $`\epsilon _0^{}`$ is the bare impurity energy. We set $`U\mathrm{}`$, which leads to the constraint $`_\sigma c_\sigma ^{}c_\sigma ^{}1`$.
Defining the angle $`\theta _{q\alpha }\mathrm{tan}^1(\mathrm{\Delta }_\alpha /\xi _{q\alpha }^{})`$ and the usual BCS coherence factors $`u_{q\alpha }^{}=\mathrm{cos}(\frac{1}{2}\theta _{q\alpha }^{})`$, $`v_{q\alpha }^{}=\mathrm{sin}(\frac{1}{2}\theta _{q\alpha }^{})\mathrm{exp}(i\delta _\alpha )`$, we express $`𝒦=𝒦_0+𝒦_1`$ in terms of the Bogoliubov quasiparticle operators:
$`𝒦_0`$ $`=`$ $`{\displaystyle \underset{q,\alpha }{}}E_{q\alpha }^{}(\gamma _{q\alpha }^{}\gamma _{q\alpha }^{}+\gamma _{q\alpha }^{}\gamma _{q\alpha }^{})+\epsilon _0^{}(c_{}^{}c_{}^{}+c_{}^{}c_{}^{})`$
$`+Uc_{}^{}c_{}^{}c_{}^{}c_{}^{}`$
$`𝒦_1`$ $`=`$ $`{\displaystyle \underset{q,\alpha }{}}t_{q\alpha }^{}\{u_{q\alpha }^{}(\gamma _{q\alpha }^{}c_{}^{}+\gamma _{q\alpha }^{}c_{}^{}+c_{}^{}\gamma _{q\alpha }^{}+c_{}^{}\gamma _{q\alpha }^{})`$
$`+v_{q\alpha }^{}(\gamma _{q\alpha }^{}c_{}^{}\gamma _{q\alpha }^{}c_{}^{})+v_{q\alpha }^{}(\gamma _{q\alpha }^{}c_{}^{}\gamma _{q\alpha }^{}c_{}^{})\},`$
where $`E_{q\alpha }^{}=\sqrt{\xi _{q\alpha }^2+\mathrm{\Delta }_\alpha ^2}`$, and $`t_{q\alpha }^{}|t_\alpha ^{}|/\sqrt{𝒩_\alpha }`$.
We now two variational many-body states for the $`U=\mathrm{}`$ limit: a singlet,
$`|\mathrm{S}`$ $``$ $`\{A+{\displaystyle \underset{q,\alpha }{}}{\displaystyle \frac{1}{\sqrt{2}}}B_{q\alpha }^{}(\gamma _{q\alpha }^{}c_{}^{}\gamma _{q\alpha }^{}c_{}^{})`$ (6)
$`+{\displaystyle \underset{\genfrac{}{}{0pt}{}{q,\alpha }{q^{},\alpha ^{}}}{}}C_{qq^{}}^{\alpha \alpha ^{}}\gamma _{q\alpha }^{}\gamma _{q^{}\alpha ^{}}^{}\}|\mathrm{\hspace{0.17em}0},`$
and a doublet,
$`|\mathrm{D}\{\stackrel{~}{A}c_{}^{}+{\displaystyle \underset{q,\alpha }{}}\stackrel{~}{B}_{q\alpha }^{}\gamma _{q\alpha }^{}+{\displaystyle \underset{\genfrac{}{}{0pt}{}{q,\alpha }{q^{},\alpha ^{}}}{}}[\stackrel{~}{C}_{qq^{}}^{\alpha \alpha ^{}}\gamma _{q\alpha }^{}\gamma _{q^{}\alpha ^{}}^{}c_{}^{}`$ (7)
$`+{\displaystyle \frac{1}{\sqrt{3}}}\stackrel{~}{D}_{qq^{}}^{\alpha \alpha ^{}}(\gamma _{q\alpha }^{}\gamma _{q^{}\alpha ^{}}^{}c_{}^{}\gamma _{q\alpha }^{}\gamma _{q^{}\alpha ^{}}^{}c_{}^{})]\}|\mathrm{\hspace{0.17em}0},`$ (8)
where $`|\mathrm{\hspace{0.17em}0}`$ is the fermion vacuum (the other doublet state $`|\mathrm{D}`$ is obtained by rotating the spins by $`\pi `$ about the $`y`$-axis). Here, $`C_{q^{}q}^{\alpha ^{}\alpha }=C_{qq^{}}^{\alpha \alpha ^{}}`$, $`\stackrel{~}{C}_{q^{}q}^{\alpha ^{}\alpha }=\stackrel{~}{C}_{qq^{}}^{\alpha \alpha ^{}}`$, and $`\stackrel{~}{D}_{q^{}q}^{\alpha ^{}\alpha }=\stackrel{~}{D}_{qq^{}}^{\alpha \alpha ^{}}`$. We next set to zero the variations
$$\delta \mathrm{\Psi }|𝒦_0+𝒦_1|\mathrm{\Psi }E\delta \mathrm{\Psi }|\mathrm{\Psi }=0,$$
(9)
where $`|\mathrm{\Psi }`$ is $`|\mathrm{S}`$ or $`|\mathrm{D}`$, to obtain equations relating the variational coefficients. We find that $`A`$ and the matrix $`C_{qq^{}}^{\alpha \alpha ^{}}`$ may be expressed in terms of the coefficients $`B_{q\alpha }^{}`$, and similarly $`\stackrel{~}{A}`$ and the matrices $`\stackrel{~}{C}_{qq^{}}^{\alpha \alpha ^{}}`$ and $`\stackrel{~}{D}_{qq^{}}^{\alpha \alpha ^{}}`$ may be expressed in terms of the coefficients $`\stackrel{~}{B}_{q\alpha }^{}`$. We then obtain the two eigenvalue equations for the singlet and doublet energies $`E`$ and $`\stackrel{~}{E}`$, respectively:
$`{\displaystyle \underset{q^{}\alpha ^{}}{}}`$ $`\left[{\displaystyle \frac{2v_{q\alpha }^{}v_{q^{}\alpha ^{}}^{}t_{q\alpha }^{}t_{q^{}\alpha ^{}}^{}}{E}}+{\displaystyle \frac{u_{q\alpha }^{}u_{q^{}\alpha ^{}}^{}t_{q\alpha }^{}t_{q^{}\alpha ^{}}}{EE_{q\alpha }^{}E_{q^{}\alpha ^{}}^{}}}\right]B_{q^{}\alpha ^{}}^{}`$ (11)
$`=\left[EE_{q\alpha }^{}\epsilon _0^{}{\displaystyle \underset{q^{},\alpha ^{}}{}}{\displaystyle \frac{u_{q^{}\alpha ^{}}^2t_{q^{}\alpha ^{}}^2}{EE_{q\alpha }^{}E_{q^{}\alpha ^{}}^{}}}\right]B_{q\alpha }^{}`$
and
$`{\displaystyle \underset{q^{}\alpha ^{}}{}}`$ $`\left[{\displaystyle \frac{u_{q\alpha }^{}u_{q^{}\alpha ^{}}^{}t_{q\alpha }^{}t_{q^{}\alpha ^{}}^{}}{\stackrel{~}{E}\epsilon _0^{}}}+{\displaystyle \frac{v_{q\alpha }^{}v_{q^{}\alpha ^{}}^{}t_{q\alpha }^{}t_{q^{}\alpha ^{}}}{\stackrel{~}{E}\epsilon _0^{}E_{q\alpha }^{}E_{q^{}\alpha ^{}}^{}}}\right]\stackrel{~}{B}_{q^{}\alpha ^{}}^{}`$ (13)
$`=\left[\stackrel{~}{E}E_{q\alpha }^{}{\displaystyle \underset{q^{},\alpha ^{}}{}}{\displaystyle \frac{2|v_{q^{}\alpha ^{}}|^2t_{q^{}\alpha ^{}}^2}{\stackrel{~}{E}\epsilon _0^{}E_{q\alpha }^{}E_{q^{}\alpha ^{}}^{}}}\right]\stackrel{~}{B}_{q\alpha }^{}`$
We solve these equations numerically for the symmetric case $`\mathrm{\Delta }_\mathrm{L}^{}=\mathrm{\Delta }_\mathrm{R}^{}=\mathrm{\Delta }`$, $`t_\mathrm{L}^{}=t_\mathrm{R}^{}=t`$, $`\mathrm{\Gamma }_\mathrm{L}^{}=\mathrm{\Gamma }_\mathrm{R}^{}=\mathrm{\Gamma }`$. The normal state of each electrode is described by a flat band of width $`2W`$; we use $`W/\mathrm{\Delta }=10`$ in our calculations, but the general features are rather insensitive to the value of $`W`$. Energy versus phase difference, $`E(\delta `$), is plotted in figure 1 for $`\mathrm{\Delta }=1`$, $`\mathrm{\Gamma }=0.3\pi `$, for four different values of $`\epsilon _0^{}`$. When $`\epsilon _0^{}=1.6`$, the singlet state is the ground state for all $`\delta `$, while for $`\epsilon _0^{}=2.0`$, the ground state is always the doublet. These are $`0`$\- and $`\pi `$-junctions, respectively. For $`\epsilon _0^{}=1.8`$, the curves cross, and the ground state energy has a kink as a function of $`\delta `$. Both $`\delta =0`$ and $`\delta =\pi `$ are local minima in $`E(\delta )`$, with $`\delta =0`$ the global minimum. Using the terminology of ref. , this is a $`0^{}`$-junction. The final case of the $`\pi ^{}`$-junction is reflected in the curves for $`\epsilon _0^{}=1.9`$, where $`\delta =0`$ is a local minimum and $`\delta =\pi `$ the global minimum of the energy.
The phase diagram is displayed in fig. 2. In the HF treatment , one also finds a transition from $`\pi `$ to $`\pi ^{}`$ to $`0^{}`$ to $`0`$-junction phases as $`\mathrm{\Gamma }`$ is increased from zero. However, the critical value of $`\mathrm{\Gamma }`$ there turns out to be inversely proportional to $`\mathrm{ln}U`$. There are two reasons for this: first, the HF calculations assumed an infinite bandwidth, and second, HF is unable to describe the formation of a Kondo singlet. Here, the transition from $`0`$-junction to $`\pi `$-junction occurs roughly for $`\mathrm{\Delta }T_\mathrm{K}^{}`$, in agreement with ref. . Indeed, when $`0<\epsilon _0^{},\mathrm{\Gamma }\mathrm{\Delta }`$ this relation can be derived within perturbation theory. The perturbative value of energy of the singlet state is
$$E_0\frac{4\mathrm{\Gamma }}{\pi }\mathrm{ln}\frac{2W}{\mathrm{\Delta }},$$
(14)
while the energy of the doublet is
$$E_{}\epsilon _0^{}\frac{2\mathrm{\Gamma }}{\pi }\mathrm{ln}\frac{2W}{\mathrm{\Delta }}.$$
(15)
Comparing these two we find that the symmetry of the ground state changes when
$$\mathrm{\Delta }=2W\mathrm{exp}(\pi |\epsilon _0^{}|/2\mathrm{\Gamma })=2T_\mathrm{K}^{}.$$
(16)
This predicts a straight line for $`\mathrm{\Gamma }`$ versus $`\epsilon _0^{}`$ with a slope $`\mathrm{\Gamma }/(\epsilon _0^{})=\pi /2\mathrm{ln}(2W/\mathrm{\Delta })`$. With $`W=10\mathrm{\Delta }`$, the slope is $`0.524`$, and the line would be bounded by the $`00^{}`$ and $`\pi \pi ^{}`$ phase boundaries in fig. 2 (it is not shown for the purposes of clarity).
## III Slave Boson Mean Field Theory
We consider an extension of the model of eqn. (4) by introducing a flavor -1 $`m`$ which runs from $`1`$ to $`N_\mathrm{f}`$. The Hamiltonian is then
$`𝒦`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{m=1}{\overset{N_\mathrm{f}}{}}}{\displaystyle \underset{q}{}}\{{\displaystyle \underset{\sigma =,}{}}(\xi _{q\alpha }^{}+\mu _\mathrm{B}B\sigma )\psi _{q\alpha m\sigma }^{}\psi _{q\alpha m\sigma }^{}+`$ (20)
$`\mathrm{\Delta }_\alpha ^{}\left(e^{i\delta _\alpha }\psi _{q\alpha m}^{}\psi _{q\alpha m}^{}+e^{i\delta _\alpha }\psi _{q\alpha m}^{}\psi _{q\alpha m}^{}\right)`$
$`{\displaystyle \frac{1}{\sqrt{𝒩_\alpha }}}{\displaystyle \underset{\sigma =,}{}}(t_\alpha ^{}\psi _{q\alpha m\sigma }^{}c_{m\sigma }^{}+t_\alpha ^{}c_{m\sigma }^{}\psi _{q\alpha m\sigma }^{})\}`$
$`+{\displaystyle \underset{m=1}{\overset{N_\mathrm{f}}{}}}{\displaystyle \underset{\sigma =,}{}}(\epsilon _0^{}+\mu _\mathrm{B}B\sigma )c_{m\sigma }^{}c_{m\sigma }^{},`$
where $`\mu _\mathrm{B}B`$ is the Zeeman energy. The impurity level occupancy satisfies the constraint $`_{m,\sigma }c_{m\sigma }^{}c_{m\sigma }^{}r`$. We refer to this as model I. In a slightly different large-$`N_\mathrm{f}`$ extension (model II), we rescale the hopping amplitudes $`t_\alpha t_\alpha /\sqrt{N_\mathrm{f}}`$ and write the constraint as $`_{m,\sigma }c_{m\sigma }^{}c_{m\sigma }^{}rN_\mathrm{f}`$. In both cases, $`0r2`$.
Introducing a slave boson $`b`$ and a Lagrange multiplier $`\lambda `$ to impose the constraint, we evaluate the impurity contribution to the free energy at the mean field level, assuming both $`b`$ and $`\lambda `$ to be static. The impurity free energy per flavor is then
$`F_{\mathrm{imp}}/N_\mathrm{f}`$ $`=`$ $`\epsilon \mu _\mathrm{B}B+p(\epsilon \epsilon _0^{})(|b|^2r)`$ (22)
$`+{\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{d\omega }{\pi }}f(\omega +\mu _\mathrm{B}B)\mathrm{Im}\mathrm{ln}H(\omega +i0^+)`$
$`H(\omega )`$ $`=`$ $`\omega ^2\epsilon ^2(2|b|^2\mathrm{\Gamma }_\mathrm{a}^{})^2+{\displaystyle \frac{\left[2|b|^2\mathrm{\Gamma }_\mathrm{g}^{}\mathrm{\Delta }\mathrm{sin}(\frac{1}{2}\delta )\right]^2}{\mathrm{\Delta }^2\omega ^2}}`$ (24)
$`+{\displaystyle \frac{4|b|^2\mathrm{\Gamma }_\mathrm{a}^{}\omega ^2}{\sqrt{\mathrm{\Delta }^2\omega ^2}}}`$
where $`\mathrm{\Gamma }_\mathrm{a}^{}=\frac{1}{2}`$, $`\mathrm{\Gamma }_\mathrm{g}^{}=\sqrt{\mathrm{\Gamma }_\mathrm{L}^{}\mathrm{\Gamma }_\mathrm{R}^{}}`$ with $`\mathrm{\Gamma }_\alpha =\pi \rho _\alpha |t_\alpha |^2`$ ($`\rho _\alpha `$ is the bare density of states per unit cell in electrode $`\alpha `$), $`\delta =\delta _\mathrm{L}^{}\delta _\mathrm{R}^{}`$ is the phase difference between the two superconducting electrodes (we assume $`\mathrm{\Delta }_1=\mathrm{\Delta }_2\mathrm{\Delta }`$), and $`f(\omega )=[\mathrm{exp}(\omega /T)+1]^1`$ is the Fermi function. The Lagrange multiplier $`\lambda `$ is absorbed into the renormalized impurity level energy, $`\epsilon \epsilon _0^{}+\lambda `$. The factor $`p`$ in the second term is $`1/N_\mathrm{f}`$ in model I, while $`p=1`$ in model II. Hence it is model II which generates a true large-$`N_\mathrm{f}`$ expansion,u with all terms in the impurity free energy $`F_{\mathrm{imp}}`$ proportional to $`N_\mathrm{f}`$.
From $`H(0)<0`$, $`H(\mathrm{\Delta }^{})=+\mathrm{}`$, and $`dH/d\omega >0`$, we conclude that there is a unique solution to the equation $`H(\omega )=0`$ on the interval $`\omega [0,\mathrm{\Delta }]`$. Call this root $`\mathrm{\Omega }`$. We use the external field $`B`$ as a Lagrange multiplier so that we may fix the total value of $`S^z`$. Making a Legendre transformation to $`G_{\mathrm{imp}}^{}(S^z)=F_{\mathrm{imp}}^{}(B)2\mu _\mathrm{B}BS^z`$, with $`S^z=\frac{1}{2},0,+\frac{1}{2}`$, we set $`F_{\mathrm{imp}}^{}/(\mu _\mathrm{B}B)=0`$, and obtain, at $`T=0`$ (assuming $`0<\mu _\mathrm{B}B<\mathrm{\Delta }`$),
$`G_{\mathrm{imp}}^0/N_\mathrm{f}`$ $`=`$ $`\epsilon +p(\epsilon \epsilon _0^{})(|b|^2r)+𝒜`$ (25)
$`𝒜`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _\mathrm{\Delta }^\omega _{}}𝑑\omega \mathrm{Im}\mathrm{ln}H(\omega i0^+)\mathrm{\Omega }\delta _{S^z,0},`$ (26)
where $`\omega _{}=\sqrt{W^2+\mathrm{\Delta }^2}`$; $`W`$ is the half-bandwidth in the electrodes. The mean field equations, obtained by setting $`G_{\mathrm{imp}}^0/|b|^2=0`$ and $`G_{\mathrm{imp}}^0/\epsilon =0`$ are
$`1+p(|b|^2q)+{\displaystyle \frac{𝒜}{\epsilon }}`$ $`=`$ $`0`$ (27)
$`p(\epsilon \epsilon _0^{})+{\displaystyle \frac{𝒜}{|b|^2}}`$ $`=`$ $`0,`$ (28)
and the Josephson current is
$$I=\frac{2e}{\mathrm{}}\frac{𝒜}{\delta }.$$
(29)
Thus, for $`0\mu _\mathrm{B}B<\mathrm{\Omega }`$, the ground state has $`S^z=0`$, while for $`\mathrm{\Omega }<|\mu _\mathrm{B}B|<\mathrm{\Delta }`$, the ground state has $`S^z=\frac{1}{2}\mathrm{sgn}(B)`$.
For $`p=r=1`$, we can show that the ground state energies satisfy $`\mathrm{\Delta }G_{\mathrm{imp}}^0G_{\mathrm{imp}}^0G_{\mathrm{imp}}^0>0`$, which means that the mean field slave boson theory cannot describe the $`\pi `$ or $`\pi ^{}`$ phases. The energy difference $`\mathrm{\Delta }G_{\mathrm{imp}}^0`$ is minimized at $`\delta =\pi `$. The value of $`\mathrm{\Delta }G_{\mathrm{imp}}^0`$ is an increasing function of the bare impurity level energy $`\epsilon _0^{}`$. However, there are no solutions to the mean field equations for $`\epsilon _0^{}<\epsilon _0^{\mathrm{min}}=4\pi ^1\mathrm{\Gamma }_\mathrm{a}^{}\mathrm{cosh}^1(\omega _{}/\mathrm{\Delta })2\mathrm{\Gamma }_\mathrm{a}^{}\delta _{S^z,0}`$. Rather, the endpoint solution $`\epsilon =|b|^2=0`$ holds. Thus, the best we can do is $`\mathrm{\Delta }G_{\mathrm{imp}}^0=0`$, but in this case $`G_{\mathrm{imp}}^0=\epsilon _0^{}`$, independent of $`\delta `$, and there is no Josephson current.
## IV Large $`\mathrm{\Delta }`$ Limit
The case $`\mathrm{\Delta }\mathrm{}`$ (with $`U`$ finite) may be solved exactly. We begin with the Hamiltonian of eqn. (4), integrating out the fermion degrees of freedom in the superconductors . This generates an induced action,
$`S_{\mathrm{ind}}`$ $`=`$ $`{\displaystyle \underset{\omega _m}{}}\overline{\mathrm{\Psi }}_i(\omega _m)[\sigma ^z𝒢(\omega _m)\sigma ^z]_{ij}\mathrm{\Psi }_j(\omega _m)`$ (30)
$`\mathrm{\Psi }(\omega _m)`$ $`=`$ $`\left(\begin{array}{c}c_{}^{}(\omega _m)\\ \overline{c}_{}^{}(\omega _m)\end{array}\right)`$ (31)
$`𝒢(\omega _m)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \frac{\mathrm{\Gamma }_\alpha }{\sqrt{\omega _m^2+\mathrm{\Delta }_\alpha ^2}}}\left(\begin{array}{cc}i\omega _m& \mathrm{\Delta }e^{i\delta _\alpha }\\ \mathrm{\Delta }e^{i\delta _\alpha }& i\omega _m\end{array}\right).`$ (32)
If the dynamics all occur on frequency scales $`\omega \mathrm{\Delta }`$, we may ignore the Matsubara frequencies $`\omega _m`$ in comparison with $`\mathrm{\Delta }`$. Adding the induced action to the bare action for the impurity, we find the resultant action is that for a Hamiltonian
$$_{\mathrm{eff}}=\epsilon _0^{}(c_{}^{}c_{}^{}+c_{}^{}c_{}^{})\chi c_{}^{}c_{}^{}\chi ^{}c_{}^{}c_{}^{}+Uc_{}^{}c_{}^{}c_{}^{}c_{}^{},$$
(33)
where
$$\chi =\mathrm{\Gamma }_\mathrm{L}^{}e^{i\delta _\mathrm{L}^{}}+\mathrm{\Gamma }_\mathrm{R}^{}e^{i\delta _\mathrm{R}^{}}.$$
(34)
The ground state energy is $`E_\mathrm{D}^0=\epsilon _0^{}`$ for the Kramers doublet, and
$$E_\mathrm{S}^0=\epsilon _0^{}+\frac{1}{2}U\sqrt{(\epsilon _0^{}+\frac{1}{2}U)^2+|\chi |^2}$$
(35)
for the singlet. Thus, for $`U<U_\mathrm{c}`$, where
$$U_\mathrm{c}=4\sqrt{\mathrm{\Gamma }_\mathrm{a}^2\mathrm{\Gamma }_\mathrm{g}^2\mathrm{sin}^2(\frac{1}{2}\delta )},$$
(36)
the Coulomb repulsion is too weak to overcome hybridization effects, and the ground state is a singlet for all $`\epsilon _0^{}`$. For $`U>U_\mathrm{c}`$, the ground state will be a doublet provided
$$U\sqrt{U^2U_\mathrm{c}^2}<\epsilon _0^{}<U+\sqrt{U^2U_\mathrm{c}^2}.$$
(37)
This allows the possibility of a level crossing as a function of $`\delta `$ if $`E_\mathrm{s}^0<E_\mathrm{d}^0<E_\mathrm{s}^0`$. However, the doublet energy is independent of $`\delta `$ in this model, hence there is no Josephson coupling in the doublet ground state.
## V Conclusions
Using a generalization of the variational wavefunctions applied in the study of the Kondo effect , we have demonstrated that the $`U=\mathrm{}`$ interacting impurity Josephson junction exhibits a first order phase transition for $`\mathrm{\Delta }T_\mathrm{K}^{}`$. When $`\mathrm{\Delta }<T_\mathrm{K}^{}`$, the ground state of the system is a singlet, and $`E(\delta )`$ is minimized at $`\delta =0`$ and maximized at $`\delta =\pi `$. When $`\mathrm{\Delta }>T_\mathrm{K}^{}`$, the ground state is a Kramers doublet, the stability of $`\delta =0`$ and $`\delta =\pi `$ is reversed, and the system forms a $`\pi `$-junction. For $`\mathrm{\Delta }T_\mathrm{K}^{}`$, the singlet and doublet energy curves may cross, in which case the ground state energy has a kink as a function of $`\delta `$ and is strongly nonsinusoidal.
We also solved for the junction’s properties within a slave boson mean field theory. Unfortunately, this approach is unable to identify a phase transition, and the singlet state always is lower in energy than the doublet. One must go beyond the static mean field approach, as has recently been accomplished in ref. , to see the transition.
## VI Bibliography |
warning/0001/math0001093.html | ar5iv | text | # 0 Introduction
## 0 Introduction
Hyperbolic complex manifolds have been studied extensively during the last 30 years (see, for example, , ). However, it is still an important problem in hyperbolic geometry to understand the algebro-geometric and the differential-geometric meaning of hyperbolicity. The use of jet bundles has become a powerful tool to attack this problem. For example, Green and Griffiths () explained an approach to establish Bloch’s Theorem on the algebraic degeneracy of holomorphic maps into abelian varieties by constructing negatively curved pseudometrics on jet bundles and by applying Ahlfors’ Lemma. Siu and Yeung () clarified this approach. Moreover, they gave a Second Main Theorem for divisors in abelian varieties, which was, very recently, clarified and generalized to the case of semi-tori by Noguchi, Winkelmann and Yamanoi ().
Demailly () presented a new construction of projective jets and pseudo-metrics on them which realizes directly the approach to Bloch’s theorem given in . These projective jets are closer to the geometry of holomorphic curves than the usual jets, since the action of the group of reparametrizations of germs of curves, which is geometrically redundant, is divided out. Using these pseudometrics on projective jets, Demailly and El Goul (, see also Mc Quillan ()) were able to show that a (very) generic surface $`X`$ in $`𝐏^3`$ of degree $`d21`$ is Kobayashi hyperbolic. As a corollary one obtains that the complement of a (very) generic curve in $`𝐏^2`$ of degree $`d21`$ is hyperbolic and hyperbolically embedded, a result first proved by Siu and Yeung () for much higher degree, using jet bundles and value distribution theory. In both papers this quasi-projective case is treated by proving hyperbolicity of a branched cover over the compactification.
However, it is desirable to have also a direct approach to deal with quasiprojective varieties, since one can hope to get easier proofs and even better results <sup>1</sup><sup>1</sup>1Actually, El Goul told the first named author that, using the results of this paper, he succeeded to drop the degree in from 21 to 15.. So one should also consider the case of logarithmic jet bundles. Noguchi () did this already for the case of the jet bundles used by Green-Griffiths. Via these bundles he generalized Bloch’s theorem to semi-abelian varieties.
The main purpose of the present paper is to generalize Demailly’s construction of projective jet bundles and strictly negatively curved pseudometrics on them to the logarithmic case. In sections 1 to 3, we establish this logarithmic generalization of Demailly’s construction explicitly via coordinates, just as Noguchi’s generalization of the jets used by Green-Griffiths. These explicit coordinates should be very useful for further applications. We also have another, more intrinsic way to obtain the same generalization in , which is much shorter, but does not give coordinates right away. In section 4 we prove the Ahlfors Lemma and the Big Picard Theorem for logarithmic projective jet bundles.
In section 5, we use our method to give a metric proof of Lang’s Conjecture for semi-abelian varieties and of a Big Picard analogue of it. The first result is due to Siu and Yeung () and Noguchi (), who used value distribution theory while we use negatively curved jet metrics. However, a common ingredient, due to Siu and Yeung (), is to construct a special jet differential (which naturally lives on the jet space constructed by Demailly) from theta functions on an abelian variety, the existence of which on a semi-abelian variety we cite from Noguchi (). Hence, the main importance of this section is the method of proof. In fact, we have to overcome some small technical difficulties to make our method work in this case: For example, we have to introduce a d-operator for sections over logarithmic projective jet bundles, and we have to deal with the case of a divisor which can have worse singularities than normal crossing, and with the precise relations between two different logarithmic structures (the one coming from the boundary divisor of a semi-abelian variety, the other coming from its union with an arbitrary reduced algebraic divisor). In this way, section 5 can also serve as a complement to sections 1 to 4.
We would like to thank J.P. Demailly and J. Noguchi for many discussions on this subject. We also would like to thank the JSPS, the SFB, the DFG, the MSRI and the universities of Göttingen, of Waterloo and of Osaka for their support during the preparation of this work. Finally we would like to thank the referee for his proposals, which led to a substantial improvement of this article.
## 1 Log-directed jet bundles
### 1.1 Logarithmic jet bundles
In this subsection we recall some basic setup and results of Noguchi in . For the proofs we refer to this article. Furthermore, in sections 1 to 3 we denote open subsets of a manifold by $`O`$, in order to distinguish them from open neighborhoods of a given point, usually with fixed coordinates centered at this point, which we denote by $`U`$.
Let $`X`$ be a complex manifold. Let $`xX`$. We consider germs $`f:(𝐂,0)(X,x)`$ of holomorphic curves through $`x`$. Two such germs are considered to be equivalent if they have the same Taylor expansions of order $`k`$ in some (and hence, any) local coordinate around $`x`$. Denote the equivalence class of $`f`$ by $`j_k(f)`$. We define $`J_kX_x=\{j_k(f)|f:(𝐂,0)(X,x)\}`$ and $`J_kX=_{xX}J_kX_x`$. Let $`\pi :J_kXX`$ be the natural projection. Then $`J_kX`$ carries the structure of a holomorphic fiber bundle over $`X`$. It is called the $`k`$-jet bundle over $`X`$. If no confusion arises, we will denote the sheaf of sections of $`J_kX`$ also by $`J_kX`$. There exist, for $`kl`$, canonical projection maps
$$\pi _{l,k}:J_kXJ_lX;j_k(f)j_l(f),$$
(1)
and $`J_1X`$ is canonically isomorphic to the holomorphic tangent bundle $`TX`$ over $`X`$. If $`F:XY`$ is a holomorphic map to another complex manifold $`Y`$, then it induces a holomorphic map
$$F_k:J_kXJ_kY;j_k(f)j_k(Ff)$$
(2)
over $`F`$.
Let $`\mathrm{\Omega }X`$ be the holomorphic cotangent bundle over $`X`$. Take a holomorphic section $`\omega H^0(O,\mathrm{\Omega }X)`$ for some open subset $`OX`$. For $`j_k(f)J_kX|_O`$ we put $`f^{}\omega =Z(t)dt`$. Then the derivatives $`\frac{d^jZ}{dt^j}(0)`$, $`0jk1`$ are well defined, independantly of the representative $`f`$ for $`j_k(f)`$. Hence, we have a well defined mapping
$$\stackrel{~}{\omega }:J_kX|_O𝐂^k;j_k(f)(\frac{d^jZ}{dt^j}(0))_{0jk1}$$
(3)
which is holomorphic. If, moreover, $`\omega ^1`$,…,$`\omega ^n`$ with $`n=dimX`$ are holomorphic 1-forms on $`O`$ such that $`\omega ^1\mathrm{}\omega ^n`$ does not vanish anywhere, then we have a biholomorphic map
$$(\stackrel{~}{\omega }^1,\mathrm{},\stackrel{~}{\omega }^n)\times \pi :J_kX|_O(𝐂^k)^n\times O$$
(4)
which we call the trivialization associated with $`\omega ^1`$,…,$`\omega ^n`$. More generally, if $`\omega `$ is a section over $`O`$ in the sheaf of meromorphic 1-forms, then the map $`\stackrel{~}{\omega }`$ defined as in equation (3) induces a meromorphic vector valued function
$$\stackrel{~}{\omega }:J_kX|_O𝐂^k.$$
(5)
Let $`\overline{X}`$ be a complex manifold with a normal crossing divisor $`D`$. This means that around any point $`x`$ of $`\overline{X}`$, there exist local coordinates $`z_1,\mathrm{},z_n`$ centered at $`x`$ such that $`D`$ is defined by $`z_1z_2\mathrm{}z_l=0`$ in some neighborhood of $`x`$ and for some $`ln`$. We note that $`l`$ depends on $`x`$, which is implicitly assumed. The pair $`(\overline{X},D)`$ will be called a log-manifold. Let $`X=\overline{X}D`$.
Following Iitaka (), we define the logarithmic cotangent sheaf $`\overline{\mathrm{\Omega }}X=\mathrm{\Omega }(\overline{X},\mathrm{log}D)`$ as the locally free subsheaf of the sheaf of meromorphic 1-forms on $`\overline{X}`$, whose restriction to $`X`$ is $`\mathrm{\Omega }X`$ (where we identify vector bundles and their sheaves of sections) and whose localization at any point $`xD`$ is given by
$$\overline{\mathrm{\Omega }}X_x=\underset{i=1}{\overset{l}{}}𝒪_{\overline{X},x}\frac{dz_i}{z_i}+\underset{j=l+1}{\overset{n}{}}𝒪_{\overline{X},x}dz_j,$$
(6)
where the local coordinates $`z_1,\mathrm{},z_n`$ around $`x`$ is chosen as before. Its dual, the logarithmic tangent sheaf $`\overline{T}X=T(\overline{X},\mathrm{log}D)`$, is a locally free subsheaf of the holomorphic tangent bundle $`T\overline{X}`$ over $`\overline{X}`$. Its restriction to $`X`$ is identical to $`TX`$, and its localization at any $`xD`$ is given by
$$\overline{T}X_x=\underset{i=1}{\overset{l}{}}𝒪_{\overline{X},x}z_i\frac{}{z_i}+\underset{j=l+1}{\overset{n}{}}𝒪_{\overline{X},x}\frac{}{z_j}.$$
(7)
Given log-manifolds $`(\overline{X^{}},D^{})`$ and $`(\overline{X},D)`$, a holomorphic map $`F:\overline{X^{}}\overline{X}`$ such that $`F^1DD^{}`$ will be called a log-morphism from $`(\overline{X^{}},D^{})`$ to $`(\overline{X},D)`$. If no confusion arises, we will simply write $`F:X^{}X`$ for the log-morphism $`F:(\overline{X^{}},D^{})(\overline{X},D)`$. It induces (see ) vector bundle morphisms,
$$F^{}:\overline{\mathrm{\Omega }}XF^1\overline{\mathrm{\Omega }}X^{}\overline{\mathrm{\Omega }}X^{}\text{and}F_{}:\overline{T}X^{}F^1\overline{T}X\overline{T}X,$$
(8)
where we have again identified locally free sheaves and vector bundles.
Let $`sH^0(O,J_k\overline{X})`$ be a holomorphic section over an open subset $`O\overline{X}`$. We say that $`s`$ is a logarithmic $`k`$-jet field if the map $`\stackrel{~}{\omega }s|_O^{}:O^{}𝐂^k`$ is holomorphic for all $`\omega H^0(O^{},\overline{\mathrm{\Omega }}X)`$ for all open subsets $`O^{}`$ of $`O`$ and where the map $`\stackrel{~}{\omega }`$ is defined as in equation (5). The set of logarithmic $`k`$-jet fields over open subsets of $`\overline{X}`$ defines a subsheaf of the sheaf $`J_k\overline{X}`$, which we denote by $`\overline{J}_kX`$. By a) of the following proposition, $`\overline{J}_kX`$ is the sheaf of sections of a holomorphic fiber bundle over $`\overline{X}`$, which we denote again by $`\overline{J}_kX`$, and which we call the logarithmic $`k`$-jet bundle of $`(\overline{X},D)`$.
###### Proposition 1.1 (see )
* $`\overline{J}_kX`$ is the sheaf of sections of a holomorphic fiber bundle over $`\overline{X}`$. (However, it is only a subsheaf and not a subbundle of $`J_k\overline{X}`$.)
* We have a canonical identification of $`(\overline{J}_kX)|_X`$ with $`J_kX`$.
* Let $`O\overline{X}`$ be an open set and $`\theta `$ be any meromorphic function on $`O`$ such that the support of its divisor $`(\theta )`$ is contained in $`D`$. Let $`d^l\mathrm{log}\theta `$ be the $`l`$-th component of the map $`\stackrel{~}{\mathrm{\Theta }}:\overline{J}_kX|_O𝐂^k`$, where $`\mathrm{\Theta }=d\mathrm{log}\theta `$ (see equation (3) and (5)). Then the differentials $`d^l\mathrm{log}\theta `$, $`l=1,\mathrm{},k`$, define holomorphic functions on $`\overline{J}_kX|_O`$. Moreover, outside the support of $`(\theta )`$, we have $`(d^l\mathrm{log}\theta )(j_k(f))=\frac{^l\mathrm{log}(\theta f)}{t^l}(0)`$.
* There exists, for $`kl`$, a canonical projection map $`\pi _{l,k}:\overline{J}_kX\overline{J}_lX`$, which extends the map $`(\pi _{l,k}|_{J_kX}):J_kXJ_lX`$ (see equation (1)), and $`\overline{J}_1X`$ is canonically isomorphic to $`\overline{T}X`$.
* A log-morphism $`F:X^{}X`$ induces a canonical map $`F_k:\overline{J}_kX^{}\overline{J}_kX`$, which extends the map $`F_k|_{J_kX}:J_kX^{}J_kX`$ (see equation (2)).
Finally, we want to express the local triviality of $`\overline{J}_kX`$ explicitly in terms of coordinates. Let $`z_1,\mathrm{},z_n`$ be coordinates in an open set $`U\overline{X}`$ in which $`D=\{z_1z_2\mathrm{}z_l=0\}`$. Let $`\omega ^1=\frac{dz_1}{z_1},\mathrm{},\omega ^l=\frac{dz_l}{z_l},\omega ^{l+1}=dz_{l+1},\mathrm{},\omega ^n=dz_n`$. Then we have a biholomorphic map (see equations (4) and (5))
$$(\stackrel{~}{\omega }^1,\mathrm{},\stackrel{~}{\omega }^n)\times \pi :\overline{J}_kX|_U(𝐂^k)^n\times U.$$
(9)
Let $`sH^0(U,\overline{J}_kX)`$ be given by $`s(x)=(Z(x);x)`$ in this trivialization with
$$Z=(Z_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k},$$
where the $`Z_j^i(x)`$ are holomorphic functions on $`U`$ and the indices $`j`$ correspond to the orders of derivatives. Then the same $`s`$, considered as an element of $`H^0(U,J_k\overline{X})`$ and trivialized by $`\omega ^1=dz_1,\mathrm{},\omega ^n=dz_n`$ (see equation (4)) is given by $`s(x)=(\widehat{Z}(x);x)`$ with $`\widehat{Z}=(\widehat{Z}_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k}`$, where
$$\widehat{Z}_j^i=\{\begin{array}{ccc}z_i(Z_j^i+g_j(Z_1^i,\mathrm{},Z_{j1}^i))& :& il\\ Z_j^i& :& il+1\end{array}.$$
(10)
Here, the $`g_j`$ are polynomials in the variables $`Z_1^i,\mathrm{},Z_{j1}^i`$ with constant coefficients and without constant terms (in particular $`g_1=0`$), which are obtained by expressing first the different components $`Z_j^i`$ of $`\stackrel{~}{(\frac{dz_i}{z_i})}s(x)`$ in terms of the components $`\widehat{Z}_j^i`$ of $`\stackrel{~}{dz_i}s(x)`$ by using the chain rule, and then by expressing the $`Z_j^i`$ in terms of the $`\widehat{Z}_j^i`$ by inverting this system of polynomial equations. This clarifies equation (1.13) in , where, for $`il`$, only the leading term $`z_iZ_j^i`$ is given. This also exhibits the sheaf inclusion $`\overline{J}_kX|_UJ_k\overline{X}|_U`$ explicitly in terms of coordinates. By abuse of notation, we also consider the $`Z_j^i`$’s as the holomorphic functions defined on $`\overline{J}_kX|_U`$ given by equation (9), so that $`Z_1^1,\mathrm{},Z_k^n;z_1,\mathrm{},z_n`$ form a holomorphic coordinate system on $`\overline{J}_kX|_U`$.
We remark that a trivialization of $`\overline{J}_kX|_U`$ is also obtained if we replace the special $`\omega `$’s used in equation (9) by any $`\omega ^1,\mathrm{},\omega ^nH^0(U,\overline{\mathrm{\Omega }}X)`$ with
$$\omega ^1\mathrm{}\omega ^n=\frac{a(x)}{z_1z_2\mathrm{}z_l}dz_1\mathrm{}dz_n,$$
(11)
where $`a(x)`$ is a nowhere vanishing holomorphic function on $`U`$.
### 1.2 Log-directed jet bundles
We first follow Demailly (). Let $`X`$ be a complex manifold together with a holomorphic subbundle $`VTX`$. The pair $`(X,V)`$ is called a directed manifold. If $`(X,V)`$ and $`(Y,W)`$ are two such manifolds, then a holomorphic map $`F:XY`$ which satisfies $`F_{}(V)W`$ is called a directed morphism.
Let $`(X,V)`$ be a directed manifold. The subset $`J_kV`$ of $`J_kX`$ is defined to be the set of $`k`$-jets $`j_k(f)J_kX`$ for which there exists a representative $`f:(𝐂,0)(X,x)`$ such that $`f^{}(t)V_{f(t)}`$ for all $`t`$ in a neighborhood of $`0`$. We will show in the next subsection that $`J_kV`$ is a fiber bundle over $`X`$, which we call the directed k-jet bundle $`J_kV`$ of $`(X,V)`$. If $`F:(X,V)(Y,W)`$ is a directed morphism, then equation (2) induces a holomorphic map
$$F_k:J_kVJ_kW;j_k(f)j_k(Ff)$$
(12)
over $`F`$, since the restriction of $`F_k:J_kXJ_kY`$ to $`J_kV`$ maps to $`J_kW`$ as $`(Ff)^{}(t)=F_{}(f^{}(t))W_{Ff(t)}`$ if $`f^{}(t)V_{f(t)}`$.
We now generalize Demailly’s directed $`k`$-jet bundles to the logarithmic context. We define a log-directed manifold to be the triple $`(\overline{X},D,\overline{V})`$, where $`(\overline{X},D)`$ is a log-manifold together with a subbundle $`\overline{V}`$ of $`\overline{T}X`$. A log-directed morphism between log-directed manifolds $`(\overline{X^{}},D^{},\overline{V}^{})`$ and $`(\overline{X},D,\overline{V})`$ is a log-morphism $`F:(\overline{X}^{},D^{})(\overline{X},D)`$ such that $`F_{}(\overline{V}^{})\overline{V}`$.
Let $`(\overline{X},D,\overline{V})`$ be a log-directed manifold and set $`V=\overline{V}|_X`$. By Proposition 1.1 we can canonically identify $`(\overline{J}_kX)|_X`$ with $`J_kX`$. Hence, the directed $`k`$-jet bundle $`J_kV`$ of $`(X,V)`$ can be considered as a subset of the logarithmic $`k`$-jet bundle $`\overline{J}_kX`$ over $`\overline{X}`$. We define the log-directed $`k`$-jet bundle $`\overline{J}_kV`$ of $`(\overline{X},D,\overline{V})`$ to be the topological closure $`\overline{J_kV}\overline{J}_kX`$ of $`J_kV`$ in $`\overline{J}_kX`$. If $`F:(\overline{X}^{},D^{},\overline{V}^{})(\overline{X},D,\overline{V})`$ is a log-directed morphism, it induces a map
$$F_k:\overline{J}_kV^{}\overline{J}_kV$$
(13)
over $`F`$ which is holomorphic. It is the restriction of the canonical map $`F_k:\overline{J}_kX\overline{J}_kX^{}`$ (see Proposition 1.1) to $`\overline{J}_kV^{}`$ and is also an extension of the map $`F_k|_X:J_kV^{}J_kV`$ (see equation (12)) to $`\overline{J}_kV^{}`$.
### 1.3 Structure of log-directed jet bundles
In this subsection, we study the local structure of $`\overline{J}_kV\overline{J}_kX`$ over $`\overline{X}`$. In particular, we show that $`\overline{J}_kV\overline{J}_kX`$ is a submanifold of $`\overline{J}_kX`$ which itself is a locally trivial bundle. This justifies the name of log-directed $`k`$-jet bundle for $`\overline{J}_kV`$ introduced in the previous subsection.
First, we consider the directed manifold $`(X,V)`$. For any point $`x_0X`$, there is a coordinate system $`(z_1,\mathrm{},z_n)`$, centered at $`x_0`$, on a neighborhood $`U`$ of $`x_0`$ such that the fibers $`V_x`$ for $`xU`$ can be defined by linear equations
$$V_x=\{\xi =\underset{1in}{}\xi _i\frac{}{z_i}|\xi _i=\underset{1mr}{}a_{im}(x)\xi _m\mathrm{for}i=r+1,\mathrm{},n\}.$$
(14)
We fix $`x_0`$, $`U`$ and these coordinates from now on. If we trivialize $`TX=J_1X`$ by $`\omega ^1=dz_1,\mathrm{},\omega ^n=dz_n`$ over $`U`$ as in equation (4), we obtain
$$V_x=\{(Z_1^1,\mathrm{},Z_1^n)|Z_1^i=\underset{1mr}{}a_{im}(x)Z_1^m\mathrm{for}i=r+1,\mathrm{},n\}.$$
(15)
If we trivialize $`J_kX`$ by the same forms, we obtain more generally:
###### Proposition 1.2
* Let $`P_h^i`$ be the polynomials in the variables $`Z_j^i`$ with coefficients depending holomorphically on $`x`$ obtained by formally differentiating the equations
$$f_i^{}(t)=\underset{1mr}{}a_{im}(f(t))f_m^{}(t)$$
$`h1`$ times with respect to $`t`$, using the fact that $`Z_j^i(j_k(f))=f_i^{(j)}(t)`$ in our trivialization. Then we have
$$(J_kV)_x=\{(Z_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k}|Z_h^i=P_h^i(x,Z_1^1,\mathrm{},Z_1^n,\mathrm{},Z_{h1}^1,\mathrm{}$$
$$\mathrm{},Z_{h1}^n,Z_h^1,\mathrm{},Z_h^r)\mathrm{for}h=1,\mathrm{},k,i=r+1,\mathrm{},n\}.$$
(16)
* $`J_kVJ_kX`$ is a submanifold, and the canonical projection
$$K:U𝐂^r;(z_1,\mathrm{},z_n)(z_1,\mathrm{},z_r)$$
induces a bundle isomorphism
$$K_k:J_kV|_UK^1(J_k𝐂^r).$$
Proof for a) Let $`j_k(f)J_kV`$. By definition, there exists a representative $`f`$ such that $`f^{}(t)V_{f(t)}`$ for all $`t`$ in a neighborhood of $`0𝐂`$, namely
$$f_i^{}(t)=\underset{1mr}{}a_{im}(f(t))f_m^{}(t).$$
Now it follows from the chain rule that $`j_k(f)`$ satisfies equations of the form $`Z_h^i=P_h^i`$, $`h=1,\mathrm{},k,i=r+1,\mathrm{},n`$, given in equation (a)).
Conversely let $`Z_j^i𝐂`$, $`i=1,\mathrm{},n,j=1,\mathrm{},k`$ be given satisfying the equations $`Z_h^i=P_h^i`$, $`h=1,\mathrm{},k,i=r+1,\mathrm{},n`$ of equation (a)). For $`xX`$ fixed, define, for $`i=1,\mathrm{},r`$, holomorphic functions
$$f_i:𝐂𝐂;tz_i(x)+\underset{\nu =1}{\overset{k}{}}\frac{Z_\nu ^i}{\nu !}t^\nu .$$
Now we integrate the system of differential equations
$$f_i^{}(t)=\underset{1mr}{}a_{im}(f_1(t),\mathrm{},f_n(t))f_m^{}(t)i=r+1,\mathrm{},n,$$
to obtain a germ $`f:(𝐂,0)(X,x)`$ with $`z_i(f)=f_i`$, $`i=1,\mathrm{},n`$. We see by construction (as $`f^{}(t)V_{f(t)}`$) that
$$(\stackrel{~}{\omega }^1,\mathrm{},\stackrel{~}{\omega }^n)(j_k(f))=(Z_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k}.$$
Proof for b) If one replaces successively in the $`P_h^i`$ all the $`Z_j^i`$ with $`ir+1`$ and $`jh1`$ by their expressions in terms of the $`Z_j^i`$ with $`ir`$ via equation (a)), we get from a) that $`J_kV|_U`$ is the graph of these new functions $`P_h^i`$, $`i=r+1,\mathrm{},n;h=1,\mathrm{},k`$ in the variables $`z_1,\mathrm{},z_n`$ and $`Z_j^i`$, $`i=1,\mathrm{},r;h=1,\mathrm{},k`$. These are in turn coordinates for $`K^1(J_k𝐂^r)`$. $``$$``$
Let now $`(\overline{X},D,\overline{V})`$ be a log-directed manifold. Let $`x_0\overline{X}`$ and let $`z_1,\mathrm{},z_n`$ be a coordinate system centered at $`x_0`$ on a neighborhood $`U`$ of $`x_0`$ where $`D`$ is defined by $`z_1z_2\mathrm{}z_l=0`$ for some $`ln`$. If we trivialize $`\overline{T}X=\overline{J}_1X`$ over $`U`$ by $`\omega ^1=\frac{dz_1}{z_1},\mathrm{},\omega ^l=\frac{dz_l}{z_l},\omega ^{l+1}=dz_{l+1},\mathrm{},\omega ^n=dz_n`$ as in equations (9) and (10), we obtain
$$V_x=\{(Z_1^1,\mathrm{},Z_1^n)|Z_1^i=\underset{mA}{}a_{im}(x)Z_1^m\mathrm{for}iB\}$$
(17)
for all $`xU`$, where, after permuting $`z_1,\mathrm{},z_l`$ respectively $`z_{l+1},\mathrm{},z_n`$, we have $`A=\{1,\mathrm{},a,l+1,\mathrm{},l+b\}`$ and $`B=\{1,\mathrm{},n\}A`$ with $`a+b=r=\mathrm{rank}V`$. We fix this setup for the rest of this section.
First, we generalize the projection $`K`$ to log-directed manifolds:
###### Proposition 1.3
With $`E=\{z_1\mathrm{}z_a=0\}`$, the log-directed projection
$$K:(\overline{X},D,\overline{V})|_U(𝐂^r,E,\overline{T}𝐂^r);(z_1,\mathrm{},z_n)(z_1,\mathrm{}z_a,z_{l+1},\mathrm{},z_{l+b})$$
has bijective differential map $`(K_{})_x`$ for all $`xU`$.
Proof We trivialize $`\overline{T}𝐂^r`$ by the forms $`\omega ^1=\frac{dz_1}{z_1},\mathrm{},\omega ^a=\frac{dz_a}{z_a},\omega ^{l+1}=dz_{l+1},\mathrm{},\omega ^{l+b}=dz_{l+b}`$. We claim that $`K_{}`$ is given by the projection map
$$(K_{})_x:(\overline{T}X)_x(\overline{T}𝐂^r)_x;(Z_1^1,\mathrm{},Z_1^n)(Z_1^1,\mathrm{},Z_1^a,Z_1^l,\mathrm{},Z_1^{l+b})$$
(18)
in these coordinates. In fact, by analytic continuation it suffices to prove equation (18) for $`xX=\overline{X}D`$. Let $`(Z_1^1,\mathrm{},Z_1^n)(\overline{T}X)_x=(TX)_x`$ be a vector in the logarithmic coordinate system. If we retrivialize $`(TX)_x`$ respectively $`(T𝐂^r)_{K(x)}`$ with the forms $`dz_i`$ ($`i=1,\mathrm{},n`$ respectively $`iA`$) instead, then the given vector is expressed by $`(z_1Z_1^1,\mathrm{},z_lZ_1^l,Z_1^{l+1},\mathrm{},Z_1^n)`$ (see equations (9) and (10)). Furthermore, in the latter trivialization, the map $`(K_{})_x`$ is just the projection to the components given by $`A`$. So equation (18) follows. Hence, the assertion follows from equation (17). $``$$``$
If we trivialize $`\overline{J}_kX`$ by the same forms as in equation (17), we obtain the following extension of Proposition 1.2.
###### Proposition 1.4
Let the setup be as above.
* There are polynomials $`Q_h^i`$ in the variables $`Z_j^i`$ with coefficients which are holomorphic functions on $`U`$ such that
$$\overline{J}_kV_x=\{(Z_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k}|Z_h^i=Q_h^i(x,Z_1^1,\mathrm{},Z_1^n,\mathrm{},Z_{h1}^1,\mathrm{}$$
$$\mathrm{},Z_{h1}^n,Z_h^1,\mathrm{},Z_h^r)\mathrm{for}h=1,\mathrm{},k,iB\}.$$
(19)
* $`\overline{J}_kV\overline{J}_kX`$ is a submanifold and the projection map $`K`$ defined in Proposition 1.3 induces a bundle isomorphism
$$K_k:\overline{J}_kV|_UK^1(\overline{J}_k(𝐂^rE)).$$
Proof for a) By using the coordinate system $`(z_1,\mathrm{},z_n)`$ on $`U`$, the map
$$\mathrm{\Psi }:𝐂^n𝐂^n;(w_1,\mathrm{},w_n)(e^{w_1},\mathrm{},e^{w_l},w_{l+1},\mathrm{},w_n)=(z_1,\mathrm{},z_n)$$
induces a locally biholomorphic map $`\mathrm{\Psi }:\mathrm{\Psi }^1(U)UD`$. Let $`\widehat{V}=\mathrm{\Psi }_{}^1(V)T(\mathrm{\Psi }^1(U))`$ and let $`W_j^i`$, $`i=1,\mathrm{},n;j=1,\mathrm{},k`$ be the components of the first part of the trivialization map
$$(\stackrel{~}{dw_1},\mathrm{},\stackrel{~}{dw_n})\times \pi :J_k(\mathrm{\Psi }^1(U))(𝐂^k)^n\times \mathrm{\Psi }^1(U).$$
###### Lemma 1.5
On $`J_k(\mathrm{\Psi }^1(U))`$, we have
$$W_j^i=Z_j^i\mathrm{\Psi }_ki=1,\mathrm{},n;j=1,\mathrm{},k.$$
Proof of Lemma 1.5 Let $`j_k(f)J_k(\mathrm{\Psi }^1(U))`$ and let $`f=(f_1,\mathrm{},f_n):(𝐂,0)\mathrm{\Psi }^1(U)`$ represent it. We put $`f^{}dw_i=df_i(t)=C_i(t)dt`$. Then we have $`W_j^i(j_k(f))=\frac{^{j1}C_i(t)}{t^{j1}}|_{t=0}`$. On the other hand,
$$(\mathrm{\Psi }f)^{}\omega ^i=df_i(t)=C_i(t)dt$$
independently of $`i`$, and hence,
$$Z_j^i\mathrm{\Psi }_k(j_k(f))=\frac{^{j1}C_i(t)}{t^{j1}}|_{t=0}=W_j^i(j_k(f)).$$
$``$$``$
Since $`\widehat{V}T(\mathrm{\Psi }^1(U))`$ is the inverse image of $`VT(UD)`$, we have
$$\widehat{V}_w=\{(W_1^1,\mathrm{},W_1^n)|W_1^i=\underset{mA}{}a_{im}\mathrm{\Psi }(w)W_1^m\mathrm{for}iB\},$$
(20)
for $`w\mathrm{\Psi }^1(U)`$. Using Proposition 1.2, we get that
$$(J_k\widehat{V})_w=\{(W_j^i)_{i=1,\mathrm{},n;j=1,\mathrm{},k}|W_h^i=P_h^i(w,W_1^1,\mathrm{},W_1^n,\mathrm{},W_{h1}^1,\mathrm{}$$
$$\mathrm{},W_{h1}^n,W_h^1,\mathrm{},W_h^r)\mathrm{for}h=1,\mathrm{},k,i=r+1,\mathrm{},n\},$$
where the $`P_h^i`$ are the polynomials in the variables $`W_j^i`$ with coefficients depending holomorphically on $`w\mathrm{\Psi }^1(U)`$ obtained by formally differentiating
$$f_i^{}(t)=\underset{1mr}{}a_{im}\mathrm{\Psi }(f(t))f_m^{}(t)$$
(21)
$`h1`$ times. The important point is now that the coefficient functions factor through $`\mathrm{\Psi }`$ by holomorphic functions which are still holomorphic for all $`xU`$:
###### Main Lemma 1.6
The coefficients of the polynomials $`P_h^i`$ factor through $`\mathrm{\Psi }`$ by holomorphic functions which are defined on all of $`U`$. Namely,
$$P_h^i(w,W_1^1,\mathrm{},W_1^n,\mathrm{},W_{h1}^1,\mathrm{},W_{h1}^n,W_h^1,\mathrm{},W_h^r)$$
$$=Q_h^i(x,W_1^1,\mathrm{},W_1^n,\mathrm{},W_{h1}^1,\mathrm{},W_{h1}^n,W_h^1,\mathrm{},W_h^r)$$
for $`x=\mathrm{\Psi }(w)UD`$, where the $`Q_h^i`$ are polynomials in the variables $`W_j^i`$, with coefficients which are holomorphic in $`x`$ on all of $`U`$.
Proof of the Main Lemma If $`\alpha :U𝐂`$ is holomorphic, we have:
$$\frac{}{t}\alpha \mathrm{\Psi }(f(t))=\underset{\mu =1}{\overset{n}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))\frac{}{t}(\mathrm{\Psi }_\mu f)(t)$$
$$=\underset{\mu =1}{\overset{l}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))\frac{}{t}e^{f_\mu }(t)+\underset{\mu =l+1}{\overset{n}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))\frac{}{t}f_\mu (t)$$
$$=\underset{\mu =1}{\overset{l}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))(z_\mu \mathrm{\Psi })(f(t))f_\mu ^{}(t)+\underset{\mu =l+1}{\overset{n}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))f_\mu ^{}(t)$$
$$=\underset{\mu =1}{\overset{l}{}}((z_\mu \frac{\alpha }{z_\mu })\mathrm{\Psi })(f(t))f_\mu ^{}(t)+\underset{\mu =l+1}{\overset{n}{}}(\frac{\alpha }{z_\mu }\mathrm{\Psi })(f(t))f_\mu ^{}(t).$$
Now the assertion follows by induction on $`h`$ as the coefficients of the polynomials $`P_h^i`$ are obtained by formally differentiating the equation in equation (21) $`h1`$ times and that the functions $`a_{im}`$ are holomorphic on all of $`U`$. $``$$``$
Finally we patch together these results and obtain the proof of Proposition 1.4 (a): Using Lemma 1.5, the Main Lemma and the local isomorphisms $`\mathrm{\Psi }^1`$ we see that this assertion holds for all $`xUD`$, with equations $`Z_h^i=Q_h^i(\mathrm{\Psi }(\mathrm{\Psi }^1(x)),\mathrm{}W_j^i\mathrm{})=Q_h^i(x,\mathrm{}Z_j^i\mathrm{})`$ which are independent of the choice of the local isomorphism $`\mathrm{\Psi }^1`$. Moreover, their coefficient functions are still holomorphic on $`U`$. Since $`\overline{J}_kV`$ is defined as the closure of $`J_kV`$ in $`\overline{J}_kX`$, the structure of the equations $`Z_h^i=Q_h^i`$ implies the assertion for all $`xU`$.
Proof for Proposition 1.4 b) It is verbatim that of Proposition 1.2 b). $``$$``$
### 1.4 Regular jets
Let $`(X,V)`$ be a directed manifold. The subset $`J_kV^{\mathrm{sing}}J_kV`$ of singular k-jets is defined to be the subset of $`k`$-jets $`j_k(f)J_kV`$ of germs $`f:(𝐂,0)(X,x)`$ such that $`f^{}(0)=0`$. Its complement $`J_kV^{\mathrm{reg}}=J_kVJ_kV^{\mathrm{sing}}`$ defines the regular k-jets.
Let now $`(\overline{X},D,\overline{V})`$ be a log-directed manifold. Define $`\overline{J}_kV^{\mathrm{sing}}\overline{J}_kV`$ to be the closure $`\overline{J_kV^{\mathrm{sing}}}\overline{J}_kV`$ of $`J_kV^{\mathrm{sing}}`$ in $`\overline{J}_kV`$ and set $`\overline{J}_kV^{\mathrm{reg}}=\overline{J}_kV\overline{J}_kV^{\mathrm{sing}}`$.
###### Proposition 1.7
* $`\overline{J}_kV^{\mathrm{sing}}\overline{J}_kV`$ is a smooth submanifold of codimension $`r=\mathrm{rank}\overline{V}`$. In terms of the local coordinates of $`\overline{J}_kV\overline{J}_kX`$ (see Proposition 1.4), this submanifold is given by the equations
$$Z_1^i=0,iA.$$
* The bundle isomorphism $`K_k`$ given in Proposition 1.4 b) respects the singular and regular jets.
Proof for a) Using equations (9), (10) and (17), we see that $`J_kV^{\mathrm{sing}}`$ in $`\overline{J}_kV`$ is given locally by the equations $`Z_1^i=0,iA`$. So the assertion follows.
Proof for b) This follows directly from the proof of Proposition 1.4 b). $``$$``$
## 2 Log-Demailly-Semple jet bundles
A natural notion of higher order contact structures was introduced on a firm setting by Demailly () in the holomorphic category for the study of complex hyperbolic geometry. These structures realize natural “quotient” spaces of directed jet bundles. Demailly called them Semple Jet bundles after Semple (), who constructed and worked with these bundles over $`𝐏^2`$. In this section, we generalize these bundles to the logarithmic case and prove some important properties, like functoriality and local triviality. Their connection with log-directed jet bundles will be discussed in section 3.
### 2.1 Definition of log-Demailly-Semple jet bundles
We begin with a log-directed manifold $`X_0=(\overline{X}_0,D_0,\overline{V}_0)`$. We inductively define $`(\overline{X}_k,D_k,\overline{V}_k)`$ as follows. Let $`\overline{X}_k=𝐏(\overline{V}_{k1})`$ with its natural projection $`\pi _k`$ to $`\overline{X}_{k1}`$. Set $`D_k=\pi _k^1(D_{k1})`$ and $`X_k=\overline{X}_kD_k`$. Let $`𝒪_{\overline{X}_k}(1)`$ be the tautological subbundle of $`\pi _k^1\overline{V}_{k1}\pi _k^1\overline{T}X_{k1}`$, and set
$$\overline{V}_k=(\pi _k)_{}^1\left(𝒪_{\overline{X}_k}(1)\right).$$
(1)
Equivalently, $`\overline{V}_k\overline{T}X_k`$ is defined, for every point $`(x,[v])\overline{X}_k=𝐏(\overline{V}_{k1})`$ associated with a vector $`v(\overline{V}_{k1})_x`$ for $`x\overline{X}_{k1}`$, by
$$(\overline{V}_k)_{(x,[v])}=\{\xi (\overline{T}X_k)_{(x,[v])}:(\pi _k)_{}\xi 𝐂v\},𝐂v(\overline{V}_{k1})_x(\overline{T}X_{k1})_x.$$
Since $`(\pi _k)_{}:\overline{T}X_k\pi _k^1\overline{T}X_{k1}`$ has maximal rank everywhere as it is a bundle projection, we see that $`\overline{V}_k`$ is a subbundle of $`\overline{T}X_k`$ giving a log-directed structure for $`X_k`$ and also for $`\pi _k`$, thus completing our inductive definition.
We set $`\overline{P}_kV=\overline{X}_k`$, $`P_kV=X_k`$, $`\overline{P}_kX=\overline{P}_kTX`$ and $`P_kX=P_kTX`$. Let
$$\pi _{j,k}=\pi _{j+1}\mathrm{}\pi _{k1}\pi _k:\overline{P}_kV\overline{P}_jV$$
for $`j<k`$. We also put $`(\overline{P}_kV)_x=(\pi _{0,k})^1(x)`$ and $`(\overline{V}_k)_x=\overline{V}_k|_{(\overline{P}_kV)_x}`$ for $`x\overline{X}`$.
Note that $`\mathrm{ker}(\pi _k)_{}=T_{\overline{P}_kV/\overline{P}_{k1}V}`$ by definition. This gives the following short exact sequence of vector bundles over $`\overline{P}_kV`$:
$$0T_{\overline{P}_kV/\overline{P}_{k1}V}\overline{V}_k\stackrel{(\pi _k)_{}}{}𝒪_{\overline{P}_kV}(1)0.$$
(2)
Furthermore, we have the Euler exact sequence for projectivized bundles (applied to the bundle $`𝐏(\overline{V}_{k1})\overline{P}_{k1}V=\overline{X}_{k1}`$)
$$0𝒪_{\overline{P}_kV}\pi _k^1\overline{V}_{k1}𝒪_{\overline{P}_kV}(1)T_{\overline{P}_kV/\overline{P}_{k1}V}0.$$
(3)
The composition of vector bundle morphisms over $`\overline{P}_kV`$
$$𝒪_{\overline{P}_kV}(1)\pi _k^1\overline{V}_{k1}\stackrel{(\pi _k)^1(\pi _{k1})_{}}{}\pi _k^1𝒪_{\overline{P}_{k1}V}(1)$$
yields an effective divisor $`\mathrm{\Gamma }_k`$ corresponding to a section of
$$𝒪_{\overline{P}_kV}(1)\pi _k^1𝒪_{\overline{P}_{k1}V}(1)=𝒪(\mathrm{\Gamma }_k).$$
(4)
There is a canonical divisor on $`\overline{P}_kV`$ given by
$$\overline{P}_kV^{\mathrm{sing}}=\underset{2jk}{}\pi _{j,k}^1(\mathrm{\Gamma }_j)\overline{P}_kV.$$
Finally, we set
$$\overline{P}_kV^{\mathrm{reg}}=\overline{P}_kV\overline{P}_kV^{\mathrm{sing}}\text{ and }𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}=\left(𝒪_{\overline{P}_kV}(1)|_{\overline{P}_kV^{\mathrm{reg}}}\right)\overline{P}_kV,$$
where the last $`\overline{P}_kV`$ denotes the zero section.
### 2.2 Properties of log-Demailly-Semple jet bundles
###### Proposition 2.1
Let $`F:(\overline{X}^{},D^{},\overline{V}^{})(\overline{X},D,\overline{V})`$ be a log-directed morphism.
* For all $`k0`$ there exist log-directed meromorphic maps (log-directed morphisms outside the locus of indeterminacy)
$$F_k:(\overline{P}_kV^{},D_k^{},\overline{V}_k^{})\mathrm{}(\overline{P}_kV,D_k,\overline{V}_k)$$
which commute with the projections, more specifically for all $`0lk1`$ one has
$$\pi _{l,k}F_k=F_l\pi _{l,k}^{}.$$
These maps in turn induce meromorphic maps
$$(F_k)_{}:𝒪_{\overline{P}_{k+1}V^{}}(1)\mathrm{}𝒪_{\overline{P}_{k+1}V}(1)$$
(holomorphic where $`F_k`$ is) which also commute with the projections.
* If the differential map $`F_{}:\overline{V}^{}F^1(\overline{V})`$ is injective over a point $`x_0\overline{X}^{}`$, then there exists a neighborhood $`U`$ of $`x_0`$ in $`\overline{X}^{}`$ over which the maps $`F_k`$ are log-directed morphisms and the induced maps
$$F_k:\overline{P}_kV^{}F^1(\overline{P}_kV)$$
are holomorphic embeddings and the induced maps between line bundles
$$(F_k)_{}:𝒪_{\overline{P}_{k+1}V^{}}(1)F^1(𝒪_{\overline{P}_{k+1}V}(1))$$
over these embeddings are injective.
* If the differential map $`F_{}:\overline{V}^{}F^1(\overline{V})`$ is bijective in a point $`x_0\overline{X}^{}`$, then there exists a neighborhood $`U`$ of $`x_0`$ in $`\overline{X}^{}`$ over which the maps $`F_k`$ are log-directed morphisms and the induced maps
$$F_k:\overline{P}_kV^{}F^1(\overline{P}_kV)$$
$$(F_k)_{}:𝒪_{\overline{P}_{k+1}V^{}}(1)F^1(𝒪_{\overline{P}_{k+1}V}(1))$$
are all bundle isomorphisms over $`U`$.
Combining Proposition 2.1 with Proposition 1.3 yields the following, which we will use in the next section to study $`\overline{P}_kV`$ by studying $`K^1(\overline{P}(𝐂^rE))`$:
###### Proposition 2.2
Let $`(\overline{X},D,\overline{V})`$ be a log-directed manifold. There exists a neighborhood $`U`$ of $`x_0`$ in $`\overline{X}`$ and a log-directed projection
$$K:(\overline{X},D,\overline{V})|_U(𝐂^r,E,\overline{T}𝐂^r);(z_1,\mathrm{},z_n)(z_1,\mathrm{}z_a,z_{l+1},\mathrm{},z_{l+b}),$$
with $`E=\{z_1\mathrm{}z_a=0\}`$ and $`a+b=r=\mathrm{rank}V`$, which induces
$$(K_k)_{}:𝒪_{\overline{P}_{k+1}V}(1)|_UK^1(𝒪_{\overline{P}_{k+1}(𝐂^rE)}(1)),$$
$$K_k:\overline{P}_kV|_UK^1(\overline{P}_k(𝐂^rE))$$
as bundle isomorphisms. $`\mathrm{}`$
Proof for a) We proceed by induction on $`k`$. The case $`k=0`$ holds by assumption. Assume the case $`k1`$ holds. If $`(F_{k1})_{}:\overline{V}_{k1}^{}\overline{V}_{k1}`$, define
$$F_k:=𝐏((F_{k1})_{}):\overline{P}_kV^{}\overline{P}_kV.$$
Then by definition of $`D_k^{}=(\pi _k^{})^1D_{k1}^{}`$ and $`D_k=\pi _k^1D_{k1}`$ the map $`F_k`$ is a log-meromorphic morphism which commutes with projections. Let $`(F_k)_{}:\overline{T}P_kV^{}\overline{T}P_kV`$ be the log-differential map defined as in equation (8). If $`\xi (\overline{V}_k^{})_{(w,[v])}`$ with $`w\overline{P}_{k1}V^{}`$ and $`v(\overline{V}_{k1}^{})_w`$, then
$$(\pi _k)_{}((F_k)_{}\xi )=(\pi _kF_k)_{}\xi =(F_{k1}\pi _{k1}^{})_{}\xi $$
$$=(F_{k1})_{}(\pi _{k1}^{})_{}\xi (F_{k1})_{}𝐂v=𝐂((F_{k1})_{}v),$$
hence, $`(F_k)_{}\xi (\overline{V}_k)_{(F_{k1}(w),[(F_{k1})_{}v])}`$, so $`F_k`$ is a log-directed meromorphic morphism. The second part of the assertion is clear.
Proof for b and c) $`(F_{})_x:(\overline{V}^{})_x(\overline{V})_{F(x)}`$ remains injective (respectively bijective) for all $`x`$ in a neighborhood $`U(x_0)`$. By the bundle structures it suffices to prove that for all $`xU(x_0)`$, the maps $`F_k:(\overline{P}_kV^{})_x(\overline{P}_kV)_{F(x)}`$ are holomorphic embeddings (respectively biholomorphic maps) and the maps $`(F_k)_{}:(V_k^{})_x(V_k)_{F(x)}`$ are injective (respectively bijective) bundle maps over them. We prove this by induction on $`k`$. The case $`k=0`$ holds by assumption. Assume the case $`k1`$ holds. By projectivizing the injective (respectively bijective) bundle map $`(F_{k1})_{}`$ and by a), we get that $`F_k:(\overline{P}_kV^{})_x(\overline{P}_kV)_{F(x)}`$ is an injective (respectively bijective) log-directed morphism. Furthermore, since $`F_{k1}:(\overline{P}_{k1}V^{})_x(\overline{P}_{k1}V)_{F(x)}`$ is a holomorphic embedding (resp biholomorphic) and since $`(F_{k1})_{}\overline{V}_k^{}\overline{V}_k|_{F_{k1}(\overline{P}_{k1}V^{})}`$ is a holomorphic subbundle (respectively the same bundle) the map $`F_k`$ is an embedding (respectively biholomorphic). It remains to show that $`(F_k)_{}`$ is again injective (respectively bijective). Let $`\xi (\overline{V}_k^{})_w`$ for $`w\overline{P}_kV^{}`$ such that $`(F_k)_{}\xi =0`$. Then
$$0=(\pi _k)_{}(F_k)_{}\xi =(F_{k1})_{}(\pi _k^{})_{}\xi .$$
Since the map $`(F_{k1})_{}`$ is injective we get $`\xi \mathrm{ker}(\pi _k)_{}`$. But the subbundle $`\mathrm{ker}(\pi _k)_{}\overline{V}_k^{}\overline{T}P_kV^{}`$ can be canonically identified with the relative tangent bundle $`T_{\overline{P}_kV^{}/\overline{P}_{k1}V^{}}`$, which is a subbundle of $`T\overline{P}_KV^{}`$. Since we have shown that $`F_k`$ is a holomorphic embedding, $`(F_k)_{}`$ is injective on $`(T\overline{P}_kV^{})_w`$, which contains $`\xi `$. As $`(F_k)_{}\xi =0`$ this forces $`\xi =0`$. So $`(F_k)_{}`$ is injective on $`(\overline{V}_k^{})_w`$. Moreover, if the assumption in c) holds, then $`\mathrm{rank}V_k^{}=\mathrm{rank}V^{}=\mathrm{rank}V=\mathrm{rank}V_k`$, and so $`(F_k)_{}`$ is bijective. $``$$``$
###### Corollary 2.3
We have $`F_k((\overline{P}_kV^{})^{\mathrm{sing}})\overline{P}_kV^{\mathrm{sing}}.`$ Moreover, if $`F_{}:V^{}F^1(V)`$ is injective at a point $`x_0\overline{X}^{}`$, then there is a neighborhood of $`x_0`$ over which $`(\overline{P}_kV^{})^{\mathrm{sing}}`$ is isomorphic to $`F_k^1(\overline{P}_kV^{\mathrm{sing}})`$.
Proof By the definition of the singular locus and Proposition 2.1 a) it suffices to show $`F_j(\mathrm{\Gamma }_j^{})\mathrm{\Gamma }_j`$ (respectively $`F_j^1(\mathrm{\Gamma }_j)=\mathrm{\Gamma }_j^{}`$) for $`2jk`$. Moreover, since the $`\mathrm{\Gamma }_j^{}`$ and $`\mathrm{\Gamma }_j`$ are divisors without vertical components, it suffices to prove the assertions where all maps $`F_i`$, $`ij`$ are holomorphic. The first assertion follows immediately from the definitions of $`\mathrm{\Gamma }_j^{}`$ and $`\mathrm{\Gamma }_j`$ and the equation
$$(\pi _{l1})_{}(F_{l1})_{}=(F_{l2})_{}(\pi _{l1}^{})_{}.$$
(5)
The second follows from this equation and the injectivity of $`(F_{l2})_{}`$. $``$$``$
###### Corollary 2.4
Let $`(\overline{X},D,\overline{V})`$ be a log-directed manifold. If $`\overline{V}\overline{W}\overline{T}X`$ are holomorphic subbundles, then we have natural inclusions of submanifolds
$$\overline{P}_kV\overline{P}_kW\overline{P}_kX$$
and the associated maps over these inclusions of the line bundles
$$𝒪_{\overline{P}_kV}(1)𝒪_{\overline{P}_kW}(1)𝒪_{\overline{P}_kX}(1)$$
are line bundle restrictions.
Proof We apply Proposition 2.1 to the log-directed morphism $`i:(\overline{X},D,\overline{V})(\overline{X},D,\overline{W})`$, where $`i:\overline{X}\overline{X}`$ is the identity map and induces the bundle inclusion $`i_{}:\overline{V}\overline{W}`$. By Proposition 2.1 a) and b) we obtain a log-directed morphism $`i_k:\overline{P}_kV\overline{P}_kW`$ which locally over $`\overline{X}`$ is, moreover, a holomorphic embedding $`i_k:\overline{P}_kVi^1(\overline{P}_kW)=\overline{P}_kW`$. Hence, $`i_k:\overline{P}_kV\overline{P}_kW`$ is a holomorphic embedding. The other statements follow in a similar way. $``$$``$
## 3 Log-directed jet differentials
### 3.1 Demailly-Semple jet bundles and jet differentials
In this subsection we recall parts of some basic results of the work of Demailly on his construction of the Demailly-Semple jet bundles, which we generalize to logarithmic setting in the next subsections.
Let $`(X,V)`$ be a directed manifold. Without loss of generality we assume $`r=\mathrm{rank}V2`$ in section 3, for the situation is trivial otherwise. Let
$$G_k=J_k𝐂_0^{\mathrm{reg}}=\{t\varphi (t)=\underset{i=1}{\overset{k}{}}a_it^i,a_1𝐂^{},a_i𝐂,i2\}$$
be the group of reparametrizations. Elements $`\varphi G_k`$ act on $`J_kV`$ as holomorphic automorphisms by
$$\varphi :J_kVJ_kV;j_k(f)j_k(f\varphi ).$$
In particular, $`𝐂^{}`$ acts on $`J_kV`$.
Every nonconstant germ $`f:(𝐂,0)X`$ tangent to $`V`$ lifts to a unique germ $`f_{[k]}:(𝐂,0)P_kV`$ tangent to $`V_k`$. $`f_{[k]}`$ can be defined inductively to be the projectivization of $`f_{[k1]}^{}:(𝐂,0)V_{k1}`$. As such we also have a germ
$$f_{[k1]}^{}:(𝐂,0)𝒪_{P_kV}(1).$$
This construction is actually a special case of our construction in Proposition 2.1, since $`P_k𝐂=𝐂`$. We get, moreover:
###### Proposition 3.1
Let $`F:(X^{},V^{})(X,V)`$ be a directed morphism. Let $`f:(𝐂,0)X^{}`$ be a germ tangent to $`V^{}`$ such that the germ $`Ff:(𝐂,0)X`$ is nonconstant. Then there exists a neighborhood $`U`$ of $`0𝐂`$ such that for all $`tU`$, $`t0`$, and for all $`k0`$, the map $`F_k:P_kV^{}P_kV`$ (see Proposition 2.1) is a morphism around $`f_{[k]}(t)`$, and we have on $`U`$:
$$(Ff)_{[k]}=F_k(f_{[k]}).$$
(1)
Proof Since the germ $`Ff:(𝐂,0)X`$ is nonconstant, we can find a neighborhood $`U`$ of $`0𝐂`$ such that $`(Ff)^{}(t)0`$ for all $`tU`$, $`t0`$. From the equation $`(Ff)^{}(t)=(\pi _{0,k})_{}(Ff)_{[k]}^{}(t)`$, we get that $`(Ff)_{[k]}^{}(t)0`$ for all $`k0`$. We now proceed by induction on $`k`$. The case $`k=0`$ is trivial. Assume the case $`k1`$. This means that for all $`tU`$, $`t0`$, the map $`F_{k1}:P_{k1}V^{}P_{k1}V`$ is a morphism at $`f_{[k1]}(t)`$, and we have on $`U`$:
$$(Ff)_{[k1]}=F_{k1}(f_{[k1]}).$$
Taking the derivative, we obtain
$$(Ff)_{[k1]}^{}(t)=(F_{k1})_{}(f_{[k1]})^{}(t).$$
Now the left hand side is nonzero for $`t0`$, so the right hand side is nonzero, too, and we just can projectivize and obtain the assertion for $`t0`$. Finally equation (1) still holds for $`t=0`$ by analytic continuation. $``$$``$
The bundle of directed invariant jet differentials of order $`k`$ and degree $`m`$, denoted by $`E_{k,m}V^{}`$, is defined as follows: Its sheaf of sections $`𝒪(E_{k,m}V^{})`$ over $`X`$ consists of holomorphic functions on $`J_kV|_O`$ which satisfy
$$Q(j_k(f\varphi ))=\varphi ^{}(0)^mQ(j_k(f))j_k(f)J_kV^{\mathrm{reg}}|_O\text{ and }\varphi G_k$$
(2)
as $`O`$ varies over the open subsets of $`X`$. We remark that equation (2) implies that the functions $`Q`$, restricted to the fibers of $`J_kV`$, are polynomials of weighted degree $`m`$ with respect to the $`𝐂^{}`$-action, so that our definition coincides with the usual one.
###### Theorem 3.2 (Demailly ())
Let $`(X,V)`$ be a directed manifold.
* The maps
$$\stackrel{~}{\alpha }_k:J_kV^{\mathrm{reg}}𝒪_{P_kV}(1)^{\mathrm{reg}},j_k(f)f_{[k1]}^{}(0),$$
$$\alpha _k:J_kV^{\mathrm{reg}}P_kV^{\mathrm{reg}},j_k(f)f_{[k]}(0)$$
are well defined, holomorphic and surjective.
* If $`\varphi G_k`$ is a reparametrization, one has
$$(f\varphi )_{[k1]}^{}(0)=f_{[k1]}^{}(0)\varphi ^{}(0),$$
$$(f\varphi )_{[k]}(0)=f_{[k]}(0).$$
* The quotient $`J_kV^{\mathrm{reg}}/G_k`$ of $`J_kV^{\mathrm{reg}}`$ by $`G_k`$ has the structure of a locally trivial fiber bundle over $`X`$, and the map
$$\alpha _k/G_k:J_kV^{\mathrm{reg}}/G_kP_kV$$
is a holomorphic embedding which identifies $`J_kV^{\mathrm{reg}}/G_k`$ with $`P_kV^{\mathrm{reg}}`$.
* The direct image sheaf
$$(\pi _{0,k})_{}𝒪_{P_kV}(m)𝒪(E_{k,m}V^{})$$
can be identified with the sheaf $`𝒪(E_{k,m}V^{})`$.
###### Corollary 3.3
* The group $`G_k^o=\{t\varphi (t)=_{i=1}^ka_it^i,a_1=1\}`$ acts transitively on the fibers of $`\stackrel{~}{\alpha }_k`$.
* The maps $`\stackrel{~}{\alpha }_k`$ and $`\alpha _k`$ are holomorphic submersions.
Proof for 1) By Theorem 3.2 b) and c), the group $`G_k`$ acts transitively on the fibers of $`\alpha _k`$. So for any two points $`p`$ and $`q`$ of a fixed fiber of $`\stackrel{~}{\alpha }_k`$ we find $`\varphi G_k`$ such that $`\varphi (p)=q`$. Again by b) we have $`\varphi ^{}(0)=1`$, so $`\varphi G_k^o`$.
Proof for 2) It suffices to prove the statement for $`\alpha _k`$, since by the action of $`𝐂^{}=G_k/G_k^o`$ on $`J_kV^{\mathrm{reg}}`$ and Theorem 3.2 b), it also follows for $`\stackrel{~}{\alpha }_k`$. The assertion is equivalent with the existence of local sections for $`\alpha _k`$ through every point $`j_k(f)J_kV^{\mathrm{reg}}`$.
Let $`j_k(f)J_kV^{\mathrm{reg}}`$ be given, and let $`w_0=\alpha (j_k(f))`$. Since $`f^{}(0)0`$, we get that $`f_{[k1]}^{}(0)0`$. Then by Corollary 5.12 in Demailly’s paper and its proof we find a neighborhood $`U(w_0)P_kV^{\mathrm{reg}}`$ and a holomorphic family of germs $`(f_w)`$, $`wU(w_0)`$, such that $`(f_w)_{[k]}(0)=w`$ and $`f_{w_0}=f`$. After possibly shrinking $`U(w_0)`$, we may assume that $`f_w^{}(0)0`$ for all $`wU(w_0)`$. Thus $`wj_k(f_w)`$ defines a local holomorphic section $`s:U(w_0)J_kV^{\mathrm{reg}};wj_k(f_w)`$ with $`s(\alpha _k(j_k(f)))=s(w_0)=j_k(f_{w_0})=j_k(f)`$. $``$$``$
### 3.2 Local trivializations
###### Proposition 3.4
Let $`z_1,\mathrm{},z_r`$ be the standard coordinates of $`𝐂^r`$, let $`ar`$, let $`E=\{z_1\mathrm{}z_a=0\}`$ and $`P=(\stackrel{a}{\stackrel{}{1,\mathrm{},1}},\stackrel{ra}{\stackrel{}{0,\mathrm{},0}})𝐂^r`$.
* The trivialization of $`\overline{J}_k(𝐂^rE)`$ by the forms $`\omega ^1=\frac{dz_1}{z_1},\mathrm{},\omega ^a=\frac{dz_a}{z_a},\omega ^{a+1}=dz_{a+1},\mathrm{},`$ $`\omega ^r=dz_r`$, induce an isomorphism
$$\overline{J}_k(𝐂^rE)J_k(𝐂^r)_P\times 𝐂^r$$
(3)
which respects regular and singular jets and commutes (outside $`E`$) with $`G_k`$.
* For the log-manifold $`(𝐂^r,E)`$ there exists a line bundle isomorphism
$$𝒪_{\overline{P}_k(𝐂^rE)}(1)𝒪_{P_k𝐂^r}(1)_P\times 𝐂^r$$
(4)
which respects regular and singular jets and such that the diagram
$$\begin{array}{ccc}\overline{J}_k(𝐂^rE)^{\mathrm{reg}}|_{𝐂^rE}& & J_k(𝐂^r)_P^{\mathrm{reg}}\times (𝐂^rE)\\ & & \\ \stackrel{~}{\alpha }_k& & (\stackrel{~}{\alpha }_k)_P\times id\\ & & \\ 𝒪_{\overline{P}_k(𝐂^rE)}(1)^{\mathrm{reg}}|_{𝐂^rE}& & 𝒪_{P_k(𝐂^r)}(1)_P^{\mathrm{reg}}\times (𝐂^rE)\end{array}$$
(5)
commutes.
Proof for a) The composition
$$(J_k𝐂^r)_P=J_k(𝐂^rE)_P\overline{J}_k(𝐂^rE)(𝐂^k)^r$$
is an isomorphism, where the last morphism is given by the first factor of the trivialization map in equation (9). We compose the isomorphism of equation (9) with the inverse of the above to obtain the isomorphism
$$\overline{J}_k(𝐂^rE)J_k(𝐂^r)_P\times 𝐂^r;$$
$$((Z_j^i)_{i=1,\mathrm{},r;j=1,\mathrm{},k};x)(((Z_j^i)_{i=1,\mathrm{},r;j=1,\mathrm{},k};P);x).$$
This isomorphism respects regular and singular jets, since the subset of the singular jets is given in every fiber by $`\{Z_1^i=0,i=1,\mathrm{},r\}`$ by Proposition 1.7.
Let us understand this isomorphism, restricted to $`𝐂^rE`$, in a more geometric way. As in the proof of Proposition 1.4, let
$$\mathrm{\Psi }:𝐂^r𝐂^r;(w_1,\mathrm{},w_r)(e^{w_1},\mathrm{},e^{w_a},w_{a+1},\mathrm{},w_r)=(z_1,\mathrm{},z_r),$$
and let $`W_j^i`$, $`i=1,\mathrm{},r;j=1,\mathrm{},k`$ be the components of the first part of
$$(\stackrel{~}{dw_1},\mathrm{},\stackrel{~}{dw_r})\times \pi :J_k(𝐂^r)(𝐂^k)^r\times 𝐂^r.$$
Then we claim that the above isomorphism, restricted to $`𝐂^rE`$, is given by
$$J_k(𝐂^rE)J_k(𝐂^r)_P\times 𝐂^rE;j_k(f)(j_k(\mathrm{\Psi }(\mathrm{\Psi }^1f\mathrm{\Psi }^1f(0))),f(0)),$$
where subtraction means subtraction in $`𝐂^r`$. Note that $`\mathrm{\Psi }^1`$ is only defined up to addition of summands $`2\pi im,m\mathrm{Z}\mathrm{Z}`$ for the first $`a`$ components, but the germ $`\mathrm{\Psi }^1f\mathrm{\Psi }^1f(0)`$ is well defined. In fact, by Lemma 1.5 the above isomorphism is given by trivial shift with respect to the coordinates $`W_j^i`$, but this is, by the definition of these coordinates, just subtraction of the value of the germ for $`t=0`$. It follows that the above isomorphism commutes with the action of $`G_k`$. In fact, reparametrization does not depend on the coordinates and so it commutes with $`\mathrm{\Psi }`$ and $`\mathrm{\Psi }^1`$. Furthermore, it commutes with subtraction of constants in $`𝐂^r`$. This proves a).
Proof for b) We use the following strategy: Using some results of Demailly’s paper we first define the isomorphism of equation (4) on $`𝐂^rE`$ similarly to our geometric way above. It is then easy to verify the diagram of equation (5). We then extend this isomorphism over the complement of a divisor in the bundle which does not contain any entire fiber over $`𝐂^r`$. For this, we introduce an explicit local coordinate chart in $`\overline{P}_k(𝐂^rE)`$ the complement of which is a divisor which does not contain any fiber over $`𝐂^r`$. In order to extend over the remaining codimension two locus we use the fact that our objects are defined inductively by projectivizing vector bundles, and that for vector bundle maps, the Riemann Extension Theorem holds. This way is not the shortest possible (see Lemma 5.10, which gives a much shorter and intrinsic proof of this extension over the divisor $`E`$), but it explains well the geometric contents of the isomorphism in equation (4) via explicit local coordinates (see also Corollary 3.7), which is useful for applications.
By Corollary 5.12 of , for all points $`wP_k(𝐂^rE)`$, there exists a germ $`f:(𝐂,0)𝐂^rE`$ such that $`f_{[k]}(0)=w`$ and $`f_{[k1]}^{}(0)0`$. We claim that by composing this germ with the map $`tat+t^2`$, $`a𝐂`$, the vector $`f_{[k1]}^{}(0)`$ can be made equal to an arbitrary vector in the complex line $`𝒪_{P_k(𝐂^rE)}(1)`$ over $`w`$. This follows from Theorem 3.2 b) for $`a0`$. Since the image of the germ $`f_{[k]}`$ does not change, we get, after an easy computation, that $`f_{[k1]}^{}(0)=0`$ for $`a=0`$. So every vector of $`𝒪_{P_k(𝐂^rE)}(1)`$ is obtained this way.
As above, the map
$$𝒪_{P_k(𝐂^rE)}(1)(𝒪_{P_k(𝐂^rE)}(1))_P\times (𝐂^rE);$$
$$f_{[k1]}^{}(0)(\mathrm{\Psi }(\mathrm{\Psi }^1f\mathrm{\Psi }^1f(0))_{[k1]}^{}(0);f(0))$$
is a well defined isomorphism, its inverse being given by the well defined map
$$(𝒪_{P_k(𝐂^rE)}(1))_P\times 𝐂^rE𝒪_{P_k(𝐂^rE)}(1);$$
$$(f_{[k1]}^{}(0),p)\mathrm{\Psi }(\mathrm{\Psi }^1f+\mathrm{\Psi }^1(p))_{[k1]}^{}(0).$$
Now, by Proposition 5.11 of , the singular locus of $`𝒪_{P_k(𝐂^rE)}(1)`$ can be characterized by $`f_{[k1]}^{}=0`$ along with the multiplicities of $`f_{[j]}`$, $`j=0,\mathrm{},k1`$, which remain invariant under changes of coordinates or additions by constants. So this isomorphism respects the regular and singular jet loci. So we can restrict it to the regular loci. If now $`j_k(f)J_k(𝐂^rE)^{\mathrm{reg}}`$, then this is mapped to $`(\mathrm{\Psi }(\mathrm{\Psi }^1f\mathrm{\Psi }^1f(0))_{[k1]}^{}(0),f(0))`$ by both compositions of the maps in the diagram in equation (5), so this diagram commutes.
We now carry out the above strategy via the following lemma. Let $`\overline{V}_k^{\mathrm{reg}}=\overline{V}_k\overline{P}_kV|_{\overline{P}_kV^{\mathrm{reg}}}`$, where the $`\overline{P}_kV`$ denotes the zero section of $`\overline{V}_k`$. Then we have a canonical identification
$$𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}\stackrel{}{}\overline{V}_{k1}^{\mathrm{reg}}.$$
We now introduce $`r`$ coordinate charts on $`\overline{V}_{k1}^{\mathrm{reg}}`$ in the same way as Demailly did for $`P_kV^{\mathrm{reg}}`$ in equations (4.9) and (5.7) and Theorem 6.8 of .
###### Lemma 3.5
Let $`((Z_j^i)_{i=1,\mathrm{},r;j=1,\mathrm{},k};(z_1,\mathrm{},z_r))`$ be the coordinates of $`\overline{J}_k(𝐂E)`$, and let
$$\overline{A}_{k,r}=(\underset{j=2}{\overset{k}{}}\{Z_j^r=0\})(\overline{J}_k(𝐂E)\{Z_1^r=0\}).$$
Then the map
$$\stackrel{~}{\alpha }_k:\overline{A}_{k,r}|_{𝐂^rE}\overline{V}_{k1}|_{𝐂^rE}$$
(6)
extends over $`E`$ to a map which is biholomorphic onto its image $`\stackrel{~}{\alpha }_k(A_{k,r})`$, such that this image contains the complement of a divisor in $`\overline{V}_{k1}`$ which is nowhere dense in $`(\overline{V}_{k1})_x`$ for all $`x𝐂^r`$. More precisely:
Claim S(k): $`\overline{V}_{k1}\overline{P}_{k1}(𝐂^rE)`$ is a vector bundle of rank $`r`$ over a $`(k1)`$-stage tower of $`𝐏^{r1}`$-bundles, and we can introduce inhomogenous coordinates on these bundles corresponding to the coordinates $`(z_1,\mathrm{},z_r)`$ of $`𝐂^r`$, in which the map $`\stackrel{~}{\alpha }_k`$ of equation (6) is given by
$$((Z_j^i)_{i=1,\mathrm{},r1;j=1,\mathrm{},k},Z_1^r;(z_1,\mathrm{},z_r))((\frac{Z_1^1}{Z_1^r},\mathrm{},\frac{Z_1^{r1}}{Z_1^r}),(\frac{Z_2^1}{(Z_1^r)^2},\mathrm{},\frac{Z_2^{r1}}{(Z_1^r)^2}),\mathrm{},$$
$$(\frac{Z_{k1}^1}{(Z_1^r)^{k1}},\mathrm{},\frac{Z_{k1}^{r1}}{(Z_1^r)^{k1}}),(\frac{Z_k^1}{(Z_1^r)^{k1}},\mathrm{},\frac{Z_k^{r1}}{(Z_1^r)^{k1}},Z_1^r);(z_1,\mathrm{},z_r)).$$
Proof of Lemma 3.5 It suffices to prove S(k). We prove this by induction. The statement S(1) is trivial. Assume by induction that S(k-1) holds. Then the corresponding inhomogenous coordinates of $`\overline{P}_{k1}(𝐂^rE)`$ are given by
$$((\frac{Z_1^1}{Z_1^r},\mathrm{},\frac{Z_1^{r1}}{Z_1^r}),(\frac{Z_2^1}{(Z_1^r)^2},\mathrm{},\frac{Z_2^{r1}}{(Z_1^r)^2}),\mathrm{},(\frac{Z_{k1}^1}{(Z_1^r)^{k1}},\mathrm{},\frac{Z_{k1}^{r1}}{(Z_1^r)^{k1}});(z_1,...,z_r)).$$
(7)
In order to find coordinates over this affine chart, we proceed in two steps:
The first step is to find the coordinates of the logarithmic tangent bundle
$$\overline{T}P_{k1}(𝐂^rE)\overline{P}_{k1}(𝐂^rE)$$
over our affine chart. Note that, in the coordinates of equation (7), the divisor $`E_{k1}=\pi _{0,k1}(E)`$ is given by $`\{z_1\mathrm{}z_a=0\}`$ and, hence, is independent of any of the fiber coordinates $`Z_j^i`$ or of their quotients $`Z_j^i/(Z_1^r)^j`$. So the coordinates of of $`\overline{T}P_{k1}(𝐂^rE)`$ are given by those of equation (7) and their differentials, except for $`z_1,\mathrm{},z_a`$, where the log-differentials are needed.
The second step is to restrict this coordinate system to the subbundle $`\overline{V}_{k1}\overline{T}P_{k1}(𝐂^rE)`$. By the definition of this subbundle in equation (1) (see also Demailly’s equation (5.7) in ) we choose the differentials of the $`r1`$ coordinate functions of equation (7) which describe the fibers of the map $`\pi _{k1}:\overline{P}_{k1}(𝐂^rE)\overline{P}_{k2}(𝐂^rE)`$ and of an extra component which corresponds to how we have chosen the inhomogenous coordinates: These are the (nonlog-) differentials $`d(\frac{Z_{k1}^i}{(Z_1^r)^{k1}})`$, $`i=1,\mathrm{},r`$ plus the extra component, corresponding to the log- (in case $`a=r`$) or nonlog- (in case $`a<r`$) differential of $`z_r`$, which is $`Z_1^r`$ (see equations (3) and (9)).
It remains to express these coordinates without the differential as in claim S(k). We have for $`1ir1`$:
$$d(\frac{Z_{k1}^i}{(Z_1^r)^{k1}})=\frac{dZ_{k1}^i}{(Z_1^r)^{k1}}\frac{Z_{k1}^i}{(Z_1^r)^{k1}}(k1)\frac{dZ_1^r}{Z_1^r}.$$
By equations (3) and (9) we get
$$dZ_{k1}^i(j_k(f))=(\frac{d}{dt}\frac{d^{k2}}{dt^{k2}}\frac{f^{}\omega ^i}{dt})|_{t=0}=(\frac{d^{k1}}{dt^{k1}}\frac{f^{}\omega ^i}{dt})|_{t=0}=Z_k^i(j_k(f)).$$
As we only work on the submanifold $`\overline{A}_{k,r}`$, we have $`j_k(z_rf)=a_0(x)+a_1(x)t`$. We now again use equations (3) and (9), and distinguish two cases:
If $`a<r`$, then
$$dZ_r^1(j_k(f))=(\frac{d}{dt}\frac{f^{}dz_r}{dt})|_{t=0}=(\frac{d}{dt}a_1(x))|_{t=0}=0.$$
If $`a=r`$, then
$$dZ_r^1(j_k(f))=(\frac{d}{dt}\frac{f^{}dz_r}{(z_rf)dt})|_{t=0}=(\frac{d}{dt}\frac{a_1(x)}{a_0(x)+a_1(x)t})|_{t=0}=$$
$$=(\frac{a_1(x)}{a_0(x)})^2=(\frac{f^{}dz_r}{(z_rf)dt}|_{t=0})^2=(Z_1^r)^2(j_k(f)).$$
So we have
$$d(\frac{Z_{k1}^i}{(Z_1^r)^{k1}})=\frac{Z_k^i}{(Z_1^r)^{k1}}+(k1)\{\begin{array}{ccc}0& :& a<r\\ Z_1^r& :& a=r\end{array}.$$
These coordinates can be expressed in those of claim S(k) and vice versa. $``$$``$
Now we use Lemma 3.5 to extend the diagram in equation (4), and, thus, the isomorphism of equation (4): The diagram becomes
$$\begin{array}{ccc}\overline{A}_{k,r}|_{𝐂^rE}& & (A_{k,r})_P\times (𝐂^rE)\\ & & \\ \stackrel{~}{\alpha }_k& & (\stackrel{~}{\alpha }_k)_P\times id\\ & & \\ \overline{V}_{k1}|_{𝐂^rE}& & (V_{k1})_P\times (𝐂^rE)\end{array}$$
where the two vertical arrows are now biholomorphic onto their images. By Lemma 3.5 these isomorphisms extend to isomorphisms over $`𝐂^r`$. So the isomorphism of equation (4) extends over $`E`$ outside a horizontal divisor which is nowhere dense in all fibers, giving an isomorphism outside an analytic set of codimension at least two.
We finally prove, by induction over $`k`$, that these isomorphisms extend to
$$\begin{array}{ccc}\overline{V}_{k1}& & (V_{k1})_P\times 𝐂^r\\ & & \\ & & \\ & & \\ \overline{P}_{k1}V& & (P_{k1}V)_P\times 𝐂^r\end{array}$$
which induce the desired isomorphisms of equation (4). The case S(1) is trivial (there is nothing to extend any more in this case). Assume by induction that S(k-1) is true. Then by projectivizing we have an isomorphism $`\overline{P}_{k1}V(P_{k1}V)_P\times 𝐂^r`$, and over this we have an isomorphism $`\overline{V}_{k1}(V_{k1})_P\times 𝐂^r`$ up to a subvariety of codimension two. Now this isomorphism extends, since for vector bundle maps, the Riemann Extension Theorem holds. (For any point $`w\overline{P}_{k1}V`$, take a dual basis of $`\overline{V}_{k1}`$ around $`w`$. Then the extension of the maps, in both directions, is reduced to extension of holomorphic functions once we compose these maps with the dual vectors.) The fact that the extended maps are still inverses to each other follows from the Identity Theorem. This ends the proof of Proposition 3.4. $``$$``$
Important Remark 1 The local isomorphisms of equations (3) and (4) are fiber bundle isomorphisms. But they are not induced by (directed) morphisms. As a result, these local isomorphisms have a priori no functoriality, and every compatibility which one needs has to be proved explicitly, which we proceed to do.
###### Proposition 3.6
Let $`(\overline{X},D,\overline{V})`$ be a log-directed manifold. Let $`x_0\overline{X}`$ and let $`U`$ and the log-directed projection
$$K:(\overline{X},D,\overline{V})|_U(𝐂^r,E,\overline{T}𝐂^r);(z_1,\mathrm{},z_n)(z_1,\mathrm{}z_a,z_{l+1},\mathrm{},z_{l+b}),$$
with $`E=\{z_1\mathrm{}z_a=0\}`$ and $`a+b=r=\mathrm{rank}V`$, be as in Proposition 2.2. Let, without loss of generality, $`P=(\stackrel{l}{\stackrel{}{1,\mathrm{},1}},\stackrel{nl}{\stackrel{}{0,\mathrm{},0}})U`$.
* The isomorphisms of Proposition 3.4, Proposition 1.4 and Proposition 2.2 induce isomorphisms
$$\overline{J}_kV|_UJ_kV_P\times U$$
and
$$𝒪_{\overline{P}_kV}(1)|_U𝒪_{P_kV}(1)_P\times U$$
respecting regular and singular jets in such a way that the first isomorphism commutes (outside $`D`$) with the action of $`G_k`$ and that the diagrams
$$\begin{array}{cccc}\hfill K^1\left(\overline{J}_k\left(𝐂^rE\right)\right)|_U& & \left(K^1\left(J_k𝐂^r\right)\right)_P\times U& \\ & & & \\ \hfill & & & \\ & & & \\ & \overline{J}_kV|_U& & \hfill J_kV_P\times U\end{array}$$
(8)
and
$$\begin{array}{cccc}\hfill K^1\left(𝒪_{\overline{P}_k\left(𝐂^rE\right)}\left(1\right)\right)|_U& & \left(K^1\left(𝒪_{P_k𝐂^r}\left(1\right)\right)\right)_P\times U& \\ & & & \\ \hfill & & & \\ & & & \\ & 𝒪_{\overline{P}_kV}\left(1\right)|_U& & \hfill 𝒪_{P_kV}\left(1\right)_P\times U\end{array}$$
(9)
commute.
* Moreover, outside the divisor $`D`$, they induce the following cubic diagram
$$\begin{array}{cccc}K^1\left(\overline{J}_k\left(𝐂^rE\right)^{\mathrm{reg}}\right)|_{UD}& & \left(K^1\left(J_k𝐂^r\right)^{\mathrm{reg}}\right)_P\times \left(UD\right)& \\ & & & \\ & & & \\ & & & \\ & \overline{J}_kV^{\mathrm{reg}}|_{UD}& & J_kV_P^{\mathrm{reg}}\times \left(UD\right)\\ & & & \\ \stackrel{~}{\alpha }_k& & \left(\stackrel{~}{\alpha }_k\right)_P\times id& \\ & & & \\ & & & \\ & \stackrel{~}{\alpha }_k& & \left(\stackrel{~}{\alpha }_k\right)_P\times id\\ & & & \\ & & & \\ K^1\left(𝒪_{\overline{P}_k\left(𝐂^rE\right)}\left(1\right)^{\mathrm{reg}}\right)|_{UD}& & \left(K^1\left(𝒪_{P_k𝐂^r}\left(1\right)\right)^{\mathrm{reg}}\right)_P\times \left(UD\right)& \\ & & & \\ & & & \\ & & & \\ & 𝒪_{\overline{P}_kV}\left(1\right)^{\mathrm{reg}}|_{UD}& & 𝒪_{P_kV}\left(1\right)_P^{\mathrm{reg}}\times \left(UD\right)\end{array}$$
(10)
* By combining with the canonical line bundle projections we get the same isomorphisms and diagrams with $`\overline{P}_k(𝐂^rE)`$, $`\overline{P}_kV`$ and $`\alpha _k`$ instead of $`𝒪_{\overline{P}_k(𝐂^rE)}(1)`$, $`𝒪_{\overline{P}_kV}(1)`$ and $`\stackrel{~}{\alpha }_k`$.
Proof for a) We define the isomorphisms $`\overline{J}_kV|_UJ_kV_P\times U`$ respectively $`𝒪_{\overline{P}_kV}(1)|_U𝒪_{P_kV}(1)_P\times U`$ by the other three arrows of the respective diagrams. In this way we obtain trivializations which, by definition, make the diagrams commutative. By Proposition 3.4, Proposition 1.7, Proposition 1.4 and Corollary 2.3, the regular and singular loci are preserved. The first isomorphism commutes with the action of $`G_k`$. In fact, by Proposition 3.4 a) this is true for the upper line of the diagram in equation (8). Furthermore, the isomorphism $`K_k`$ in the vertical arrows is, outside $`D`$, just given by $`K_k(j_k(f))=j_k(Kf)`$, and this trivially commutes with the action of $`G_k`$.
Proof for b) The back side of this cubic diagram (the side with the $`K^1`$) is just the pull back the diagram in equation (5). The upper and the lower sides of the cubic diagram are the restrictions of the diagrams in equations (8) respectively (9) to the regular locus over $`UD`$. The two vertical arrows on the front side are defined by the left hand side respectively the right hand side of the cubic diagram, so these two sides commute by definition. It is an easy exercise to see that then the the front side of the cubic diagram commutes also and, furthermore, that the whole diagram commutes.
Proof for c) This is clear from the diagrams. $``$$``$
Important Remark 2 For all local isomorphisms given by the horizontal left-to-right arrows in the above diagrams, our Important Remark 1 also applies. However, the local isomorphisms induced by $`K`$ are functorial.
###### Corollary 3.7
* The fiber bundles $`\overline{P}_kV`$, $`𝒪_{\overline{P}_kV}(1)`$ and $`\overline{V}_k`$ and their regular and singular jet loci are all locally trivialized over $`\overline{X}`$ in a way which is compatible, through the maps $`\alpha _k`$ respectively $`\stackrel{~}{\alpha _k}`$, with the trivialization of $`\overline{J}_kV`$ by using local logarithmic coordinates .
* Let $`U\overline{X}`$ and $`K`$ be like in Proposition 2.2. Let $`\overline{A}_{k,i}\overline{J}_k(𝐂^rE)`$, $`i=1,\mathrm{},r`$, be like in Lemma 3.5, and let $`\overline{B}_{k,i}=\overline{A}_{k,i}\{Z_1^i=1\}`$. Then there exist $`r`$ coordinate charts
$$K^1(\overline{A}_{k,i})\overline{V}_{k1}(\mathrm{respectively}K^1(\overline{A}_{k,i})𝒪_{\overline{P}_kV}(1));$$
$$((Z_j^i)_{i=1,\mathrm{},r1;j=1,\mathrm{},k},Z_1^r;(z_1,\mathrm{},z_n))((\frac{Z_1^1}{Z_1^r},\mathrm{},\frac{Z_1^{r1}}{Z_1^r}),(\frac{Z_2^1}{(Z_1^r)^2},\mathrm{},\frac{Z_2^{r1}}{(Z_1^r)^2}),\mathrm{},$$
$$(\frac{Z_{k1}^1}{(Z_1^r)^{k1}},\mathrm{},\frac{Z_{k1}^{r1}}{(Z_1^r)^{k1}}),(\frac{Z_k^1}{(Z_1^r)^{k1}},\mathrm{},\frac{Z_k^{r1}}{(Z_1^r)^{k1}},Z_1^r);(z_1,\mathrm{},z_n))$$
which cover $`\overline{V}_{k1}^{\mathrm{reg}}`$ (respectively $`𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}`$), and $`r`$ coordinate charts
$$K^1(\overline{B}_{k,i})\overline{P}_{k1}V;$$
$$((Z_j^i)_{i=1,\mathrm{},r1;j=1,\mathrm{},k};(z_1,\mathrm{},z_n))$$
$$((Z_1^1,\mathrm{},Z_1^{r1}),(Z_2^1,\mathrm{},Z_2^{r1}),\mathrm{},(Z_k^1,\mathrm{},Z_k^{r1});(z_1,\mathrm{},z_n))$$
which cover $`\overline{P}_{k1}V^{\mathrm{reg}}`$.
Proof a) is contained in Proposition 3.6. The existence of the coordinate charts in b) follows from Proposition 3.6 and Lemma 3.5. These charts cover the locus of regular jets by Theorem 3.2, a) outside $`D`$. Thus by our local trivializations which are compatible with the charts and with the locus of regular jets, these charts cover the locus of regular jets everywhere. $``$$``$
Remark The coordinates can also be obtained directly without Lemma 3.5.
### 3.3 Log-directed jets and log-Demailly-Semple jets
This subsection extends Theorem 3.2 a), b) and c) to the log-directed case.
###### Proposition 3.8
Let $`(\overline{X},D,\overline{V})`$ be a log-directed manifold.
* The maps $`\stackrel{~}{\alpha }_k`$ and $`\alpha _k`$ of Theorem 3.2 a) extend to holomorphic and surjective maps
$$\stackrel{~}{\alpha }_k:\overline{J}_kV^{\mathrm{reg}}𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}},$$
$$\alpha _k:\overline{J}_kV^{\mathrm{reg}}\overline{P}_kV^{\mathrm{reg}}.$$
* The action of $`\varphi G_k`$ extends to an automorphism of $`\overline{J}_kV`$ leaving $`\overline{J}_kV^{\mathrm{reg}}`$ and $`\overline{J}_kV^{\mathrm{sing}}`$ invariant and satisfying
$$\stackrel{~}{\alpha }_k\varphi =\stackrel{~}{\alpha }\varphi ^{}(0),\alpha _k\varphi =\alpha _k.$$
* The quotient $`\overline{J}_kV^{\mathrm{reg}}/G_k`$ has the structure of a locally trivial fiber bundle over $`\overline{X}`$, and the map
$$\alpha _k/G_k:\overline{J}_kV^{\mathrm{reg}}/G_k\overline{P}_kV$$
is a holomorphic embedding which identifies $`\overline{J}_kV^{\mathrm{reg}}/G_k`$ with $`\overline{P}_kV^{\mathrm{reg}}`$.
Proof for a) By Theorem 3.2, the map $`\stackrel{~}{\alpha }_k`$ is defined outside $`D`$:
$$\stackrel{~}{\alpha }_k:\overline{J}_kV^{\mathrm{reg}}|_{\overline{X}D}𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}|_{\overline{X}D}.$$
(11)
Let $`xD`$. By Proposition 3.6 b), there exists a neighborhood $`U`$ of $`x`$ with
$$\begin{array}{ccc}\overline{J}_kV^{\mathrm{reg}}|_{UD}& & J_kV_P^{\mathrm{reg}}\times (UD)\\ & & \\ \stackrel{~}{\alpha }_k& & (\stackrel{~}{\alpha }_k)_P\times id\\ & & \\ 𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}|_{UD}& & 𝒪_{P_kV}(1)_P^{\mathrm{reg}}\times (UD)\end{array}$$
(12)
Here the horizontal arrows are isomorphisms, which, by Proposition 3.6a), extend as isomorphisms over $`U`$, and $`(\stackrel{~}{\alpha }_k)_P\times id`$ is clearly extendable to $`U`$ to a surjective holomorphic map on the right hand side. So $`\stackrel{~}{\alpha _k}`$ is also extendable to a surjective holomorphic map over $`U`$ on the left hand side. Since $`xD`$ is arbitrary, and since by equation (11) the extension of $`\stackrel{~}{\alpha }_k`$ is unique if it exists, we obtain a well defined surjective holomorphic map
$$\stackrel{~}{\alpha }_k:\overline{J}_kV^{\mathrm{reg}}𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}.$$
(13)
By combining with the canonical line bundle projections we get in the same way a surjective holomorphic map
$$\alpha _k:\overline{J}_kV^{\mathrm{reg}}\overline{P}_kV^{\mathrm{reg}}$$
(14)
which extends the corresponding map $`\alpha _k`$ of Theorem 3.2 from $`\overline{X}D`$ to $`\overline{X}`$.
Proof for b) If $`\varphi G_k`$ is a reparametrization, one has on $`\overline{J}_kV^{\mathrm{reg}}|_{\overline{X}D}`$ by Theorem 3.2:
$$\stackrel{~}{\alpha }_k\varphi =\stackrel{~}{\alpha }_k\varphi ^{}(0),\alpha _k\varphi =\alpha _k,$$
(15)
where in the first equation the multiplication $`\stackrel{~}{\alpha _k}\varphi ^{}(0)`$ denotes the multiplication with scalars in the line bundle $`𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}|_{\overline{X}D}`$. By Proposition 3.6 a), we have the diagram
$$\begin{array}{ccc}\overline{J}_kV|_{UD}& & J_kV_P\times (UD)\\ & & \\ \varphi & & \varphi _P\times id\\ & & \\ \overline{J}_kV|_{UD}& & J_kV_P\times (UD)\end{array}.$$
(16)
By a similar argument as in a), the map $`\varphi `$ extends to a holomorphic automorphism on $`\overline{J}_kV`$. From this diagram, it also follows that $`\varphi `$ maps $`\overline{J}_kV^{\mathrm{reg}}`$ onto itself, since all arrows of this diagram preserve regular and singular jets, and by Proposition 1.7 this remains true over $`D`$. Finally, equation (15) extends from $`\overline{J}_kV^{\mathrm{reg}}|_{\overline{X}D}`$ to $`\overline{J}_kV^{\mathrm{reg}}`$ by the Identity Theorem.
Proof for c) By b), the quotient $`\overline{J}_kV^{\mathrm{reg}}/G_k`$ is well defined (as set). By the diagrams of equations (15) and (16), we obtain from Proposition 3.6 c):
$$\begin{array}{ccc}\overline{J}_kV^{\mathrm{reg}}/G_k|_{UD}& & (J_kV^{\mathrm{reg}}/G_k)_P\times (UD)\\ & & \\ \alpha _k/G_k& & (\alpha _k/G_k)_P\times id\\ & & \\ \overline{P}_kV^{\mathrm{reg}}|_{UD}& & P_kV_P^{\mathrm{reg}}\times (UD)\end{array}$$
(17)
By Demailly (), the vertical arrows in this diagram are isomorphisms. By a similar argument as in a), one obtains a holomorphic isomorphism
$$\alpha _k/G_k:\overline{J}_kV^{\mathrm{reg}}/G_k\overline{P}_kV^{\mathrm{reg}}$$
(18)
over $`\overline{X}`$. Equation (17) shows that this isomorphism makes $`\overline{J}_kV^{\mathrm{reg}}/G_k`$ into a holomorphic fiber bundle over $`\overline{X}`$. $``$$``$
### 3.4 Characterization of log-directed jet differentials
In this subsection we generalize Theorem 3.2 d). More precisely we prove:
###### Proposition 3.9
A holomorphic (respectively meromorphic) function $`Q`$ on $`\overline{J}_kV|_O`$ for some connected open subset $`O\overline{X}`$ which satisfies
$$Q(j_k(f\varphi ))=\varphi ^{}(0)^mQ(j_k(f))j_k(f)J_kV^{\mathrm{reg}}\text{ and }\varphi G_k$$
(19)
over some open subset of $`O^{}`$ of $`OD`$ defines a holomorphic (respectively meromorphic) section of $`𝒪_{\overline{P}_kV}(m)`$ over $`O`$, and vice versa.
Proof Let $`Q:\overline{J}_kV|_O𝐂`$ be a meromorphic function which satisfies
$$Q\varphi =\varphi ^{}(0)^mQ\varphi G_k$$
(20)
over $`O^{}`$. Since $`\overline{J}_kV^{\mathrm{reg}}|_O`$ is connected, equation (20) holds over $`O`$ by the Identity Theorem. Since $`\stackrel{~}{\alpha }_k`$ and $`\alpha _k`$ were obtained over $`D`$ by trivial extensions in the diagrams of Proposition 3.6, the results of Corollary 3.3 extend also over $`D`$. In particular, $`G_k^o`$ of Corollary 3.3 acts transitively on the fibers of $`\stackrel{~}{\alpha }_k`$. Since the function $`Q`$ is invariant under the action of $`G_k^o`$ by equation (20), there exists a Zariski-densely defined function $`\stackrel{~}{Q}:𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}|_O𝐂`$ such that $`Q=\stackrel{~}{Q}\stackrel{~}{\alpha }_k`$. Again by Corollary 3.3, $`\stackrel{~}{\alpha }_k`$ has local holomorphic sections everywhere. So $`\stackrel{~}{Q}`$ is a meromorphic function on $`𝒪_{\overline{P}_kV}(1)^{\mathrm{reg}}|_O`$. By equation (20), this function is $`m`$-linear with respect to the $`𝐂^{}`$-action and so corresponds to a meromorphic section $`s`$ of $`𝒪_{\overline{P}_kV}(m)^{\mathrm{reg}}|_U`$. In order to extend this section to the singular locus, we have to redo an argument of Demailly (): In a neighborhood $`W`$ of any point $`w_0\overline{P}_kV|_{OD}`$ we can find a holomorphic family of germs $`f_w`$ such that $`(f_w)_{[k]}(0)=w`$ and $`(f_w)_{[k1]}^{}(0)0`$. Then we get $`s(w)=Q(f_w^{},\mathrm{},f_w^{(k)})(0)((f_w)_{[k1]}^{})^m`$ on $`W\overline{P}_kV|_{OD}`$. Now the right hand side extends to a section of $`𝒪_{\overline{P}_kV}(m)|_{UD}`$, so the left hand side does, too. So $`s`$ is a meromorphic section of $`𝒪_{\overline{P}_kV}(m)|_{\overline{P}_kV^{\mathrm{reg}}\overline{P}_kV|_{OD}}`$. The complement of the latter set is of codimension two in $`\overline{P}_kV|_O`$, so $`s`$ extends to a section of $`𝒪_{\overline{P}_kV}(m)|_O`$.
Conversely let $`\stackrel{~}{Q}`$ be an $`m`$-linear meromorphic function on $`𝒪_{\overline{P}_kV}(1)|_O`$ corresponding to a meromorphic section of $`𝒪_{\overline{P}_kV}(m)|_O`$. Then $`Q:=\stackrel{~}{Q}\stackrel{~}{\alpha }_k`$ is a meromorphic function on $`\overline{J}_kV^{\mathrm{reg}}|_O`$. By the Riemann Extension Theorem, it extends to a meromorphic function on $`\overline{J}_kV|_U`$. It satisfies equation (20) on $`\overline{J}_kV^{\mathrm{reg}}|_O`$ since $`\stackrel{~}{Q}`$ corresponds to a section of $`𝒪_{\overline{P}_kV}(m)|_O`$ and since the fibers of $`\stackrel{~}{\alpha }_k`$ are invariant under the action of $`G_k^o`$. Hence, it satisfies equation (20) over $`O`$ by the Identity Theorem.
Finally we remark that if we start with a holomorphic rather than a meromorphic function (respectively, section) in the arguments above, we would obtain a holomorphic section (respectively, function) as a result. $``$$``$
## 4 Log-directed jet metrics
### 4.1 The case of 1-jets
This case was already treated in the second named author’s thesis. We recall the basic results after some definitions.
For a line bundle $`L`$ over a complex variety $`X`$, let $`E_L`$ be the union of the base locus
$$\text{Bs}|L|:=\{xX:s(x)=0\text{ for all }sH^0(X,L)\},$$
of $`L`$ and the restricted exceptional locus
$$\{xX\text{Bs}|L|:dim_x\phi _L^1(\phi _L(x))>0\}$$
of the rational map
$$\phi _L:=[s_1:\mathrm{}:s_n]:X\mathrm{}𝐏^{n1},$$
where $`\{s_1,\mathrm{},s_n\}`$ is a basis of $`H^0(X,L)`$. We will call $`E_L`$ the basic locus of $`L`$. Define the stable basic locus of $`L`$ to be
$$S_L:=\underset{m>0}{}E_{mL}.$$
A standard argument (worked out in details in the appendix) shows that for any line bundle $`H`$ on a normal variety $`X`$, we have
$$\text{Bs}|mLH|S_L$$
for some sufficiently large $`m`$. Let $`X_{\mathrm{reg}}`$ be the smooth part of $`X`$. Let $`L^{}`$ be the dual bundle of $`L`$. Recall that a continuous function $`g:L^{}[0,\mathrm{}]`$ such that
$$g(cv)=|c|^2g(v)$$
(1)
for all $`c𝐂`$ and $`vL^{}`$ is called a singular metric on $`L^{}`$. By equation (1), the set $`g^1(0)g^1(\mathrm{})`$ consists of the zero section of $`L^{}`$ and the inverse image in $`L^{}`$ of a closed subset $`\mathrm{\Sigma }_g`$ of $`X`$. For our purpose, we will always assume that the open set $`U=X_{\mathrm{reg}}\mathrm{\Sigma }_g`$ is dense in $`X`$ and that $`g`$ is twice differentiable on $`L^{}|_U`$. Then $`dd^c\mathrm{log}g`$ is a real $`(1,1)`$ form outside the zero section of $`L^{}|_U`$ invariant under the $`C^{}`$ action given by equation (1) and is thus the pull back of a real $`(1,1)`$ form on $`U`$ denoted by
$$\mathrm{Ric}(g)=\mathrm{\Theta }_{g^1}=\mathrm{\Theta }_{g^1}(L),$$
which is known as the curvature form of $`g`$. By convention, $`g`$ is called a pseudometric if $`g^1(\mathrm{})=\mathrm{}`$ and $`g`$ is called a metric if $`\mathrm{\Sigma }_g=\mathrm{}`$.
Let now $`(\overline{X},D)`$ be a log-manifold, and set again $`X=\overline{X}D`$. A Kähler metric $`\omega `$ on $`X`$ gives a metric on $`TX`$ which in turn gives a metric $`g_\omega `$ on $`𝒪_{\overline{P}_1X}(1)|_{P_1X}`$. If $`\omega `$ behaves logarithmically along $`D`$, then $`g_\omega `$ extends to a metric on $`𝒪_{\overline{P}_1X}(1)`$ which we can use to dominate a scalar multiple of any pseudometics on $`𝒪_{\overline{P}_1X}(1)`$ by appealing to the compactness of $`\overline{X}`$. This is the basic strategy used to obtain the following result (Proposition 1 of ). We remark that Noguchi () had already similar results in the case $`X`$ is compact under the assumption that $`𝒪_{P_1X}(m)`$ is spanned by global sections everywhere on $`𝒪_{P_1X}(m)`$ for $`m`$ large enough.
###### Proposition 4.1
Assume $`(\overline{X},D)`$ is a log-manifold and that $`\overline{X}`$ is projective. Let $`\pi _1:\overline{P}_1X\overline{X}`$ be the natural projection, let $`\overline{\mathrm{\Xi }}`$ be a subvariety of $`\overline{P}_1X`$ and let $`\sigma :\overline{Z}\overline{\mathrm{\Xi }}\overline{P}_1X`$ be the normalization of $`\overline{\mathrm{\Xi }}`$. Let $`\overline{L}_\sigma =\sigma ^1𝒪_{\overline{P}_1X}(1)`$, $`Z=\sigma ^1(P_1X)`$ and $`L_\sigma =\overline{L}_\sigma |_Z`$. Then there is a pseudometric $`g`$ on $`L_\sigma ^{}=\sigma ^1(𝒪_{P_1X}(1))`$ with $`\mathrm{\Sigma }_gS_{\overline{L}_\sigma }`$, such that $`\mathrm{Ric}(g)`$ is the pullback of a Kähler metric $`\omega `$ on $`X`$, specifically
$$\mathrm{Ric}(g)=(\sigma \pi _1)^{}\omega ,$$
such that this Kähler metric $`\omega `$ dominates $`g`$, in the sense that
$$(\sigma ^{}g_\omega )(\xi )g(\xi )$$
(2)
for all $`\xi L_\sigma ^{}`$ outside $`S_{\overline{L}_\sigma }`$ and $`\sigma ^1(\mathrm{Sing}(\overline{\mathrm{\Xi }}))`$.
By the usual definition of holomorphic sectional curvature, we see that equation (2) says precisely that $`g`$, as a “length” function on $`X`$ in the tangent directions defined by $`\overline{\mathrm{\Xi }}`$, has holomorphic sectional curvature bounded from above by $`1`$, and $`g`$ is nonvanishing outside $`S_{\overline{L}_\sigma }`$.
Hence, the usual Ahlfors’ Lemma applies to show that if $`f:\mathrm{\Delta }X`$ is any holomorphic map from the unit disk $`\mathrm{\Delta }𝐂`$ whose lifting $`f_{[1]}`$ has image in $`\overline{\mathrm{\Xi }}`$ but not completely in $`\sigma (S_{\overline{L}_\sigma })\mathrm{Sing}(\overline{\mathrm{\Xi }})`$, then $`f`$ must satisfy the distance decreasing property.
From this, the following result is derived by elementary arguments in (see also Noguchi ()), which we quote.
###### Theorem 4.2
With the same setup as in Proposition 4.1, we let $`\mathrm{\Delta }^{}=\mathrm{\Delta }\{0\}`$ be the punctured unit disk and set $`\overline{L}_0=𝒪_{\overline{P}_1X}(1)|_{\overline{\mathrm{\Xi }}}`$.
(a) (Distance decreasing property) If $`f:\mathrm{\Delta }X`$ is a holomorphic map whose lift $`f_{[1]}`$ has values in $`\mathrm{\Xi }`$ but not entirely in $`S_{\overline{L}_0}\mathrm{Sing}(\overline{\mathrm{\Xi }})`$, then $`f^{}g\rho `$ (where $`\rho `$ is the Poincaré metric on $`\mathrm{\Delta }`$).
(b) (Degeneracy of Holomorphic Curve) If $`f:𝐂X`$ is holomorphic such that $`f_{[1]}`$ has values in $`\overline{\mathrm{\Xi }}`$, then $`f_{[1]}(𝐂)S_{\overline{L}_0}\mathrm{Sing}(\overline{\mathrm{\Xi }})`$.
(c) (Big Picard Theorem) If $`f:\mathrm{\Delta }^{}X`$ is holomorphic such that $`f_{[1]}`$ has values in $`\overline{\mathrm{\Xi }}`$ but not entirely in $`S_{\overline{L}_0}\mathrm{Sing}(\overline{\mathrm{\Xi }})`$, then $`f`$ extends to a holomorphic map $`\overline{f}:\mathrm{\Delta }\overline{X}`$.
### 4.2 The general case
We call a pseudometric $`h`$ on $`𝒪_{P_kV}(1)`$ a $`k`$-jet pseudometric on $`(X,V)`$, and define $`B_k=S_{𝒪_{\overline{P}_kV}(1)}`$, the stable basic locus of $`𝒪_{\overline{P}_kV}(1)`$.
###### Theorem 4.3
With the notations as in Theorem 4.2, assume that $`(\overline{X},D,\overline{V})`$ is a log-directed manifold and that $`\overline{X}`$ is projective.
* If $`B_k\overline{P}_kV`$, then there exists a k-jet pseudometric $`h`$ on $`(X,V)`$ with $`\mathrm{\Sigma }_hB_k`$ such that $`h`$ has curvature bounded from above by $`1`$ in the sense that $`\mathrm{Ric}(h)=\pi _k^{}\omega `$ is the pullback of a Kähler metric $`\omega `$ on $`P_{k1}V`$ such that $`g_\omega `$ dominates $`h`$. In particular, we have
$$<\mathrm{\Theta }_{h^1},|\xi |^2>=<\mathrm{Ric}(h),|\xi |^2>h((\pi _k)_{}\xi )\mathrm{for}\xi V_k.$$
* If $`f:𝐂\overline{X}D`$ is holomorphic with $`f_{}(T𝐂)\overline{V}`$, then $`f_{[k]}(𝐂)B_k`$.
* If $`f:\mathrm{\Delta }^{}\overline{X}D`$ is holomorphic with $`f_{}(T\mathrm{\Delta }^{})\overline{V}`$, then:
Either $`f`$ extends to a holomorphic map $`\overline{f}:\mathrm{\Delta }\overline{X}`$ or $`f_{[k]}(\mathrm{\Delta }^{})B_k`$.
Moreover, let $`Y\overline{P}_kV`$ be any subvariety. We define $`B_k(Y)=S_{𝒪_{\overline{P}_kV}(1)|_Y}`$. If $`f`$ lifts to a map with values in $`Y`$, then b) and c) hold with $`B_k(Y)\mathrm{Sing}(Y)`$ instead of $`B_k`$.
Proof Apply Proposition 4.1 and Theorem 4.2 to the log-manifold $`(\overline{P}_{k1}V,`$ $`D_{k1})`$ and the subvariety $`\overline{\mathrm{\Xi }}=\overline{P}_kV`$ (or $`\overline{\mathrm{\Xi }}=Y\overline{P}_kV`$). Note that, since $`\overline{V}_k\overline{T}(P_{k1}V)`$ is a holomorphic subbundle, $`\overline{P}_1V_k=\overline{P}_kV`$ is a submanifold in $`\overline{P}_1(P_{k1}V)`$ and $`𝒪_{P_kV}(1)=𝒪_{\overline{P}_1(P_{k1}V)}(1)|_{P_kV}`$ by Proposition 2.4. $``$$``$
## 5 Logarithmic Bloch’s and Lang’s Conjecture
In this section we apply our method to the special case of semi-abelian varieties where our Ahlfors-Schwarz Lemma (Theorem 4.3) gives a logarithmic version of Bloch’s Theorem and our big Picard Theorem yields a big Picard version of Bloch’s Theorem. By using the Wronskian associated to the theta function of an effective divisor in a semi-abelian variety (, ), we affirm furthermore a logarithmic version of Lang’s Conjecture and a big Picard analogue of it, all via metric geometry on Logarithmic Demailly-Semple jets.
### 5.1 Statement of the results
We first recall the definition and some basic facts on semi-abelian varieties (see , , ) needed to state our results.
A quasiprojective variety $`G`$ is called a semi-abelian variety if it is a commutative group which admits an exact sequence of groups
$$0(𝐂^{})^{\mathrm{}}GA0,$$
where $`A`$ is an abelian variety of dimension $`𝗆`$.
Taking the pushforward of $`(𝐂^{})^{\mathrm{}}G`$ with the natural embedding $`(𝐂^{})^{\mathrm{}}(𝐏^1)^{\mathrm{}}`$, we obtain a smooth completion
$$\overline{G}=(𝐏^1)^{\mathrm{}}\times _{(𝐂^{})^{\mathrm{}}}G$$
of $`G`$ with boundary divisor $`S`$, which has only normal crossing singularities. We denote the natual action of $`G`$ on $`\overline{G}`$ on the right as addition. It follows that the exponential map from the Lie algebra $`𝐂^n`$ is a group homomorphism and, hence, it is also the universal covering map of $`G=𝐂^n/\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }=\mathrm{\Pi }_1(G)`$ is a discrete subgroup of $`𝐂^n`$ and $`n=𝗆+\mathrm{}`$.
Following Iitaka (), we have the following trivialization of the logarithmic tangent respectively cotangent bundles of $`\overline{G}`$: Let $`z_1,\mathrm{},z_n`$ be the standard coordinates of $`𝐂^n`$. Since $`dz_1,\mathrm{},dz_n`$ are invariant under the group action of translation on $`𝐂^n`$, they descend to forms on $`G`$. There they extend to logarithmic forms on $`\overline{G}`$ along $`S`$, which are elements of $`H^0(\overline{G},\overline{\mathrm{\Omega }}G)`$. These logarithmic $`1`$-forms are everywhere linearly independent on $`\overline{G}`$. Thus, they globally trivialize the vector bundle $`\overline{\mathrm{\Omega }}G`$. Finally, we note that these logarithmic forms are invariant under the group action of $`G`$ on $`\overline{G}`$, and, hence, the associated trivialization of $`\overline{\mathrm{\Omega }}G`$ over $`\overline{G}`$ is also invariant.
We now state the main theorems of this section. With the above setup, let $`f:\mathrm{\Gamma }G`$ be a holomorphic map, where $`\mathrm{\Gamma }`$ is either $`𝐂`$ or the punctured disk $`\mathrm{\Delta }^{}`$. Denote by $`\overline{X}(f)`$ the Zariski closure of $`f(\mathrm{\Gamma })`$ in $`\overline{G}`$ and let $`X(f)=\overline{X}(f)G`$. Furthermore, let $`DG`$ be a reduced algebraic divisor in $`G`$, which we regard as a union of codimension one algebraic subvarieties of $`G`$. We note that an algebraic subvariety of $`G`$ which is also a subgroup is necessarily a semi-abelian variety as well, see .
###### Theorem 5.1
* Let $`f:𝐂G`$ define a holomorphic curve. Then $`X(f)`$ is a translate of an algebraic subgroup of $`G`$.
* Let $`f:𝐂(GD)`$ be holomorphic. Then $`X(f)D=\mathrm{}`$.
###### Corollary 5.2
If $`D`$ has nonempty intersection with any translate of an algebraic subgroup of $`G`$ (of positive dimension), then $`GD`$ is Brody hyperbolic.
In particular, this holds if $`G=A`$ is an abelian variety and $`D`$ is ample. $``$$``$
Theorem 5.1 (a) is a logarithmic version of Bloch’s Theorem, first proved by Noguchi (), (b) is a logarithmic version of Lang’s Conjecture. Both Theorem 5.1 and Corollary 5.2 were obtained by Noguchi (), and were, in the nonlogarithmic case, first proved by Siu-Yeung ().
###### Theorem 5.3
* Let $`f:\mathrm{\Delta }^{}G`$ be a holomorphic map. Then either it extends to a holomorphic map $`\overline{f}:\mathrm{\Delta }\overline{G}`$, or there exists a maximal algebraic subgroup $`G^{}`$ of $`G`$ of positive dimension such that $`X(f)`$ is foliated by translates of $`G^{}`$.
* Let $`f:\mathrm{\Delta }^{}(GD)`$ be holomorphic. Then one of the following holds:
(i) $`f`$ extends to $`\overline{f}:\mathrm{\Delta }\overline{G}`$.
(ii) $`X(f)D=\mathrm{}`$.
(iii) There is an algebraic subgroup $`G^{\prime \prime }G^{}`$ of positive dimension such that $`X(f)D`$ foliated by translates of $`G^{\prime \prime }`$.
* Let now $`f:\mathrm{\Delta }^{}(AD)`$, where $`G`$ is the special case of an abelian variety $`A`$. Then one of the following holds:
(i) $`f`$ extends to $`\overline{f}:\mathrm{\Delta }A`$.
(ii) There exists an algebraic subgroup $`G^{\prime \prime }G^{}`$ of positive dimension such that $`D`$ is foliated by translates of $`G^{\prime \prime }`$.
###### Corollary 5.4
If $`G=A`$ is an abelian variety and $`D`$ is ample, then $`f:\mathrm{\Delta }^{}AD`$ extends to a holomorphic map $`\overline{f}:\mathrm{\Delta }A`$.
We remark that Theorem 5.3 and Corollary 5.4 are big Picard type Theorems. Aside from (a), which can be found in Noguchi (), these are, to our knowledge, new to the literature.
Proof of Corollary 5.4 Corollary 5.4 follows from Theorem 5.3 (c) and the fact that an ample divisor $`D`$ in an abelian variety $`A`$ cannot be foliated by translates of an algebraic subgroup $`A^{\prime \prime }`$ of $`A`$ of positive dimension. For assume it were. Then $`D=q^1(\overline{D})`$, where $`\overline{D}`$ is a divisor in $`A/A^{\prime \prime }`$ and $`q:AA/A^{\prime \prime }`$ is the quotient map. But then $`𝒪(D)=q^1𝒪(\overline{D})`$ is trivial along $`A^{\prime \prime }`$, since $`A^{\prime \prime }`$ is a fiber of the map $`q`$. This is a contradiction. $``$$``$
Remark The last part of Corollary 5.2 follows from Corollary 5.4 as follows. The $`𝐂\mathrm{\Delta }`$ is biholomorphic to $`\mathrm{\Delta }^{}`$. So we can conclude from Corollary 5.4 that any entire curve $`f:𝐂A`$ extends to a holomorphic map $`\overline{f}:𝐏^1A`$. Hence, $`\overline{f}`$ must be constant, since all coordinate 1-forms on $`A`$ must pull back to the zero on $`𝐏^1`$ as $`𝐏^1`$ has no nontrivial 1-forms.
However, Corollary 5.4 does not follow from Corollary 5.2. It would if $`AD`$ were hyperbolically embedded in $`A`$. This would be the case, for example, if $`D`$ were hyperbolic (see for example ). But even a very ample divisor in $`A`$ is not hyperbolic in general. To see this, choose any translate $`T`$ of an algebraic subvariety which is of codimension at least 3 in $`A`$. Then there always exists an irreducible ample divisor in $`A`$ which contains $`T`$, as can be deduced by applying the Bertini’s Theorem 7.19 in . Note that a hyperbolic open subset $`V`$ in a projective variety $`\overline{V}`$ need not be hyperbolically embedded in general, as one can easily see by blowing up a point of $`\overline{V}V`$.
Remark Let $`G=(𝐂^{})^n𝐏^n`$, $`n4`$. Let $`\overline{D}`$ be a generic hyperplane in $`𝐏^n`$ and $`\overline{H}`$ be another hyperplane with $`G\overline{H}\mathrm{}`$ and $`\overline{H}\overline{D}G=\mathrm{}`$. Then $`\overline{H}(G\overline{D})`$ is equal to $`𝐏^{n1}`$ minus at most $`n+2`$ hyperplanes, which contains nontrivial images of $`𝐂`$ and hence, admits maps $`f`$ from $`\mathrm{\Delta }^{}`$ which do not extend to $`\mathrm{\Delta }`$. This is because the complement of $`n+2`$ hyperplanes in $`𝐏^{n1}`$ contains nontrivial diagonals for $`n4`$, which are nonhyperbolic. So we get examples of $`f`$ for Theorems 5.1 (b) and 5.3 (b) (ii) with nontrivial $`X(f)`$.
Remark Let $`A`$ be an abelian variety, $`DA`$ a divisor and $`f_1:\mathrm{\Delta }AD`$ a holomorphic map. Let $`X(f_1)`$ be the Zariski closure of $`f_1(\mathrm{\Delta })A`$. Let $`E`$ be an elliptic curve and $`q:𝐂E`$ be the universal cover. Then
$$f(z)=(f_1(z),q\mathrm{exp}(\frac{1}{z})):\mathrm{\Delta }^{}A\times E$$
does not extend.
This easy construction provides examples which are relevant to Theorem 5.3:
* It makes Theorem 5.3 (b) (iii) and (c) (ii) sharp.
* Choose $`f_1`$ in such a way that $`X(f_1)`$ is not a translate of an algebraic subgroup in $`A`$. Then we have an example for (a) where $`X(f)`$ is not itself a translate of an algebraic subgroup of $`A\times E`$.
### 5.2 Some results on semi-abelian varieties
We first summarize some elementary properties of semi-abelian varieties.
###### Lemma 5.5
* The quotient of a semi-abelian variety $`G`$ by an algebraic subgroup $`G^{}`$ is again a semi-abelian variety, and the quotient map $`q:GG/G^{}`$ is an algebraic morphism.
* If $`XG`$ is an algebraic variety foliated by translates of $`G^{}`$, then
$$X/G^{}G/G^{}$$
is again an algebraic variety.
* If $`X`$ is an algebraic subvariety of $`G`$, and $`h:X𝐏^1`$ is a rational function, then the closed subgroup
$$G^{}=\{aG:X=(X+a)\}\{aG:h(x)=h(x+a)\mathrm{for}\mathrm{all}xX\}$$
is again an algebraic subvariety.
Proof Lemma 5.5 should be well known, but since we do not know a precise reference we indicate a proof. From the fact that connected algebraic subgroups of a semi-abelian variety are again semi-abelian, it is easy to see that one can consider quotients of $`G`$ by $`G^{}`$ by taking the quotient on the abelian and the $`(𝐂^{})^{\mathrm{}}`$ factors separately. Now the quotient of the abelian factor by an algebraic subgroup is abelian by isogeny, and the quotient of $`(𝐂^{})^{\mathrm{}}`$ by a connected algebraic subgroup is likewise a product of $`𝐂^{}`$, see . Hence, we obtain a $`(𝐂^{})^l`$ bundle over an abelian variety for some $`l`$, which projectivizes to a $`𝐏^l`$ bundle over a projective variety, and, therefore, must be projective. From this the entire Lemma 5.5 follows. $``$$``$
###### Lemma 5.6
Let $`A`$ be an abelian variety and $`DA`$ a reduced algebraic divisor. Let $`A^{}`$ be an algebraic subgroup of $`A`$ and $`T`$ a translate of $`A^{}`$ in $`A`$. Assume $`TD=\mathrm{}`$. Then $`D`$ is foliated by translates of $`A^{}`$.
Proof Without loss of generality we may assume that $`D`$ is irreducible. Let $`q:AA/A^{}`$ denote the quotient map. Since $`q`$ is a proper map, $`q(D)`$ is a projective subvariety in $`A/A^{}`$. Since $`D`$ is irreducible and $`TD=\mathrm{}`$, $`q(D)`$ is an irreducible divisor. So $`\stackrel{~}{D}=q^1(q(D))A`$ is also an irreducible divisor containing $`D`$ as $`q`$ is smooth. This forces $`D=\stackrel{~}{D}`$. $``$$``$
Remark Lemma 5.6 is false for semi-abelian varieties. For we may take $`G=(𝐂^{})^2`$, $`T=\{(z_1,z_2)G:z_1=1\}`$, $`D=\{(z_1,z_2)G:(z_11)z_2=1\}.`$
### 5.3 Jet bundles on semi-abelian varieties
To simplify notation, we work exclusively with a semi-abelian variety $`G`$ and its associated log-manifold $`(\overline{G},S)`$ defined as before. We remark, however, that many the definitions and results hold for an arbitray log-manifold.
Recall from subsections 1.2 and 2.1 that $`\overline{P}_kV`$ denotes the logarithmic $`k`$-jet bundle of $`(\overline{G},S,\overline{V})`$ and that $`\overline{P}_kG`$ denotes the logarithmic jet bundle of $`(\overline{G},S)=(\overline{G},S,\overline{T}G)`$. Note that the log-directed morphism $`i:(\overline{G},S,\overline{V})(\overline{G},S,\overline{T}G)`$ induces a canonical realization of $`\overline{P}_kV`$ as a submanifold of $`\overline{P}_kG`$ as $`\overline{V}`$ is a subbundle of $`\overline{T}G`$ over $`\overline{G}`$.
Let $`DG`$ be a reduced algebraic divisor. Then, by Hironaka (), there exist a log-manifold $`(\overline{Y},E)`$ and a log-morphism $`p:(\overline{Y},E)(\overline{G},S)`$ with
* $`p^1(SD)=E,`$
* $`p:p^1(\overline{G}D)\overline{G}D`$ is biholomorphic.
Given a subbundle $`\overline{V}`$ of $`\overline{T}G`$, we have the following commutative diagram:
$$\begin{array}{ccccc}𝒪_{\overline{P}_kY}(1)& \stackrel{(p_{[k1]})_{}}{}& 𝒪_{\overline{P}_kG}(1)& \stackrel{(i_{[k1]})_{}}{}& 𝒪_{\overline{P}_k(V)}(1)\\ & & & & \\ & & & & \\ & & & & \\ \overline{P}_kY& \stackrel{p_{[k]}}{}& \overline{P}_kG& \stackrel{i_{[k]}}{}& \overline{P}_k(V)\\ & & & & \\ & & & & \\ & & & & \\ (\overline{Y},E)& \stackrel{p}{}& (\overline{G},S)& \stackrel{i}{}& (\overline{G},S,\overline{V})\end{array}$$
(1)
Here, $`i_{[k]}`$ realizes $`\overline{P}_kV`$ as a submanifold of $`\overline{P}_kG`$ by Proposition 2.1. Outside $`\pi _{0,k}^1(D)\overline{P}_kG`$, the maps $`p_{[k]}`$ (and hence, $`p_{[k1]}^{}{}_{}{}^{}`$) are isomorphisms. All other maps are holomorphic. We define
$$\overline{Z}_k:=\overline{p_{[k]}^1(\overline{P}_kV\pi _{0,k}^1(D))}^{\mathrm{Zariski}}\overline{P}_kY.$$
###### Definition 5.7
A meromorphic section $`s`$ of $`𝒪_{\overline{P}_kV}(m)`$ is said to have at most log-poles along $`D`$ if it pulls back, via the map $`(p_{[k1]})_{}|_{\overline{Z}_k}=(p_{[k1]}|_{\overline{Z}_k})_{}`$, to a holomorphic section<sup>2</sup><sup>2</sup>2With this we mean that, after pulling back the section $`s`$ over the part of $`\overline{Z}_k`$ where the meromorphic map $`(p_{[k1]})_{}`$ is holomorphic, it extends to a holomorphic section of $`𝒪_{\overline{P}_kY}(m)|_{\overline{Z}_k}`$. of $`𝒪_{\overline{P}_kY}(m)|_{\overline{Z}_k}`$.
Suppose that a meromorphic section $`s`$ of $`𝒪_{\overline{P}_kV}(m)`$ over an open subset $`U\overline{G}`$ is defined by a meromorphic function $`Q`$ on $`J_kG|_U`$ satisfying equation (2), see Proposition 3.9. Suppose more precisely that $`Q`$ is given by a polynomial in the differentials up to order $`k1`$ of sections of $`\overline{\mathrm{\Omega }}G|_U`$ as well as the differentials up to order $`k1`$ of $`d\mathrm{log}\theta `$, where $`\theta `$ is a meromorphic function on $`U`$ nonvanishing and holomorphic on $`U_0=U\{DS\}`$. Then, after composing with $`p`$ given above, these differentials, and so also the polynomial in them, become holomorphic functions on $`\overline{J}_kY|_{p^1(U)}`$ by Proposition 1.1(c). Furthermore, the resulting polynomial still satisfies equation (2) on $`p^1(U)`$. Hence, by Proposition 3.9, we obtain a holomorphic section of $`𝒪_{\overline{P}_kY}(m)|_{\stackrel{~}{U}}`$ that matches with the pullback of $`s`$ on an open set of $`\overline{Z}_k`$. Therefore, $`s|_U`$ is meromorphic with at most log-poles along $`D`$. If such a description is possible on a neighborhood $`U`$ of each point in $`\overline{G}`$, then $`s`$ is meromorphic with at most log-poles along $`D`$. We will consider examples of such an $`s`$ in subsection 5.5.
###### Lemma 5.8
Let $`(\overline{G},S,\overline{V})`$ be as above.
* There exist injective maps
$$\overline{P}_{k+l}V\overline{P}_l(P_kV)\text{and }𝒪_{\overline{P}_{k+l}V}(1)𝒪_{\overline{P}_l(P_kV)}(1)$$
which are given outside $`S`$ by $`f_{[k+l]}(f_{[k]})_{[l]}`$ and by $`f_{[k+l1]}^{}(f_{[k]})_{[l1]}^{}`$, respectively, and which realize $`\overline{P}_{k+l}V\overline{P}_l(P_kV)`$ and $`𝒪_{\overline{P}_{k+l}V}(1)𝒪_{\overline{P}_l(P_kV)}(1)`$, respectively, as submanifolds.
* Furthermore, let $`\mathrm{\Gamma }`$ be a curve and $`f:\mathrm{\Gamma }X`$ be a holomorphic map which is tangent to $`V`$. As before, let $`\overline{X}_k(f)\overline{P}_kV`$ be the Zariski closure of the image of the $`k`$-th lift $`f_{[k]}:\mathrm{\Gamma }\overline{P}_kV`$ of the map $`f`$. We denote $`\pi _{0,k}^1(S)`$ again by $`S`$ and $`\pi _{0,k}^1(G)\overline{X}_k(f)`$ by $`X_k(f)`$. We recall that by Hironaka there exists a log-morphism
$$\mathrm{\Psi }:(\overline{\stackrel{~}{X}_k}(f),\stackrel{~}{S})(\overline{P}_kV,S)\text{such that:}$$
+ $`\mathrm{\Psi }(\overline{\stackrel{~}{X_k}}(f))=\overline{X}_k(f)`$.
+ $`\mathrm{\Psi }^1(S)=\stackrel{~}{S}`$.
+ $`\mathrm{\Psi }`$ is biholomorphic outside $`\mathrm{\Psi }^1\left(\mathrm{Sing}(\overline{X}_k(f))\right)`$.
We set $`\stackrel{~}{X}_k(f)=\overline{\stackrel{~}{X}_k}(f)\stackrel{~}{S}`$. With this setup, we have the following commutative diagram:
$$\begin{array}{ccccc}\overline{P}_l(\stackrel{~}{X}_k(f))& \stackrel{\mathrm{\Psi }_{[l]}}{}& \overline{P}_l(P_kV)& & \\ & & & & \\ |& & i& & \\ & & & & \\ |& \overline{X}_{k+l}(f)& \overline{P}_{k+l}(V)& & \\ & & & & \\ & & & & \\ & & & & \\ \overline{\stackrel{~}{X}_k}(f)\stackrel{\mathrm{\Psi }}{}& \overline{X}_k(f)& \overline{P}_k(V)& & \end{array}$$
(2)
where $`\mathrm{\Psi }_{[l]}`$ may be meromorphic, all other maps in the diagram are holomorphic and the following holds:
$$i(\overline{X}_{k+l}(f))\mathrm{\Psi }_{[l]}(\overline{P}_l(\stackrel{~}{X}_k(f))).$$
* Let $`s`$, $`t`$ be meromorphic sections of the bundle $`𝒪_{\overline{P}_kV}(m)`$ with at most log-poles along $`D`$, and assume $`t`$ is not the zero section. Then $`t^2d(\frac{s}{t})`$ can be considered as a meromorphic section of the bundle $`𝒪_{\overline{P}_{k+1}V}(2m+1)`$ with at most log-poles along $`D`$.
Proof for (1) This follows directly from Corollary 2.4 and Proposition 2.1 applied to the subbundle inclusion $`\overline{V}_k\overline{T}P_{k1}V`$. The presentation of the maps outside $`S`$ follows from Proposition 3.1.
Proof for (2) $`\mathrm{\Psi }_{[l]}:\overline{P}_l(\stackrel{~}{X}_k(f))\overline{P}_l(P_kV)`$ is a proper rational map. Hence, $`\mathrm{\Psi }_{[l]}(\overline{P}_l(\stackrel{~}{X}_k(f))`$ is an algebraic subset containing $`i(f_{[k+l]}(\mathrm{\Gamma }))=(f_{[k]})_{[l]}(\mathrm{\Gamma })`$. Therefore, we also have $`i(\overline{X}_{k+l}(f))`$ $`\mathrm{\Psi }_{[l]}(\overline{P}_l(\stackrel{~}{X_k}(f))`$.
Proof for (3) $`\frac{s}{t}`$ is a rational function on $`\overline{P}_kV`$. Hence, $`d(\frac{s}{t})`$ is a rational section of $`𝒪_{\overline{P}_1(P_kV)}(1)`$, which lifts back, via the inclusion map $`𝒪_{\overline{P}_{k+1}V}(1)`$ $`𝒪_{\overline{P}_1(P_kV)}(1)`$, to a rational section of $`𝒪_{\overline{P}_{k+1}V}(1)`$. Hence, $`t^2d(\frac{s}{t})`$ is a rational section of $`𝒪_{\overline{P}_{k+1}V}(2m+1)`$. To prove that $`t^2d(\frac{s}{t})`$ again is a meromorphic section with at most log-poles along $`D`$, we pull back the sections $`s`$, $`t`$ and $`t^2d(\frac{s}{t})`$ to $`\overline{Z}_k`$ by the map $`(p_{[k1]})_{}`$. Then by definition the sections $`s`$ and $`t`$ become holomorphic sections on $`\overline{Z}_k`$. It suffices to show that the rational section $`t^2d(\frac{s}{t})`$ also is holomorphic on $`\overline{Z}_k`$. It suffices to prove this locally on $`\overline{Z}_k`$. Given any point $`w\overline{Z}_k\overline{P}_kY`$ there exists an open neighborhood $`U`$ of $`w`$ in $`\overline{P}_kY`$ such that the holomorphic sections $`s`$ and $`t`$ over $`\overline{Z}_kU`$ extend to holomorphic sections of the bundle $`𝒪_{\overline{P}_kY}(m)`$ over $`U`$, and that this bundle is trivial over $`U`$. After choosing such a trivialization, one has, by the product rule for the holomorphic functions $`s`$ and $`t`$,
$$t^2d(\frac{s}{t})=tdssdt,$$
the latter being a holomorphic section of $`𝒪_{\overline{P}_1(UE)}(1)`$. $``$$``$
Next, we want to prove that $`\overline{P}_kG`$, and even $`\overline{P}_kV`$, are trivial over $`\overline{G}`$ for certain subbundles $`\overline{V}\overline{T}G`$, which we will call special. Let $`z_1`$, …, $`z_n`$ be any linear coordinates of the universal cover $`𝐂^nG`$. We have observed that $`\overline{T}G=𝐂^n\times \overline{G}`$, where the trivialization is given by $`dz_1`$, … $`dz_n`$.
###### Definition 5.9
$`\overline{V}`$ is said to be special if $`\overline{V}=𝐂^r\times \overline{G}`$ in this trivialization.
From now on, all subbundles $`\overline{V}\overline{T}G`$ we use are assumed to be special.
###### Lemma 5.10
* The map
$$𝒪_{P_kV}(1)(𝒪_{P_kV}(1))_0\times G;f_{[k1]}^{}(0)((ff(0))_{[k1]}^{}(0),f(0))$$
(3)
gives an isomorphism $`𝒪_{P_kV}(1)𝒪_{\overline{P}_kV}(1)|_G`$, and this isomorphism is invariant under the action of $`G`$.
* This isomorphism can be extended to a $`G`$-invariant isomorphism
$$\overline{\mathrm{\Psi }}_k:𝒪_{\overline{P}_kV}(1)(𝒪_{P_kV}(1))_0\times \overline{G}$$
respecting the fibers of the line bundles.
* By combining with the canonical line bundle projections we get the same isomorphisms with $`P_kV`$ and $`\overline{P}_kV`$ instead of $`𝒪_{P_kV}(1)`$ and $`𝒪_{\overline{P}_kV}(1)`$.
Proof (a) is immediate, and (c) follows immediately from (b). To prove (b), we use that by Corollary 3.7 (a) the trivialization of (a) extends locally, so by the Identity Theorem it extends globally. The invariance under the action of $`G`$ of this trivialization extends from $`G`$ to $`\overline{G}`$ by continuity.
We would like, however, also to indicate a direct proof. It is obtained by proving the following more precise statement by induction over $`k`$.
Claim S(k) Then there exist trivialization maps $`\overline{\psi }_k`$, induced canonically by the trivialization of $`V`$ by $`dz_1,\mathrm{},dz_n`$, such that
$$\begin{array}{cccccc}& \overline{V}_k& \stackrel{\overline{\psi }_k}{}& (V_k)_0\times \overline{G}& & \\ & & & & & \\ (+)_k& \pi _k& & (\pi _k)_0\times id_{\overline{G}}& & \\ & & & & & \\ & \overline{P}_{k1}V& \stackrel{\overline{\mathrm{\Psi }}_{k1}}{}& (P_{k1}V)_0\times \overline{G}& & \end{array}$$
commutes and the upper line projectivizes to
$$\begin{array}{cccccc}& \overline{P}_kV& \stackrel{\overline{\mathrm{\Psi }}_k}{}& (P_kV)_0\times \overline{G}& & \\ & & & & & \\ (++)_k& \pi _k& & (\pi _k)_0\times id_{\overline{G}}& & \\ & & & & & \\ & \overline{P}_{k1}V& \stackrel{\overline{\mathrm{\Psi }}_{k1}}{}& (P_{k1}V)_0\times \overline{G}& & \end{array}$$
where $`\mathrm{\Psi }_k`$ extends the isomorphism in equation (3) in a $`G`$-invariant way.
Now S(1) is clear from the trivialization
$$V\stackrel{~}{}𝐂^r\times \overline{G},$$
since we are given that $`V`$ is special. Assuming, by induction, that S(k) is true. Then we get $`(++)_{k+1}`$ by projectivizing $`(+)_{k+1}`$. It remains to show $`(+)_{k+1}`$. By $`(++)_k`$ we get induced trivializations of the logarithmic tangent bundles
$$\begin{array}{ccccc}\overline{T}(P_kV)& \stackrel{(\overline{\mathrm{\Psi }}_k)_{}}{}& T((P_kV)_0)\times \overline{T}G& & T((P_kV)_0)\times 𝐂^r\times G\\ & & & & \\ (\pi _k)_{}& & ((\pi _k)_0)_{}\times (id_{\overline{G}})_{}& & \\ & & & & \\ \overline{T}(P_{k1}V)& \stackrel{(\overline{\mathrm{\Psi }}_{k1})_{}}{}& T((P_{k1}V)_0)\times \overline{T}G& & T((P_{k1}V)_0)\times 𝐂^r\times G\end{array}$$
where the isomorphisms on the right hand side are obtained by trivializing $`\overline{T}G`$ by the forms $`dz_1,\mathrm{},dz_n`$. We want to show it we restrict the isomorphism in the upper line of this diagram from $`\overline{T}(P_kV)`$ to $`\overline{V}_k`$, we also get a trivialization of $`\overline{V}_k`$ over $`\overline{G}`$. Then we can denote this trivialization by $`\psi _{k+1}`$ and the rest follows easily. The key point of the proof is now that by equation (1), namely
$$\overline{V}_k(G):=(\pi _k)_{}^1(𝒪_{\overline{P}_kV}(1))\overline{T}(P_kV),$$
the subbundle $`\overline{V}_k\overline{T}P_kV`$ is defined in an intrinsic way which is compatible with the isomorphisms of the diagram above. $``$$``$
###### Lemma 5.11
Let $`Y`$ be a complex manifold, $`ZY`$ a complex submanifold and denote by $`i:(Z,TZ)(Y,TY)`$ the directed inclusion map. Let $`g:\mathrm{\Delta }Y`$ be holomorphic with $`dg(0)0`$, where $`\mathrm{\Delta }`$ denotes again the unit disk in $`𝐂`$. Assume that $`g_{[l]}(0)`$ is in the image of the (composed) morphism
$$P_lZ\stackrel{i_{[l]}}{}i^1P_lYP_lY$$
(4)
for all $`l0`$. Then $`g(\mathrm{\Delta })Z`$.
Proof By Proposition 2.1 a) and b) the map in equation (4) is a morphism. Let $`UY`$ be a neighborhood of $`g(0)`$ and $`F:U𝐂`$ be a holomorphic function with $`F|_{ZU}0`$. It suffices to show $`Fg0`$. There exist a small disk $`\mathrm{\Delta }_ϵ`$ and a map $`h_l:\mathrm{\Delta }_ϵZ`$ with $`i_{[l]}((h_l)_{[l]}(0))=g_{[l]}(0)`$ and by Corollary 2.3 we may assume $`dh_l(0)0`$. By Proposition 3.1 we have $`i_{[l]}(h_l)_{[l]}=(ih_l)_{[l]}`$ and hence, $`(ih_l)_{[l]}(0)=g_{[l]}(0)`$ and $`d(ih_l)(0)0`$. Hence, we can reparametrize $`ih_l`$ in a way that it has the same Taylor expansion as $`g`$ up to order $`l`$. We assume that this has been done, and note that $`h_l`$ still maps a neighborhood of the origin to $`Z`$. Hence,
$$(\frac{^j}{t^j}Fg)(0)=(\frac{^j}{t^j}Fih_l)(0)=0\mathrm{for}jl.$$
Since $`l`$ is arbitrary, we get $`Fg0`$. $``$$``$
### 5.4 The Main Lemma
The following Main Lemma is the key step in proving Theorem 5.1 and Theorem 5.3. In the case $`\mathrm{\Gamma }=𝐂`$, it is a generalization of a lemma contained in , and for $`G=A`$, it generalizes a lemma in .
###### Main Lemma 5.12
With the same setup as that for Theorem 5.1 (or Theorem 5.3), let $`\overline{V}\overline{T}G`$ be a special subbundle. Assume that $`f:\mathrm{\Gamma }GD`$ is tangent to $`V`$ and, in the case $`\mathrm{\Gamma }=\mathrm{\Delta }^{}`$, that $`f`$ does not extend to $`\mathrm{\Delta }`$ as a map to $`\overline{G}`$. Let, for $`k0`$, $`\overline{X}_k(f)`$ denote the Zariski closure of $`f_{[k]}(\mathrm{\Gamma })`$ in $`\overline{P}_kV`$. Let, for $`k,m1`$, $`\mathrm{\Theta }`$ be meromorphic section of the line bundle $`𝒪_{\overline{P}_kV}(m)`$ with at most log-poles along $`D`$, there exists an algebraic subgroup $`G^{}G`$, of positive dimension, which leaves $`\overline{X}_k(f)`$ and, for $`k1`$, also $`\mathrm{\Theta }|_{\overline{X}_k(f)}`$ invariant.
Remark The same is true for finitely many different sections $`\mathrm{\Theta }`$.
The rest of this subsection will be devoted to the proof of the Main Lemma. It suffices to consider the case $`k1`$. In fact, to prove the case $`k=0`$ we apply the Main Lemma for $`k=1`$ and for $`\mathrm{\Theta }`$ being the zero section in $`𝒪_{\overline{P}_1V}(1)`$. Since the map $`\pi _1:\overline{P}_1V\overline{X}`$ is equivariant under the action of $`G^{}`$ and maps $`\overline{X}_1(f)`$ surjectively $`\overline{X}_0(f)`$, the subgroup $`G^{}`$ also leaves $`\overline{X}_0(f)`$ invariant. So for the rest of the proof of the Main Lemma we assume $`k1`$.
We fix $`u\mathrm{\Gamma }`$ to be any point for which $`df(u)0`$. Then all $`f_{[k+l]}(u)`$, $`l0`$, are regular jets. Let $`s_0H^0(\overline{P}_1V,𝒪_{\overline{P}_1V}(1))`$ be a global section, which is invariant under the action of $`G`$, and which is nonvanishing at $`f_{[1]}(u)`$. It exists because $`df(u)0`$ and $`𝒪_{\overline{P}_1V}(1)=𝒪_{P_1V_0}(1)\times \overline{G}`$ (see Lemma 5.10). Choose an infinite sequence $`\{n_0,n_1,n_2,n_3,\mathrm{}\}`$ of natural numbers such that the following two conditions hold:
$$(2(k+l)1)|n_l\mathrm{for}l0,$$
$$n_l2(n_{l1}+1)\mathrm{for}l1.$$
For example, $`n_l=2^l(2(k+l)1)m`$, where $`m=\mathrm{deg}\mathrm{\Theta }`$ will work. Let
$$\mathrm{\Theta }_0=\mathrm{\Theta }(s_0^{n_0m})$$
and, for $`l1`$, define inductively:
$$\mathrm{\Theta }_l=d(\frac{\mathrm{\Theta }_{l1}}{s_0^{n_{l1}}})(s_0^{n_l1}).$$
Then, by Lemma 5.8 (3), $`\mathrm{\Theta }_l`$ is a meromorphic section of $`𝒪_{\overline{P}_{k+l}V}(n_l)`$ with at most log-poles along $`D`$ (here $`s_0=s_0(\pi _{0,k+l1})_{}`$).
By Lemma 5.10 we may identify $`\overline{P}_kV`$ as $`P_kV_0\times \overline{G}`$. Then we have:
$$\begin{array}{ccccc}𝒪_{\overline{P}_{k+l}(V)}(1)|_{\overline{X}_{k+l}(f)}& & 𝒪_{\overline{P}_{k+l}(V)}(1)& \hfill \stackrel{\alpha }{}& 𝒪_{P_{k+l}V_0}(1)\hfill \\ & & & & \\ & & & & \hfill \\ & & & & \\ \overline{X}_{k+l}(f)& & \overline{P}_{k+l}(V)& \hfill =P_{k+l}V_0\times \overline{G}& \stackrel{p_1}{}P_{k+l}V_0\hfill \\ & & & & \\ & & & \hfill & \hfill \\ & & & & \\ \overline{X}_k(f)& & \overline{P}_k(V)& \hfill =P_kV_0\times \overline{G}& \stackrel{p_1}{}P_kV_0\hfill \\ & & & & \\ & & & \hfill & \hfill \\ & & & & \\ \mathrm{\Gamma }\stackrel{f}{}\overline{X}(f)& & \overline{G}& \hfill =\overline{G}& 0\hfill \end{array}$$
where $`0GD`$ and $`p_1`$ is projection to the first factor. Define, for $`l0`$:
$$W_l=\{aG:(f+a)_{[k+l]}(u)\overline{X}_{k+l}(f)\mathrm{and}\frac{\mathrm{\Theta }_i}{s_0^{n_i}}|_{f_{[k+i]}(u)}=\frac{\mathrm{\Theta }_i}{s_0^{n_i}}|_{(f+a)_{[k+i]}(u)},i=0,\mathrm{},l\}.$$
###### Lemma 5.13
With the hypothesis and the setup as in the Main Lemma, $`W:=_{l=0}^{\mathrm{}}W_l`$ is an algebraic subvariety of $`G`$ and $`\mathrm{dim}_0W1`$.
Proof of Lemma 5.13 $`W_l`$ is an algebraic subvariety of $`G`$ as the group action of $`G`$ on itself is algebraic. Hence, $`W`$ is also algebraic. Let $`l_1>l_2`$ and $`\pi _{k+l_2,k+l_1}:\overline{P}_{k+l_1}V\overline{P}_{k+l_2}V`$. If $`(f+a)_{[k+l_1]}(u)\overline{X}_{k+l_1}(f)`$, then
$$(f+a)_{[k+l_2]}(u)=\pi _{k+l_2,k+l_1}(f+a)_{[k+l_1]}(u)\pi _{k+l_2,k+l_1}(\overline{X}_{k+l_1}(f))=\overline{X}_{k+l_2}(f).$$
Hence, $`W_l`$, $`l0`$, is a decreasing sequence of algebraic subvarieties of $`G`$. So the proof of Lemma 5.13 is complete if we show:
$$\mathrm{dim}_0W_l1\mathrm{for}l0.$$
By the beginning of part iii) of the proof of Theorem 6.8 of Demailly in , the rational map, obtained by a basis of holomorphic sections of the line bundle $`𝒪_{P_{k+l}V_0}(2(k+l)1)`$, is a morphism on the subset $`P_{k+l}(V)_0^{\mathrm{reg}}`$ of regular jets in $`P_{k+l}V_0`$ and seperates all points there. Denote by $`L_{(2(k+l)1)}`$ the linear system obtained by the pull backs of these sections by the map $`\alpha `$ (see the last diagram). Then we have
$$(L_{(2(k+l)1)})^{\frac{n_l}{2(k+l)1}}H^0(\overline{P}_{k+l}V,𝒪_{\overline{P}_{k+l}V}(n_l)).$$
Therefore, the sections of $`H^0(\overline{P}_{k+l}V,𝒪_{\overline{P}_{k+l}V}(n_l))`$ still separate points in the subset of regular jets of each fiber of the map $`\overline{P}_{k+l}V\overline{G}`$. Let the map $`\mathrm{\Phi }_l:\overline{P}_{k+l}V𝐏^{N_l}`$ be defined by a basis of these sections. Then the fiber of the map $`\mathrm{\Phi }_l`$ through a regular jet $`\xi \overline{P}_{k+l}V`$ is necessarily of the form $`\{\xi +a,aR\}`$, where $`R\overline{G}`$ is algebraic.
To the basis of holomorphic sections which define the map $`\mathrm{\Phi }_l`$, we now add some extra sections which we allow in addition to have log-poles along the divisor $`D`$, namely the sections $`\mathrm{\Theta }_is_0^{n_ln_i}`$, $`i=0,\mathrm{},l`$. So we get a map
$$\stackrel{~}{\mathrm{\Phi }_l}:\overline{P}_{k+l}V𝐏^{N_l+l+1}.$$
This map will in general only separate the subset of regular jets of those fibers of the map $`\pi _{0,k+l}:\overline{P}_{k+l}V\overline{G}`$ which are not over the divisor $`D`$. But the fibers of the map $`\stackrel{~}{\mathrm{\Phi }}_l:\overline{P}_{k+l}V𝐏^{N_l+l+1}`$ through a regular jet $`\xi \overline{P}_{k+l}V\pi _{0,k+l}^1(D)`$ must still be of the form $`\{\xi +a,aR_\xi \}`$, where $`R_\xi R\overline{G}`$ is an algebraic subset. So, Lemma 5.13 is proved if we show that $`\stackrel{~}{\mathrm{\Phi }_l}:\overline{X}_{k+l}(f)𝐏^{N_l+l+1}`$ has positive dimensional fiber through $`\xi =f_{[k+l]}(u)`$.
For proving this, we want to use our Ahlfors Lemma 4.3. But this only applies for holomorphic sections. We first extend the diagram in equation (1) to
$$\begin{array}{cccccccc}\hfill \overline{Z}_{k+l}(f)& & \overline{P}_{k+l}Y& \stackrel{p_{[k+l]}}{}\overline{P}_{k+l}G& & \overline{P}_{k+l}(V)& & X_{k+l}(f)\hfill \\ & & & & & & & \\ & & & & & & & \\ & & & & & & & \\ \hfill \overline{Z}_k(f)& & \overline{P}_kY& \stackrel{p_{[k]}}{}\overline{P}_kG& & \overline{P}_k(V)& & X_k(f)\hfill \\ & & & & & & & \\ & & & & & & & \\ & & & & & & & \\ & & (\overline{Y},E)& \stackrel{p}{}(\overline{G},S)& & (\overline{G},S,\overline{V})& & \end{array}$$
where
$$\overline{Z}_k(f)=\overline{p_{[k]}^1(X_k(f)\pi _{0,k}^1(D))}^{\mathrm{Zariski}}\overline{Z}_k\overline{P}_kY,$$
$$\overline{Z}_{k+l}(f)=\overline{p_{[k+l]}^1(X_{k+l}(f)\pi _{0,k+l}^1(D))}^{\mathrm{Zariski}}\overline{Z}_{k+l}\overline{P}_{k+l}Y.$$
By functoriality of the jet bundles and definition 5.7, the sections which define $`\stackrel{~}{\mathrm{\Phi }_l}`$ pull back to holomorphic sections to span a linear system
$$\stackrel{~}{L}_{k+l}H^0(\overline{Z}_{k+l}(f),𝒪_{\overline{P}_{k+l}Y}(n_l)).$$
The elements of $`\stackrel{~}{L}_{k+l}`$ define the pull back of $`\stackrel{~}{\mathrm{\Phi }_l}`$ to $`\overline{Z}_{k+l(f)}`$ (which we denote again by $`\stackrel{~}{\mathrm{\Phi }_l}`$).
So we can apply our Ahlfors Lemma 4.3 to the map
$$\stackrel{~}{\mathrm{\Phi }}_l:\overline{Z}_{k+l}(f)𝐏^{N_l+l+1}$$
to conclude that $`\stackrel{~}{f}_{[k+l]}(u)B_{k+l}(\overline{Z}_{k+l}(f))`$, where $`\stackrel{~}{f}=p^1f`$ and, without loss of generality, $`u\mathrm{Sing}\stackrel{~}{X}_{k+l}`$. For otherwise, the map $`\stackrel{~}{f}`$ would be constant or (in the case of $`\mathrm{\Gamma }=\mathrm{\Delta }^{}`$) at least extendable, which would imply that $`f`$ has this property, or the image of $`\stackrel{~}{f}_{[k+l]}`$ would be contained in $`\mathrm{Sing}\overline{Z}_{k+l}(f)`$, which is impossible, since $`\overline{Z}_{k+l}(f)`$ is the proper transform of the Zariski closure of the image of $`f_{[k+l]}`$. Hence, $`\stackrel{~}{\mathrm{\Phi }}_l:\overline{Z}_{k+l}(f)𝐏^{N_l+l+1}`$ has positive dimensional fiber through $`\stackrel{~}{f}_{[k+l]}(u)`$. Since $`f_{[k+l]}(u)\pi _{0,k+l}^1(SD)`$, the map $`p_{[k+l]}^1`$ is an isomorphism around $`f_{[k+l]}(0)`$. Hence, $`\stackrel{~}{\mathrm{\Phi }_l}:X_{k+l}(f)𝐏^{N_l+l+1}`$ also has positive dimensional fiber. This ends the proof of Lemma 5.13. $``$$``$
Let us now continue with the proof of the Main Lemma. Without loss of generality we may assume that $`f_{[k]}(u)\mathrm{Sing}(X_k(f))`$. Then there exists a neighborhood $`U=U(0)G`$ such that, for all $`aU`$, we have $`(f+a)_{[k]}(u)\mathrm{Sing}(X_k(f))`$. By Lemma 5.8 (2), we get
$$((f+a)_{[k]})_{[l]}(u)\overline{P}_l(\stackrel{~}{X_k}(f))\overline{P}_l(P_k(V))$$
for $`aWU`$ and all $`l0`$, where we have omitted to write the map $`\mathrm{\Psi }`$. This is justified by the fact that around $`(f+a)_{[k]}(u)`$, for $`aU`$, the variety $`X_k(f)`$ is smooth, and so $`\mathrm{\Psi }`$ is an isomorphism there. Applying Lemma 5.11, we get that $`(f+a)_{[k]}(\mathrm{\Gamma })\stackrel{~}{X_k}(f)`$. Hence,
$$(f+a)_{[k]}(\mathrm{\Gamma })X_k(f)$$
for $`aWU`$. But this means
$$f(\mathrm{\Gamma })X_k(f)(X_k(f)+a).$$
Since $`X_k(f)`$ was the Zariski closure of $`f(\mathrm{\Gamma })`$, we get
$$X_k(f)=X_k(f)+a$$
for all $`aWU`$. We next want to show:
###### Lemma 5.14
$$\frac{\mathrm{\Theta }}{s_0^{\mathrm{deg}\mathrm{\Theta }}}|_{X_k(f)}$$
is invariant under the action of all $`aWU`$
Proof of Lemma 5.14 Let $`aWU`$ be fixed. In order to simplify notation, we denote by $`F_{(i)}`$, $`i\mathrm{I}\mathrm{N}_0`$, the following rational function on $`\overline{P}_{k+i}V`$:
$$F_{(i)}(y):=\frac{\mathrm{\Theta }_i}{s_0^{n_i}}(y)\frac{\mathrm{\Theta }_i}{s_0^{n_i}}(y+a).$$
Set $`F=F_{(0)}`$. It suffices to show that, for all $`i\mathrm{I}\mathrm{N}_0`$, we have
$$\frac{^i}{t^i}(Ff_{[k]})(u)=0.$$
(5)
For then, by analytic continuation applied to $`f_{[k]}:\mathrm{\Delta }\overline{P}_kV`$, we will have $`Ff_{[k]}(\mathrm{\Gamma })0`$, so that $`F0`$ on $`X_k(f)`$ as required by Lemma 5.14.
By abuse of notation, we have identified $`s_0(\pi _{0,k+l1})_{}`$ with $`s_0`$. Since these sections are maps from $`𝒪_{\overline{P}_{k+l}V}(1)`$ respectively $`𝒪_{\overline{P}_1V}(1)`$, this means we have
$$s_0(f_{[k+l]}(t))f_{[k+l1]}^{}(t)=s_0(f_{[1]}(t))f^{}(t).$$
(6)
Recall that the section $`s_0`$ is nonvanishing at $`f_{[1]}(u)`$, and that $`f^{}(u)0`$. So, after possibly shrinking $`\mathrm{\Delta }`$, we may reparametrize $`f`$ such that
$$s_0(f_{[k+l]}(t))f_{[k+l1]}^{}(t)=s_0(f_{[1]}(t))f^{}(t)1.$$
(7)
For the rest of this proof we fix the parameter $`t`$ in such a way that equation (7) is satisfied. We claim that for any $`i\mathrm{I}\mathrm{N}_0`$ we have:
$$\frac{^i}{t^i}(Ff_{[k]})(t)=F_{(i)}f_{[k+i]}(t).$$
(8)
We prove this by induction over $`i\mathrm{I}\mathrm{N}_0`$. The case S(0) is clear by definition. Assume that S($`i`$), $`i<l`$ is true. Then we have
$`{\displaystyle \frac{^l}{t^l}}(Ff_{[k]})(t)`$ $`=`$ $`{\displaystyle \frac{}{t}}({\displaystyle \frac{^{l1}}{t^{l1}}}(Ff_{[k]}))(t)`$
$`=`$ $`{\displaystyle \frac{}{t}}(F_{(l1)}f_{[k+l1]})(t)`$
$`=`$ $`(dF_{(l1)})((f_{[k+l1]})(t))f_{[k+l1]}^{}(t)`$
$`=`$ $`{\displaystyle \frac{(dF_{(l1)})((f_{[k+l1]})(t))f_{[k+l1]}^{}(t)}{s_0((f_{[k+l]})(t))f_{[k+l1]}^{}(t)}}`$
$`=`$ $`{\displaystyle \frac{(dF_{(l1)})((f_{[k+l]})(t))f_{[k+l1]}^{}(t)}{s_0((f_{[k+l]})(t))f_{[k+l1]}^{}(t)}}`$
$`=`$ $`{\displaystyle \frac{dF_{(l1)}}{s_0}}((f_{[k+l]})(t))=F_{(l)}((f_{[k+l]}(t)).`$
Here, we regard the differentials as linear maps on $`𝒪(1)`$, more specifically, sections of $`𝒪(1)`$. To see the second equality from below, recall that the section $`dF_{(l1)}`$ of $`𝒪_{\overline{P}_1(P_{k+l1}(V))}(1)`$ naturally restricts to a section of $`𝒪_{\overline{P}_{k+l}V}(1)`$. This proves equation (8). But from equation (8), equation (5) follows immediately. This is because, by the definition of $`W`$, $`F_{(i)}f_{[k+i]}(u)=0`$ for all $`i\mathrm{I}\mathrm{N}_0`$. $``$$``$
Now the proof of the Main Lemma is immediate. Let us define
$$\stackrel{~}{W}=\{aG:X_k(f)=X_k(f)+a,\frac{\mathrm{\Theta }}{s_0^{\mathrm{deg}\mathrm{\Theta }}}\mathrm{is}\mathrm{invariant}\mathrm{under}a\}.$$
$`\stackrel{~}{W}`$ is clearly a group, which is algebraic by Lemma 5.5. It is of positive dimension, since $`(WU)\stackrel{~}{W}`$ and $`\mathrm{dim}_0(WU)1`$. This finishes the proof of the Main Lemma 5.12. $``$$``$
### 5.5 Proof of Theorems 5.1 and 5.3
Proof of Theorems 5.1 (a) and 5.3 (a) We apply the Main Lemma 5.12 in the special case where $`k=0`$ and $`V=\overline{T}G`$. Then Theorem 5.3 (a) is immediate. Theorem 5.1 (a) is obtained by dividing out by the biggest algebraic subgroup of $`G`$ under which $`X(f)`$ is invariant. $``$$``$
In order to prove the remaining parts of Theorems 5.1 and 5.3, we first have to choose the section $`\mathrm{\Theta }`$ appropriately. We do this in the same way as Siu-Yeung () or Noguchi ().
By Noguchi (, Lemma 2.1), there exists a theta function for $`DG`$. This means the following. Let $`\pi :𝐂^nG`$ be the universal covering with a ‘semi-lattice’ $`\mathrm{\Pi }_1(G)`$. There exists an entire function $`\theta :𝐂^n𝐂`$ such that
$$(\theta )=\pi ^{}D.$$
Moreover, for any $`\gamma \mathrm{\Pi }_1(G)`$, there is an affine linear function $`L_\gamma `$ in $`x`$ with
$$\theta (x+\gamma )=e^{L_\gamma (x)}\theta (x),x𝐂^n.$$
(9)
But the proof which Noguchi gives actually yields more. Consider $`G`$ as a $`(𝐂^{})^{\mathrm{}}`$-principal fiber bundle over an abelian variety $`A`$, and denote the projection map by $`p:GA`$. Let $`\pi :𝐂^𝗆A`$ be the universal covering. Then the fibered products of $`\pi `$ with $`G`$ respectively $`\overline{G}`$ over $`A`$ are $`(𝐂^{})^{\mathrm{}}\times 𝐂^𝗆`$ respectively $`(𝐏^1)^{\mathrm{}}\times 𝐂^𝗆`$. Let $`\stackrel{~}{\pi }:𝐂^n(𝐂^{})^{\mathrm{}}\times 𝐂^𝗆`$ be the universal covering. Then Noguchi’s proof yields that there exists a holomorphic function $`\stackrel{~}{\theta }:(𝐂^{})^{\mathrm{}}\times 𝐂^𝗆𝐂`$ which extends to a meromorphic function on $`(𝐏^1)^{\mathrm{}}\times 𝐂^𝗆`$ such that $`\theta =\stackrel{~}{\theta }\stackrel{~}{\pi }`$ is a theta function for $`DG`$ as above. More precisely, if $`(w_1,\mathrm{},w_t)`$ is a multiplicative coordinate system of $`(𝐂^{})^{\mathrm{}}`$ and $`U𝐂^𝗆`$ is a small neighborhood of a point in $`𝐂^𝗆`$, then $`\stackrel{~}{\theta }|_{(𝐂^{})^{\mathrm{}}\times U}`$ can be written as
$$\underset{\mathrm{finite}}{}a_{l_1\mathrm{}l_t}(y)w_1^{l_1}\mathrm{}w_t^{l_t},$$
(10)
where the coefficients $`a_{l_1\mathrm{}l_t}(y)`$ are holomorphic functions on $`U`$. As $`(𝐏^1)^{\mathrm{}}\times 𝐂^𝗆`$ is the universal covering space of $`\overline{G}`$, $`\stackrel{~}{\theta }`$ gives a multivalued defining function for $`D`$ on $`\overline{G}`$ locally given by equation (10). It follows that, whenever an algebraic expression in $`\theta `$ descends to $`G`$, it extends to $`\overline{G}`$ to a meromorphic object, and, hence, to a rational object.
Let $`\mu :\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }`$ be the universal cover, and let $`\stackrel{~}{f}:\stackrel{~}{\mathrm{\Gamma }}𝐂^n`$ be any lift of the map $`f\mu :\stackrel{~}{\mathrm{\Gamma }}GD`$. Then we can choose a linear coordinate system $`(z_1,\mathrm{},z_n)`$ of $`𝐂^n`$ such that in these coordinates, $`\stackrel{~}{f}`$ is expressed as
$$\stackrel{~}{f}(z)=(f_1(z),\mathrm{},f_n^{}(z),0,\mathrm{},0)$$
with holomorphic functions $`f_1(z),\mathrm{},f_n^{}(z)`$, for which the functions
$`1,f_1,f_2,\mathrm{},f_n^{}`$ are linearly independent.
The set of differential equations $`dz_{n^{}+1}=0`$,…,$`dz_n=0`$ obviously defines a subbundle $`\stackrel{~}{V}T𝐂^n`$, which is invariant under translation and, hence, descends to $`G`$. It is important to remark that this subbundle extends to a special subbundle $`V\overline{T}G`$. This is true, since by definition, $`V`$ is just of the form $`V=𝐂^n^{}\times G`$ with respect to the trivialization of $`\overline{T}G`$ given by the standard logarithmic forms $`dz_1,\mathrm{},dz_n`$.
Siu-Yeung () defined the following logarithmic jet differential:
$$\mathrm{\Theta }=\left|\begin{array}{cccc}d\mathrm{log}\theta \hfill & dz_1\hfill & \mathrm{}\hfill & dz_n^{}\hfill \\ d^2\mathrm{log}\theta \hfill & d^2z_1\hfill & \mathrm{}\hfill & d^2z_n^{}\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ d^{n^{}+1}\mathrm{log}\theta \hfill & d^{n^{}+1}z_1\hfill & \mathrm{}\hfill & d^{n^{}+1}z_n^{}\hfill \end{array}\right|$$
We want to use this $`\mathrm{\Theta }`$ as the $`\mathrm{\Theta }`$ in our Main Lemma 5.12. So we need:
###### Lemma 5.15
$`\mathrm{\Theta }`$ can be considered as a meromorphic section in the line bundle
$$𝒪_{\overline{P}_{n^{}+1}V}(\frac{(n^{}+1)(n^{}+2)}{2})$$
with at most log-poles along $`D`$.
Proof of Lemma 5.15 We first show (part (a)) that $`\mathrm{\Theta }`$ is a meromorphic function on $`J_k𝐂^n`$, and, hence, also on $`J_k\stackrel{~}{V}`$, which satisfies equation (2) with order $`k=n^{}+1`$ and degree $`m=(n^{}+1)(n^{}+2)/2=k(k+1)/2`$ (recall that $`\stackrel{~}{V}T𝐂^n`$ is defined by $`dz_{n^{}+1}=\mathrm{}=dz_n=0`$). By using equation (10) it follows actually that $`\mathrm{\Theta }`$ defines such a function on $`J_k(𝐂^{\mathrm{}}\times 𝐂^𝗆)`$, which extends, again by equation (10), to a meromorphic function on $`\overline{J}_k(𝐂^{\mathrm{}}\times 𝐂^𝗆)`$. This meromorphic function descends to a multivalued function on $`\overline{J}_kG`$. Then we show (part (b)) that it is actually singlevalued on $`\overline{J}_kV\overline{J}_kG`$. So it corresponds, by Proposition 3.9, to a meromorphic section of $`𝒪_{\overline{P}_kV}(m)`$. That it is a meromorphic section of $`𝒪_{\overline{P}_kV}(m)`$ with at most log-poles along $`D`$ then follows from the local description of $`\stackrel{~}{\theta }`$ in equation (10) and from the local description of meromorphic sections with at most log-poles along $`D`$ right after definition 5.7: In fact, by this description, applied to the multivalued meromorphic fuction on $`\overline{J}_kG`$, it gives rise to a (possibly) multivalued holomorphic function on $`\overline{J}_kY`$. But, as we saw above, it is singlevalued over $`\overline{J}_kV`$, and hence, it is also singlevalued over $`\overline{J}_kY`$. Thus, it yields a meromorphic section with at most log-poles along $`D`$.
(a) We show more generally: Let $`h_1,\mathrm{},h_r:(\mathrm{\Delta },0)𝐂`$ be nonvanishing germs of holomorphic functions. Let $`g_i=\mathrm{log}h_i`$, $`i=1,\mathrm{},r`$. Then
$$\left|\begin{array}{cccc}dg_1\hfill & dg_2\hfill & \mathrm{}\hfill & dg_r\hfill \\ d^2g_1\hfill & d^2g_2\hfill & \mathrm{}\hfill & d^2g_r\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ d^rg_1\hfill & d^rg_2\hfill & \mathrm{}\hfill & d^rg_r\hfill \end{array}\right|$$
gives a jet differential on $`J_r\mathrm{\Delta }_0`$ which is equivariant under the full reparametrization group $`J_r\mathrm{\Delta }_0`$ in the sense of equation (2). (We will apply this for the case where the germs $`h_j`$ are obtained by composing $`\theta `$ and the exponentials of the $`z_i`$’s with the germ $`f:(𝐂,0)𝐂^n`$ representing the jet in our case.)
Only the equivariance is nontrivial. Let $`g=(g_i)`$. Let $`\varphi J_r𝐂_0`$. Then, by using the identity
$$(g\varphi )^{(j)}(0)=g^{(j)}(0)(\varphi ^{}(0))^j+\underset{s=1}{\overset{j1}{}}\underset{i_1+\mathrm{}+i_s=j}{}c_{i_1\mathrm{}i_s}g^{(s)}(0)\varphi ^{(i_1)}(0)\mathrm{}\varphi ^{(i_s)}(0),$$
we get by induction on $`r`$ that
$$(g\varphi )^{}\mathrm{}(g\varphi )^{(r)}(0)=g^{}\mathrm{}g^{(r)}(\varphi ^{}(0))^{\frac{r(r+1)}{2}}.$$
This gives the desired equivariance.
(b) From equation (9) we have, for $`\gamma \mathrm{\Pi }_1(G)𝐂^n`$:
$$d^i\mathrm{log}\theta (x+\gamma )=d^i\mathrm{log}\theta (x)+d^iL_\gamma (x)=d^i\mathrm{log}\theta (x)+\underset{j=1}{\overset{n^{}}{}}a_jd^ix_j+\underset{j=n^{}+1}{\overset{n}{}}a_jd^ix_j,$$
where $`a_i𝐂`$ are constants. Then, from the properties of the determinant and the fact that we restrict $`\mathrm{\Theta }`$ to $`J_{n^{}+1}\stackrel{~}{V}`$, it follows that this jet differential is invariant under the action of $`\mathrm{\Pi }_1(G)`$. Hence, it descends to $`G`$. $``$$``$
###### Lemma 5.16
Under the assumptions of Theorem 5.1 (b), or Theorem 5.3 (b) and (c) and the additional assumption that $`f`$ does not extend, the following holds:
If $`X(f)D\mathrm{}`$, then $`X(f)D`$ is foliated by translates of an algebraic subgroup $`G^{\prime \prime }G^{}`$ of positive dimension, where $`G^{}`$ is the maximal subgroup whose translates foliate $`X(f)`$.
Proof of Lemma 5.16 We may assume that $`f`$ is nonconstant and, for the case that $`\mathrm{\Gamma }=\mathrm{\Delta }^{}`$, that the map $`f`$ does not extend. We apply the Main Lemma 5.12 to get the existence of an algebraic subgroup $`G^{\prime \prime }G`$ of positive dimension which leaves $`X_{n^{}+1}(f)`$ and $`\mathrm{\Theta }|_{X_{n^{}+1}(f)}`$ invariant. As $`X_{n^{}+1}(f)`$ is invariant under the action of $`G^{\prime \prime }`$ and the projection $`\pi _{1,n^{}+1}`$ (respectively $`\pi _{0,n^{}+1}`$) maps $`X_{n^{}+1}(f)`$ surjectively onto $`X_1(f)`$ (respectively $`X(f)`$), we see that $`X_1(f)`$ and $`X(f)`$ are also invariant.
Take any $`aG^{\prime \prime }`$. Since $`\mathrm{\Theta }|_{X_{n^{}+1}(f)}`$ is invariant under translation by $`a`$,
$$\left|\begin{array}{cccc}d\mathrm{log}\frac{\theta (f)}{\theta (f+a)}\hfill & df_1\hfill & \mathrm{}\hfill & df_n^{}\hfill \\ d^2\mathrm{log}\frac{\theta (f)}{\theta (f+a)}\hfill & d^2f_1\hfill & \mathrm{}\hfill & d^2f_n^{}\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ .\hfill & .\hfill & .\hfill & .\hfill \\ d^{n^{}+1}\mathrm{log}\frac{\theta (f)}{\theta (f+a)}\hfill & d^{n^{}+1}f_1\hfill & \mathrm{}\hfill & d^{n^{}+1}f_n^{}\hfill \end{array}\right|0\mathrm{on}\mathrm{\Gamma }.$$
Now $`d\mathrm{log}\frac{\theta (x)}{\theta (x+a)}`$ is a rational differential on $`G`$. Since $`f+a`$ cannot map entirely into the zero set of $`\theta `$, because $`X(f)`$ is the Zariski closure of $`f(\mathrm{\Gamma })`$,
$$\frac{}{z}(\mathrm{log}\frac{\theta (f)(z)}{\theta (f+a)(z)}),\frac{}{z}f_1(z),\mathrm{},\frac{}{z}f_n^{}(z)$$
are well defined meromorphic functions on $`\mathrm{\Gamma }`$. The functions $`\frac{}{z}f_1(z),\mathrm{},`$ $`\frac{}{z}f_n^{}(z)`$ are linearly independent as $`1,f_1,\mathrm{},f_n^{}`$ were so. Hence, we get, by the classical Lemma of the Wronskian , that there exist complex numbers $`c_1,\mathrm{},c_n^{}`$ (which may depend on $`aG^{\prime \prime }`$) such that
$$d\mathrm{log}\frac{\theta (f)(z)}{\theta (f+a)(z)}+c_1df_1(z)+\mathrm{}+c_n^{}df_n^{}(z)0\mathrm{on}\mathrm{\Gamma }.$$
So we have
$$d\mathrm{log}\frac{\theta (x)}{\theta (x+a)}+c_1dx_1+\mathrm{}+c_n^{}dx_n^{}0$$
(11)
on $`f_{[1]}(\mathrm{\Gamma })`$. Moreover, since $`d\mathrm{log}\frac{\theta (x)}{\theta (x+a)}`$ is a rational differential on $`G`$, this equation holds on $`X_1(f)`$.
Assume now that Lemma 5.16 does not hold. Then there exists $`x_0X(f)D`$ and $`a_0G^{\prime \prime }`$ such that $`x_0+a_0D`$. We want to show that this assumption leads to a contradiction. From equation (11) we get that
$$d\mathrm{log}\frac{\theta (x)}{\theta (x+a_0)}=d\mathrm{log}\frac{\theta (x+b)}{\theta (x+a_0+b)}$$
(12)
on $`X_1(f)`$ for $`bG^{\prime \prime }`$. This means that
$$d\mathrm{log}(\frac{\theta (f)}{\theta (f+b)}\frac{\theta (f+a_0+b)}{\theta (f+a_0)})0\mathrm{on}\mathrm{\Gamma }.$$
Hence,
$$\frac{\theta (f)}{\theta (f+b)}\frac{\theta (f+a_0+b)}{\theta (f+a_0)}=c_{a_0,b}\mathrm{on}\mathrm{\Gamma },$$
where $`c_{a_0,b}𝐂`$ is a constant, which may depend on $`a_0`$ and $`b`$. Since $`\frac{\theta (x)}{\theta (x+b)}\frac{\theta (x+a_0+b)}{\theta (x+a_0)}`$ is a well defined rational function on $`G`$, we have
$$\frac{\theta (x)}{\theta (x+b)}\frac{\theta (x+a_0+b)}{\theta (x+a_0)}=c_{a_0,b}\mathrm{on}X(f),$$
(13)
where $`bG^{\prime \prime }`$. Now $`x_0+a_0D`$, but $`x_0D`$. So we get, for $`b=a_0`$ and $`x=x_0`$, that $`c_{a_0,a_0}=0`$. This means, as $`X(f)`$ is irreducible, that either $`\theta (x)0`$ or $`\theta (x+2a_0)0`$ on $`X(f)`$. But both is not true, as one sees by taking $`x=x_0+a_0`$ respectively $`x=x_0+a_02a_0`$ (remark that the latter is still in $`X(f)`$, since $`X(f)`$ is invariant under the action of $`G^{\prime \prime }`$). So our assumption was wrong, and we have proved Lemma 5.16. $``$$``$
Proof of Theorem 5.1 (b) Assume that $`X(f)D\mathrm{}`$. We want to show that this assumption leads to a contradiction. After applying a translation, we may assume, by Theorem 5.1 (a), that $`X(f)`$ again is a semi abelian variety $`G^{}`$ with nonempty divisor $`D^{}`$ in $`G^{}`$, where $`D^{}`$ is the reduction of $`X(f)D`$. Now we devide through the maximal algebraic subgroup $`\stackrel{~}{G}`$ of $`G^{}`$ which foliates $`D^{}`$. Then, by applying Lemma 5.16 to the quotient $`G^{}/\stackrel{~}{G}`$ and by taking the invers image under the quotient map, we get $`X(f)D=\mathrm{}`$, which contradicts our assumption. $``$$``$
Proof of Theorem 5.3 (b) and (c) Part (b) follows immediately from Lemma 5.16. For (c), let $`G`$ be an abelian variety. Let $`G^{}`$ again be the maximal algebraic subgroup the translates of which foliate $`X(f)`$. We may assume that all translates of $`G^{}`$ which foliate $`X(f)`$ intersect $`D`$ (in particular $`X(f)D\mathrm{}`$), for otherwise we finish the proof by using Lemma 5.6. Then there must be such a translate $`T_0`$ of $`G^{}`$ such that $`T_0D`$. Now by Lemma 5.16 we find a subgroup $`G^{\prime \prime }G^{}`$ of positive dimension which foliates $`X(f)D`$. Hence, $`TD`$ is foliated by translates of $`G^{\prime \prime }`$. But since $`T`$ is also foliated by translates of $`G^{\prime \prime }`$, there must be such a translate not hitting $`D`$ at all. This finishes the proof again by using Lemma 5.6. $``$$``$
## 6 Appendix
We use the notations of subsection 4.1. We now give a key result via which pseudometrics of negative curvature are usually constructed (see (2.7) of ). We point out that our version is sharper than the ones in for the basic locus in our definition is smaller than theirs.
###### Lemma 6.1
Let the setup be as in subsection 4.1, and assume further that $`X`$ is normal. Then given line bundles $`L`$ and $`H`$ over $`X`$, there is an integer $`l_11`$ such that $`xE_L`$ implies that $`x\mathrm{Bs}|lLH|`$ for all positive multiples $`l`$ of $`l_1`$, more specifically,
$$E_L\mathrm{Bs}|lLH|.$$
Proof Observe that we may always write $`H+H^{}=H^{\prime \prime }`$, where $`H^{}`$ and $`H^{\prime \prime }`$ are very ample divisors. Then $`\text{Bs}|lLH^{\prime \prime }|\text{Bs}|lLH|`$, as one can see from the fact that $`\text{Bs}|GG^{}|\text{Bs}|G^{}|\text{Bs}|G|`$ for arbitrary line bundles $`G`$ and $`G^{}`$. Hence, we will assume without loss of generality that $`H`$ is very ample.
Let $`xX`$ be outside $`E_L`$. Then, we may assume that $`\phi _L`$ is birational onto its image after replacing $`L`$ by a suitable multiple of $`L`$ (see 1.10 and 5.7 of ). It will be sufficient to show that $`x`$ is outside $`\mathrm{Bs}|lLH|`$ for some $`l`$, and hence, for all multiples thereof, as Lemma 6.1 would then follow by the quasi-compactness of $`XE_L`$.
Consider the ideal sheaf $`𝒪_X`$ generated by the global sections of $`L`$. By blowing up this ideal sheaf, we obtain a modification $`\sigma :\stackrel{~}{X}X`$ so that $`𝒥=\sigma ^{}𝒪_{\stackrel{~}{X}}`$ is an invertible ideal sheaf of $`𝒪_{\stackrel{~}{X}}`$ generated by a global section $`s`$ of the line bundle $`F=𝒪_\sigma (1)`$, namely, $`𝒥=\text{Im}\{𝒪(F^{})\stackrel{s}{}𝒪\}.`$ Then $`\overline{L}:=\sigma ^1LF`$ is spanned by the sections $`(\sigma ^1t)/s`$ as $`t`$ ranges over $`H^0(L)`$. Note that $`E_{\overline{L}}=\sigma ^1(E_L)`$, that $`\sigma ^1`$ is an isomorphism on the neighborhood $`XE_L`$ of $`x`$ and that any section of $`l\overline{L}\sigma ^1H`$ not vanishing on the point $`\sigma ^1(x)`$ gives rise to a section of $`lLH`$ not vanishing on $`x`$ by tensoring with $`s^l`$ (Zariski’s Main Theorem and $`s(x)0`$). Hence, replacing $`(X,L)`$ by $`(\stackrel{~}{X},\overline{L})`$ we may assume that $`\phi _L:X𝐏^n`$ is a birational morphism onto its image $`W=\phi _L(X)`$. Let $`\sigma _0:W_0W𝐏^n`$ be the normalization of $`W`$. Then $`H_0:=\sigma _0^1𝒪_{𝐏^n}(1)`$ is ample so that there is a positive integer $`d`$ such that $`dH_0`$ is very ample on $`W_0`$. As $`X`$ is normal, there is a canonical morphism $`\phi :XW_0`$ such that $`\phi _L=\sigma _0\phi `$. Noting $`\phi ^1H_0=L,`$ we see that the image $`W_1`$ of the morphism $`\phi _{dL}`$ admits a birational morphism $`r`$ to $`W_0`$ and that $`\phi =r\phi _{dL}`$. As $`\phi `$ is connnected by Zariski’s Main Theorem, $`\phi ^1(\phi (x))=\{x\}`$. Hence, replacing $`L`$ by $`dL`$ we may assume that
$$\phi _L^1(\phi _L(x))=\{x\}.$$
(1)
As $`H`$ is very ample, $`|H|`$ has an element $`D`$ such that $`xDE_L`$ by general positioning. We may now choose, thanks to equation (1), a hypersurface of sufficiently high degree $`l`$ in $`𝐏^n`$ containing $`\phi _L(D)`$ but not $`\phi _L(x)`$. This gives a divisor in $`|lLH|`$ not containing $`x`$ as desired. $``$$``$
In practice, Lemma 6.1 is all that one uses. But one can easily deduce the following strengthened version in order to complete the picture.
###### Lemma 6.2
Let $`X`$ be a normal complex projective variety with any line bundle $`H`$. For any line bundle $`L`$ over $`X`$, there is an integer $`m_0`$ such that
$$E_{m_0L}S_L\text{Bs}|m_0LH|.$$
Moreover, if $`H`$ is very ample, the both inclusions are equalities.
Proof Clearly there is an integer $`N`$ such that $`S_L=_{m=1}^NE_{mL}.`$ For each $`m`$, there is an integer $`l_m>0`$, such that $`E_{mL}\text{Bs}|lLH|`$ for all positive multiple $`l`$ of $`l_m`$ by Lemma 6.1. Letting $`m_0`$ be a common multiple of $`l_1,\mathrm{},l_N`$, we see that $`E_{m_0L}S_L\text{Bs}|m_0LH|.`$ If, furthermore, $`H`$ is very ample, one easily verifies that $`\text{Bs}|mLH|E_{mL}`$ for all $`m`$. $``$$``$
Remark As $`\text{Bs}|lLH|E_{lL}S_L`$, it follows that $`S_L=_{l>0}\text{Bs}|lLH|`$ for any very ample $`H`$.
Gerd Dethloff
Département de Mathématiques, UFR Sciences et Techniques
Université de Bretagne Occidentale
6, avenue Le Gorgeu, BP 452
29275 Brest Cedex, France.
e-mail: dethloff@univ-brest.fr
Steven Lu
Department of Pure Mathematics
University of Waterloo
Waterloo, Ontario, Canada N2L 3G1
e-mail: slu@math.uwaterloo.ca |
warning/0001/hep-th0001160.html | ar5iv | text | # The Four-fermi Coupling of the Supersymmetric Non-linear 𝜎-model on 𝐺/𝑆⊗{𝑈(1)}^𝑘
## 1 Introduction
The non-linear $`\sigma `$-model on a coset space $`G/H`$ is a low-energy effective theory of the Nambu-Goldstone(NG) bosons generated when a large group symmetry $`G`$ at high energy spontaneously breaks into a small one at low energy. If the supersymmetry exists at that high energy and if it survives in the spontaneous breaking, the NG bosons are accompanied by superpartners, called the pseudo NG fermions. They are described by the supersymmetric non-linear $`\sigma `$-model on a coset space $`G/H`$ which is kählerian. In the beginning of the 80’s Buchmüller, Love, Peccei and Yanagida proposed to identify the massless pseudo NG fermions with light quarks and leptons and used the supersymmetric non-linear $`\sigma `$-model as a low-energy effective theory for the grand unification. Recently in ref. this idea was revived to explain the neutrino mass observed in the SuperKamiokande experiment. They proposed the supersymmetric non-linear $`\sigma `$-model on the Kähler coset space $`E_7/SU(5)\{U(1)\}^3`$ as a thoery which naturally accomodates the three families of right-handed neutrinos. Namely the model contains three families of $`\mathrm{𝟏𝟎}+\mathrm{𝟓}^{}+\mathrm{𝟏}`$ and $`\mathrm{𝟓}`$ of $`SU(5)`$ as the NG supermultiplet $`(\varphi ^\alpha ,\psi ^\alpha )`$. Interactions among the pseudo NG fermions take place through the four-fermi coupling
$$R_{\alpha \overline{\sigma }\beta \overline{\delta }}(\frac{\varphi }{f},\frac{\overline{\varphi }}{f})(\overline{\psi }^{\overline{\sigma }}\psi ^\alpha )(\overline{\psi }^{\overline{\delta }}\psi ^\beta )$$
in which $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ is the Riemann curvature of the coset space and $`f`$ is a constant giving a mass scale. It is phenomelogically the most interesting part of the model. The aim of this paper is to establish a practical method to calculate the Riemann curvature of $`E_7/SU(5)\{U(1)\}^3`$.
$`E_7/SU(5)\{U(1)\}^3`$ is a fairly complicated coset space. The complication comes in twofold. Firstly the coset space is reducible. One can calculate the holomorphic Killing vectors and the Kähler potential of the reducible coset space. But they take more cumbersome forms than those of the irreducible coset space. Secondly the homogeneous group includes plural $`U(1)`$s as direct products. In such a case the complex structure of the Kähler coset space is not unique and the metric depends on as many free parameters as $`U(1)`$s. These subjects on the reducible Kähler coset space were extensively studied by the Kyoto group in ref. . The general method to construct the Kähler potential of the reducible coset space was given. The Riemann curvature may be calculated by differentiating that Kähler potential by coordinates, in principle. But it is too involved.
In this paper we employ an alternative formalism to do this more directly, which was proposed in ref. . It is based on the Killing potentials instead of the Kähler potential, which are also characteristic for the Kähler coset space $`G/H`$ . In ref. the formalism was developed for the irreducible case, and the Riemann curvature was given in terms of the holomorphic Killing vectors with no derivative by coordinates. (See eq. (4.14).) Once given a concrete form of the holomorphic Killing vectors, one can directly calculate the Riemann curvature by that formula. In this paper we extend this formalism to study the reducible Kähler coset space along the same line. We will be particularly interested in the Riemann curvature at the low-energy limit $`f\mathrm{}`$, i.e., $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}(0,0)`$. It gives the four-fermi coupling constants at low energy, which depend on as many free parameters as $`U(1)`$s of the homogeneous subgroup. Knowing the dependence explicitly is very interesting from the phenomelogical point of view.
In Section 2 we summarize the geometry of the Kähler coset space $`G/H`$. In Section 3 we briefly explain the generalized CCWZ formalism. It enables us to construct the holomorphic Killing vectors, the metric and the Kähler potential for the reducible case. In Section 4 the alternative formalism based on the Killing potentials is presented which allows us to calculate the metric and the Riemann curvature more directly. We first of all review the formalism which was developed for the irreducible case. Then it will be extended to the reducible case. We derive the general formula which expresses the Riemann curvature of the reducible Kähler coset space in terms of the holomorphic Killing vectors.(See eq. (4.32).) In Section 5 the Riemann curvatures of $`SU(3)/\{U(1)\}^2`$ and $`E_7/SU(5)\{U(1)\}^3`$ are evaluated to the leading order of $`\frac{1}{f}`$ by this general formula.
## 2 The geometry of the Kähler manifold
In this section we briefly review on the Kähler manifold, giving our notation. Consider a real $`2N`$-dimensional Riemann manifold $``$ with local coordinates $`\varphi ^a=(\varphi ^1,\varphi ^2,\mathrm{}\mathrm{},\varphi ^{2N})`$. The line element of the manifold is given by
$`ds^2=g_{ab}d\varphi ^ad\varphi ^b.`$ (2.1)
$``$ is a Kähler manifold if it is endowed with a complex structure which is covariantly constant:
$`J_{b;c}^a=0,`$ (2.2)
and satisfies $`J_b^aJ_c^b=\delta _c^a`$. We assume the metric $`g_{ab}`$ to be of type $`(1,1)`$, i.e.,
$`g_{ab}=g_{cd}J_a^cJ_b^d.`$ (2.3)
(A tensor of more general type $`(r,s)`$ will be discussed soon later.) The symplectic structure $`J_{ab}`$ is given by $`J_{ab}=g_{ac}J_b^c.`$ By (2.2) it is closed:
$`J_{ab,c}+J_{bc,a}+J_{ca,b}=0.`$ (2.4)
When the Kähler manifold is a coset space $`G/H`$, there is a set of Killing vectors
$`R^{Aa}=(R^{1a},R^{2a},\mathrm{}\mathrm{},R^{Da}),`$ (2.5)
with $`D=\mathrm{dim}G`$ , which represents the isometry $`G`$. They satisfy
$`_^𝒜R^{Ba}`$ $`=`$ $`R^{Ab}R_{,b}^{Ba}R^{Bb}R_{,b}^{Aa}=f^{ABC}R^{Ca}`$
$`(\mathrm{Isometry}),`$
$`_^𝒜g_{ab}`$ $`=`$ $`R^{Ac}g_{ab,c}+R_{,a}^{Ac}g_{cb}+R_{,b}^{Ac}g_{ca}=0`$
$`(\mathrm{Killing}\mathrm{condition}),`$
$`_^𝒜J_b^a`$ $`=`$ $`R^{Ac}J_{b,c}^aR_{,c}^{Aa}J_b^c+R_{,b}^{Ac}J_c^a=0,`$ (2.8)
in which $`_^𝒜`$ is the Lie-derivative with respect to $`R^A`$, and $`f^{ABC}`$ are structure constants of the isometry group $`G`$.
Any vector $`v^a`$ can be projected onto the $`(1,0)`$ and $`(0,1)`$ types by
$`{\displaystyle \frac{1}{2}}(1iJ)_b^av^b,{\displaystyle \frac{1}{2}}(1+iJ)_b^av^b.`$ (2.9)
A tensor of the $`(r,s)`$ type is obtained as a multi-product of these projected vectors. We may locally set the complex structure to be
$`J_b^a=\left(\begin{array}{cc}i\delta _\beta ^\alpha & 0\\ & \\ 0& i\delta _{\overline{\beta }}^{\overline{\alpha }}\end{array}\right),`$ (2.13)
with $`\alpha ,\overline{\alpha }=1,2,\mathrm{},N.`$ Then the respective vectors in (2.9) may be written as $`N`$-dimensional complex vectors $`v^\alpha `$ and $`v^{\overline{\alpha }}`$. The line element (2.1) is written as
$$ds^2=g_{\alpha \overline{\beta }}d\varphi ^\alpha d\varphi ^{\overline{\beta }},$$
by (2.3). The closure of the symplectic structure given by (2.4) reads
$`g_{\alpha \overline{\beta },\gamma }=g_{\gamma \overline{\beta },\alpha },g_{\alpha \overline{\beta },\overline{\gamma }}=g_{\alpha \overline{\gamma },\overline{\beta }}.`$ (2.14)
Then it follows that there exists a real scalar $`K(\varphi ,\overline{\varphi })`$, called Kähler potential such that
$`g_{\alpha \overline{\beta }}=K_{,\alpha \overline{\beta }}.`$ (2.15)
Furthermore (2.8) and (2.10) imply that the Killing vectors are holomorphic:
$`R_{,\alpha }^{A\overline{\beta }}=0,R_{,\overline{\alpha }}^{A\beta }=0.`$ (2.16)
Then (2.6) and (2.7) reduce respectively to
$`_^𝒜R^{B\alpha }`$ $`=`$ $`R^{A\beta }R_{,\beta }^{B\alpha }R^{B\beta }R_{,\beta }^{A\alpha }=f^{ABC}R^{C\alpha },`$
$`\mathrm{c}.\mathrm{c}.,`$
and
$`_^𝒜g_{\alpha \overline{\beta }}`$ $`=`$ $`R_{\alpha ,\overline{\beta }}^A+R_{\overline{\beta },\alpha }^A=0,`$ (2.18)
with $`R_\alpha ^A=g_{\alpha \overline{\beta }}R^{A\overline{\beta }}`$ and $`R_{\overline{\alpha }}^A=g_{\beta \overline{\alpha }}R^{A\beta }`$. From (2.15) we may find real scalars $`M^A(\varphi ,\overline{\varphi })`$, called Killing potentials, such that
$`R_\alpha ^A=iM_{,\alpha }^A,R_{\overline{\alpha }}^A=iM_{\overline{\alpha }}^A.`$ (2.19)
As shown in ref., they transform as the adjoint representation of the group $`G`$ by the Lie-variation
$`_^𝒜M^B=R^{A\alpha }M_{,\alpha }^B+R^{A\overline{\alpha }}M_{,\overline{\alpha }}^B=f^{ABC}M^C.`$ (2.20)
A manipulation of (2.17) with (2.16) leads us to write the Killing potentials in terms of the Killing vectors :
$`M^A={\displaystyle \frac{i}{𝒩_{adj}}}f^{ABC}R^{B\alpha }R^{C\overline{\beta }}g_{\alpha \overline{\beta }}.`$
Here we have used the normalization
$`f^{ABC}f^{ABD}=2𝒩_{adj}\delta ^{CD}.`$ (2.21)
These Killing potentials characterize the Kähler manifold no less than the Kähler potential, if it is a coset space $`G/H`$.
## 3 The CCWZ formalism
In this section we will explain how to construct the holomorphic Killing vectors $`R^{A\alpha }`$, the metric $`g_{\alpha \overline{\beta }}`$ and the Kähler potential $`K`$, which essentially characterize the Kähler coset space $`G/H`$. When the Kähler coset space $`G/H`$ is irreducible, they can be constructed case by case in heuristic ways. But for the reducible case we need a systematic method. It was given by generalizing the CCWZ formalism by the Kyoto group. We briefly sketch this generalized CCWZ formalism.
### 3.1 The holomorphic Killing vectors
We assume the isometry group $`G`$ is compact and semi-simple. If a coset space $`G/H`$ is kählerian , the unbroken subgroup $`H`$ contains $`U(1)`$ groups as $`H=S\{U(1)\}^k,k=1,2,\mathrm{},n`$, according to the Borel theorem. The generators $`T^A`$ of $`G`$ are decomposed as
$`\{T^A\}=\{X^a,S^I,Q^\mu \},a`$ $`=`$ $`1,2,\mathrm{},2N(=\mathrm{dim}G\mathrm{dim}H),`$
$`I`$ $`=`$ $`1,2,\mathrm{},\mathrm{dim}S(=\mathrm{dim}Hk),`$
$`\mu `$ $`=`$ $`1,2,\mathrm{},k,`$ (3.1)
in which $`S^I`$ and $`Q^\mu `$ are generators of $`S`$ and $`U(1)`$s respectively, while $`X^a`$ broken generators. Let us define a central charge as
$`Y={\displaystyle \underset{\mu =1}{\overset{k}{}}}v^\mu Q^\mu vQ,`$ (3.2)
by choosing real coefficients $`v^\mu `$ such that all the broken generators $`X^a`$ have non-vanishing $`Y`$-charges. Then the broken generators $`X^a`$ can be splitted into two parts: the generators $`X^{\overline{i}}`$ with positive $`Y`$-charge and their hermitian conjugates $`X^i`$ with negtaive charge, $`i,\overline{i}=1,2,\mathrm{},N`$. (3.1) is further decomposed as
$`\{T^A\}=\{X^{\overline{i}},X^i,S^I,Q^\mu \}.`$ (3.3)
The splitting of the broken generators determines the complex structure $`J_b^a`$ of the Kähler coset space $`G/H`$. But the splitting is not unique depending on the definition of the central charge (3.2). It implies arbitrariness of the complex structure of the coset space.
For the decompostion (3.3) the standard application of the CCWZ formalism does not give the holomorphic Killing vectors $`R^{A\alpha }`$ satisfying the Lie-algebra (2.14). Hence we extend the isometry group $`G`$ to the complex one $`G^c`$ and consider a coset space $`G^c/\widehat{H}`$ with the complex subgroup $`\widehat{H}`$ generated by $`X^i,S^I,Q^\mu `$. As explicitly given later, there is an isomorphism between this complex coset space $`G^c/\widehat{H}`$ and $`G/H`$:
$`G/HG^c/\widehat{H}.`$ (3.4)
The holomorphic Killing vectors are obtained by applying the CCWZ formalism to the complex coset space $`G^c/\widehat{H}`$. The coset space $`G^c/\widehat{H}`$ is parametrized by complex coordinates $`\varphi ^\alpha ,\alpha =1,2,\mathrm{},N`$ corresponding to the broken generators $`X^i`$. Consider a holomorphic quantity
$`\xi (\varphi )=e^{\varphi \overline{X}}G^c/\widehat{H}`$ (3.5)
with <sup>1</sup><sup>1</sup>1(3.5) should have been written as $`\xi (\varphi )=\mathrm{exp}(\frac{1}{f}\varphi \overline{X})`$ with the mass scale parameter $`f`$. But it is hereinafter set to be one to avoid unnecessary complication.
$$\varphi X^{\overline{i}}=\varphi ^1X^{\overline{1}}+\varphi ^2X^{\overline{2}}+\mathrm{}\varphi ^NX^{\overline{N}}.$$
For an element $`g`$ of the isometry group $`G`$, i.e., $`g=e^{iϵ^AT^A}G`$ with real parameters $`ϵ^A`$, there exists a compensator $`\widehat{h}(\varphi ,\overline{\varphi },g)\widehat{H}`$ such that
$`g\xi (\varphi )=\xi (\varphi ^{})\widehat{h}(\varphi ,\overline{\varphi },g).`$ (3.6)
This defines a holomorphic transformation of the coordinates $`\varphi ^\alpha `$ which realizes the isometry group non-linearly. When the real parameters $`ϵ^A`$ are infinitesimal, (3.6) yields the holomorphic Killing vectors $`R^{A\alpha }(\varphi )`$ as
$`\delta \varphi =\varphi ^\alpha (\varphi )\varphi ^\alpha =ϵ^AR^{A\alpha }(\varphi ),`$ (3.7)
which satisfy the Lie-algebra (2.14).
### 3.2 The metric
Any two points on the coset space $`G/H`$ can be related by the isometry transformation (3.6). Therefore the line element (2.1) has the same length at any point of the coset space
$`g_{ab}(\varphi ^{},\overline{\varphi }^{})d\varphi ^ad\varphi ^b=g_{ab}(\varphi ,\overline{\varphi })d\varphi ^ad\varphi ^b.`$ (3.8)
On the other hand the line element is invariant under general coordinate transformations:
$`g_{ab}(\varphi ,\overline{\varphi })d\varphi ^ad\varphi ^b=g_{ab}^{}(\varphi ^{},\overline{\varphi }^{})d\varphi ^ad\varphi ^b.`$ (3.9)
(3.8) and(3.9) require that
$`g_{ab}(\varphi ^{},\overline{\varphi }^{})=g_{ab}^{}(\varphi ^{},\overline{\varphi }^{})`$ (3.10)
which gives the Killing condition (2.7) in the infinitesimal form.
To construct the metric $`g_{ab}`$ which satisfy the condition (3.10) we have recourse to the CCWZ formalism. Consider a quantity
$`U(\varphi ,\overline{\varphi })G/H,`$ (3.11)
with $`U^{}U=UU^{}=1`$. But the standard parametrization of $`U`$, i.e, $`U(\varphi ,\overline{\varphi })=e^{\varphi \overline{X}+\overline{\varphi }X}`$ does not give the metric of the type (1,1). Therefore we employ the non-standard one, namely
$`U(\varphi ,\overline{\varphi })=\xi (\varphi )\zeta (\varphi ,\overline{\varphi }),`$ (3.12)
in which $`\xi (\varphi )`$ is the element (3.5), while $`\zeta (\varphi ,\overline{\varphi })`$ an element of $`\widehat{H}`$. We parametrize the latter as
$`\zeta (\varphi ,\overline{\varphi })=e^{a(\varphi ,\overline{\varphi })X}e^{b(\varphi ,\overline{\varphi })S}e^{c(\varphi ,\overline{\varphi })Q},`$ (3.13)
with
$`aX={\displaystyle \underset{i=1}{\overset{N}{}}}a^iX^i,bS={\displaystyle \underset{I=1}{\overset{\mathrm{dim}Hk}{}}}b^IS^I.`$
Here the function $`b(\varphi ,\overline{\varphi })`$ and $`c(\varphi ,\overline{\varphi })`$ are chosen to be real since their purely imaginary parts can be absorbed into an element of $`H`$. Then the parametrization (3.13) is determined by the unitary condition $`U^{}U=1`$ which reads
$`\xi ^{}(\overline{\varphi })\xi (\varphi )=e^{\overline{a}(\varphi ,\overline{\varphi })\overline{X}}e^{2b(\varphi ,\overline{\varphi })S}e^{2c(\varphi ,\overline{\varphi })Y}e^{a(\varphi ,\overline{\varphi })X}.`$
(3.12) is an concrete expression of the isomorphism (3.4) relating the respective elements (3.5) and (3.11) of the coset spaces $`G/H`$ and $`G^c/\widehat{H}`$.
The fundamental object to construct the metric $`g_{\alpha \overline{\beta }}`$ is the Cartan-Maurer $`1`$-form
$`\omega `$ $`=`$ $`U^1dU`$ (3.14)
$`=`$ $`e^iX^{\overline{i}}+e^{\overline{i}}X^i+\omega ^IS^I+\omega ^\mu Q^\mu ,`$
with the $`1`$-forms $`e^i(\varphi ,\overline{\varphi }),e^{\overline{i}}(\varphi ,\overline{\varphi }),\omega ^I(\varphi ,\overline{\varphi })`$ and $`\omega ^\mu (\varphi ,\overline{\varphi })`$ as coefficients of the expansion. In particular $`e^i(\varphi ,\overline{\varphi })`$ takes the form
$`e^i=e_\alpha ^id\varphi ^\alpha ,`$
with no $`d\varphi ^{\overline{\alpha }}`$, as can be seen from the parametrization (3.12). $`e^{\overline{i}}(\varphi ,\overline{\varphi })`$ is its complex conjugate. The components $`e_\alpha ^i`$ and $`e_{\overline{\alpha }}^{\overline{i}}`$ are vielbeins of the local frame of the coset space. From this it follows that
$`g_{\alpha \beta }`$ $`=`$ $`g_{\overline{\alpha }\overline{\beta }}=0,`$
$`g_{\alpha \overline{\beta }}`$ $`=`$ $`g_{\overline{\beta }\alpha }={\displaystyle \underset{i=1}{\overset{N}{}}}y_i(v)e_\alpha ^ie_{\overline{\beta }}^{\overline{i}},`$ (3.15)
where $`y_i(v)`$ is the positive $`Y`$-charge (3.2) of the broken generator $`X^{\overline{i}}`$:
$`[Y,X^i]`$ $`=`$ $`y_i(v)X^i,[Y,X^{\overline{i}}]=y_i(v)X^{\overline{i}}.`$ (3.16)
By the transformation (3.6) $`U`$ transforms as
$`gU(\varphi ,\overline{\varphi })=U(\varphi ^{},\overline{\varphi }^{})h(\varphi ,\overline{\varphi },g),`$ (3.17)
with a compensator $`hH`$. Then $`e^i`$ transform homogeneously as
$`e^i(\varphi ^{},\overline{\varphi }^{})=\rho ^{ij}(h(\varphi ,\overline{\varphi },g),g)e^j(\varphi ,\overline{\varphi }),`$ (3.18)
in which $`\rho ^{ij}(h,g)`$ is the $`N`$-dimensional representation of the subgroup $`H`$. Consequently the metric (3.15) satifies the transformation property (3.8) under (3.6) or equivalently (3.17). Furthermore (3.16) guarantees the closure property of the metric (2.11). If the Kähler coset space $`G/H`$ is reducible, the broken generators $`X^i`$ are decomposed into irreducible sets under the subgroup $`H`$, each of which may have a different $`Y`$-charge due to the Schur’s Lemma.
It can be also shown that one can write the metric (3.15) as (2.12) with the Kähler potential
$`K(\varphi ,\overline{\varphi })={\displaystyle \underset{\mu =1}{\overset{k}{}}}v^\mu c^\mu (\varphi ,\overline{\varphi }),`$ (3.19)
where $`c^\mu `$ are the functions appearing in the parametrization (3.13) and $`v^\mu `$ are the coefficients defining the $`Y`$-charge (3.2).
## 4 The Riemann curvature
The Riemann curvature of the Kähler manifold is given by
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`g_{\eta \overline{\sigma }}\mathrm{\Gamma }_{\alpha \beta ,\overline{\delta }}^\eta `$
$`=`$ $`g_{\alpha \overline{\sigma },\beta \overline{\delta }}g^{\kappa \overline{\lambda }}g_{\alpha \overline{\lambda },\beta }g_{\kappa \overline{\sigma },\overline{\delta }}.`$
To obtain it explicitly we have to compute the metric $`g_{\alpha \overline{\beta }}`$ in the first place. It may be done with (3.15) by calculating the vielbeins $`e_\alpha ^i`$ or with (2.12) by calculating the Kähler potential (3.19). Either calculation is already complicated. It is further complicated to take the derivative $`g_{\alpha \overline{\sigma },\beta \overline{\delta }}`$ to obtain the Riemann curvature. Hence in this section we will study a method which enables us to calculate the Riemann curvature in a more direct way.
### 4.1 The irreducible case
When the Kähler manifold $`G/H`$ is irreducible, all the broken generators $`X^{\overline{i}}`$ have the same $`Y`$-charge $`y(v)(>0)`$. Then the metric (3.15) becomes simple:
$`g_{\alpha \overline{\beta }}=y(v){\displaystyle \underset{i=1}{\overset{N}{}}}e_\alpha ^ie_{\overline{\beta }}^{\overline{i}},`$ (4.2)
the value of which at the origin of the manifold is
$`g_{\alpha \overline{\beta }}\underset{\varphi =\overline{\varphi }=0}{|}=y(v)\delta _{\alpha \overline{\beta }}.`$ (4.3)
It was the Killing condition (2.14) that allows us to write the metric in the form of (4.2). The Killing condition can be satisfied also by giving the metric in terms of the Killing vectors (2.5): $`g^{\alpha \overline{\beta }}R^{A\alpha }R^{A\overline{\beta }}`$. Fixing the free parameter by the initial condition (4.3) we then have
$`g^{\alpha \overline{\beta }}={\displaystyle \frac{1}{y(v)}}R^{A\alpha }R^{A\overline{\beta }},`$ (4.4)
which should be equivalent to the metric given by (4.2). Here we have used
$`R^{A\alpha }\underset{\varphi =0}{|}=i\delta ^{A\alpha },R^{A\overline{\alpha }}\underset{\overline{\varphi }=0}{|}=i\delta ^{A\overline{\alpha }},`$ (4.5)
which are obvious by the construction in Subsection 3.1. For other components of the metric we have
$`g^{\alpha \beta }=R^{A\alpha }R^{A\beta }=0,g^{\overline{\alpha }\overline{\beta }}=R^{A\overline{\alpha }}R^{A\overline{\beta }}=0.`$ (4.6)
With (4.4) the Affine connection becomes
$`\mathrm{\Gamma }_{\alpha \beta }^\eta `$ $`=`$ $`g^{\eta \overline{\sigma }}g_{\alpha \overline{\sigma },\beta }=g_{,\beta }^{\eta \overline{\sigma }}g_{\alpha \overline{\sigma }}`$ (4.7)
$`=`$ $`{\displaystyle \frac{1}{y(v)}}R_\beta ^AR_{,\alpha }^{A\eta }={\displaystyle \frac{1}{y(v)}}R_\alpha ^AR_{,\beta }^{A\eta }.`$
by using the property (2.11). Putting this into (4.1) we have
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`g_{\eta \overline{\sigma }}({\displaystyle \frac{1}{y(v)}}R_{\beta ,\overline{\delta }}^AR_{,\alpha }^{A\eta })`$ (4.8)
$`=`$ $`g_{\eta \overline{\sigma }}({\displaystyle \frac{1}{y(v)}}R_{\beta ,\overline{\delta }}^AR_{;\alpha }^{A\eta })`$
$`=`$ $`{\displaystyle \frac{1}{y(v)}}R_{\beta ,\overline{\delta }}^AR_{\overline{\sigma },\alpha }^A.`$
The second equality is due to
$`R_\beta ^AR^{A\eta }=y(v)\delta _\beta ^\eta ,`$ (4.9)
following from (4.4). Multiplying the Lie-algebra (2.14) by $`R_\gamma ^A`$ or $`R^{A\gamma }`$ yields
$`R_{;\gamma }^{B\beta }={\displaystyle \frac{1}{y(v)}}f^{ABC}R^{C\beta }R_\gamma ^A`$ (4.10)
owing to (4.7), or
$`R^{B\alpha }R^{A\gamma }R_{,\alpha }^{A\beta }=f^{ABC}R^{C\beta }R^{A\gamma }.`$ (4.11)
The former is rewritten as
$`R_{\overline{\beta },\gamma }^B={\displaystyle \frac{1}{y(v)}}f^{ABC}R_{\overline{\beta }}^CR_\gamma ^A,`$ (4.12)
while the latter becomes
$`f^{ABC}R^{C\beta }R^{A\gamma }=0,`$ (4.13)
because we have
$`R^{A\gamma }R_{,\alpha }^{A\beta }`$ $`=`$ $`R^{A\gamma }R_{;\alpha }^{A\beta }=g^{\beta \overline{\eta }}R^{A\gamma }R_{\overline{\eta },\alpha }^A`$
$`=`$ $`g^{\beta \overline{\eta }}R^{A\gamma }R_{\alpha ,\overline{\eta }}^A=0,`$
due to (4.6), (2.15) and (4.9). With (4.12) the Riemann curvature (4.8) takes the form
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`{\displaystyle \frac{1}{y(v)^3}}f^{ABE}R_\alpha ^AR_{\overline{\sigma }}^Bf^{CDE}R_\beta ^CR_{\overline{\delta }}^D`$
$`=`$ $`R_{\beta \overline{\sigma }\alpha \overline{\delta }}.`$
The last equality follows from the symmetry of the Affine connection (4.7), or directly from the Jacobi identity of the structure constants
$`f^{ADC}f^{BCE}+f^{BDC}f^{ACE}=f^{ABC}f^{CDE},`$ (4.15)
and (4.13). Contrary to (4.1) this manifests the isometry $`G`$ and includes no derivative with respect to the coordinates. By using it the Riemann curvature can be calculated algebraically, once given a concrete form of the Killing vectors $`R^{A\alpha }`$ which are proper to the Kähler manifold $`G/H`$. Thus (4.14) gives a more practical formula than (4.1) for physical applications.
### 4.2 The reducible case
When the Kähler manifold $`G/H`$ is reducible, the broken generators $`X^i`$ are decomposed into irreducible sets under the subgroup $`H`$. Each irreducible set has a different $`Y`$-charge. The metric (3.15) satisfies the initial condition
$`g_{\alpha \overline{\beta }}\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`g_{\overline{\beta }\alpha }\underset{\varphi =\overline{\varphi }=0}{|}`$ (4.21)
$`=`$ $`\left(\begin{array}{ccccc}y_1(v)& & & & \text{0}\\ & y_2(v)& & & \\ & & & & \\ & & & & \\ \text{0}& & & & y_N(v)\end{array}\right).`$
Therefore the formula (4.4) is no longer correct in this case. We have to generalize the whole arguments in the previous subsection.
First of all, with $`U`$ given by (3.12) and a real symmetric matrix $`P`$ we define the quantity
$`\mathrm{\Delta }=UPU^1`$
in the adjoint representation of the isometry group $`G`$. By (3.17) it transforms as
$`\mathrm{\Delta }(\varphi ^{}\overline{\varphi }^{})=g\mathrm{\Delta }(\varphi ,\overline{\varphi })g^1,`$
or equivalently
$`_^𝒜\mathrm{\Delta }=i[T^A,\mathrm{\Delta }],`$ (4.22)
if $`P`$ satisfies
$`hPh^1=P.`$ (4.23)
With this $`\mathrm{\Delta }`$ the metric $`g^{ab}`$ is found as a solution to the Killing condition (2.7)
$`g^{ab}`$ $`=`$ $`g^{ba}=R^{Aa}\mathrm{\Delta }^{AB}R^{Bb}R^a\mathrm{\Delta }R^b.`$ (4.24)
In the complex basis it reads
$`g^{\alpha \overline{\beta }}`$ $`=`$ $`g^{\overline{\beta }\alpha }=R^\alpha \mathrm{\Delta }R^{\overline{\beta }},`$
$`g^{\alpha \beta }`$ $`=`$ $`R^\alpha \mathrm{\Delta }R^\beta ,`$ (4.25)
$`g^{\overline{\alpha }\overline{\beta }}`$ $`=`$ $`R^{\overline{\alpha }}\mathrm{\Delta }R^{\overline{\beta }}.`$
We now assume the real symmetric matrix $`P`$ to have non-vanishing elements only in the diagonal blocks corresponding to the broken generators $`X^a=(X^{\overline{i}},X^i)`$ such that
$`P^{i\overline{j}}=P^{\overline{j}i}=\left(\begin{array}{ccccc}y_1(v)^1& & & & \text{0}\\ & y_2(v)^1& & & \\ & & & & \\ & & & & \\ \text{0}& & & & y_N(v)^1\end{array}\right).`$ (4.31)
Then $`P`$ satisfies (4.18) because the diagonal elements are decomposed into irreducible sets under the subgroup $`H`$ by the $`Y`$-charge. Evaluate these metrics in (4.20) at the origin of the coset space by (4.5) and (4.21). We find that they all satisfy the same initial conditions as the metrics given in (3.15). Thus both metrics are equivalent<sup>2</sup><sup>2</sup>2The equivalence is alternatively stated as
$$e^a=(R_\alpha UP)^ad\varphi ^\alpha +(R_{\overline{\alpha }}UP)^ad\varphi ^{\overline{\alpha }},$$
with $`e^a=(e^{\overline{i}},e^i)`$ defined by the Cartan-Maurer $`1`$-form (3.14). Namely both sides have the same Lie-derivatives with respect to $`R^{Aa}`$ and the same values at the origin of the coset space. The author is indebted to T. Kugo for the discusion on this comment , and we have
$`R^\alpha \mathrm{\Delta }R^\beta =0,R^{\overline{\alpha }}\mathrm{\Delta }R^{\overline{\beta }}=0.`$ (4.32)
This generalization of the metric requires to modify the formula (4.7)$``$
(4.14) in the previous subsection. Rewrite (4.20) as
$`R^\alpha \mathrm{\Delta }R_\beta `$ $`=`$ $`\delta _\beta ^\alpha ,R^{\overline{\alpha }}\mathrm{\Delta }R_{\overline{\beta }}=\delta _{\overline{\beta }}^{\overline{\alpha }},`$
$`R^\alpha \mathrm{\Delta }R_{\overline{\beta }}`$ $`=`$ $`0,R^{\overline{\alpha }}\mathrm{\Delta }R_\beta =0,`$ (4.33)
using (4.22). Differentiate them by the coordinates to find
$`R_{\overline{\alpha }}(\mathrm{\Delta }R_\beta )_{,\overline{\gamma }}=0,R_\alpha (\mathrm{\Delta }R_{\overline{\beta }})_{,\gamma }`$ $`=`$ $`0,`$ (4.34)
$`R_{\overline{\alpha }}(\mathrm{\Delta }R_{\overline{\beta }})_{,\overline{\gamma }}=0,R_\alpha (\mathrm{\Delta }R_\beta )_{,\gamma }`$ $`=`$ $`0.`$ (4.35)
With the metric (4.20) the Affine connection (4.7) changes the form as
$`\mathrm{\Gamma }_{\alpha \beta }^\gamma `$ $`=`$ $`R_{\alpha ,\beta }\mathrm{\Delta }R^\gamma =R_\alpha (\mathrm{\Delta }R^\gamma )_{,\beta },`$
owing to (4.23) and (4.24). Then the Riemann curvature becomes
$`R_{\alpha \beta \overline{\delta }}^\gamma `$ $`=`$ $`\mathrm{\Gamma }_{\alpha \beta ,\overline{\delta }}^\gamma `$ (4.36)
$`=`$ $`[R_{\alpha ,\overline{\delta }}(\mathrm{\Delta }R^\gamma )_{,\beta }+R_\alpha (\mathrm{\Delta }R^\gamma )_{,\beta \overline{\delta }}]`$
$`=`$ $`[(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{,\beta }^\gamma +(R_\alpha \mathrm{\Delta }_{,\beta })_{,\overline{\delta }}R^\gamma ].`$
By (4.24) and (2.15) the first piece changes the form as
$`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{,\beta }^\gamma `$ $`=`$ $`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{;\beta }^\gamma `$
$`=`$ $`g^{\gamma \overline{\sigma }}(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{\overline{\sigma },\beta }`$
$`=`$ $`g^{\gamma \overline{\sigma }}(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{\beta ,\overline{\sigma }}.`$
By means of the formulae (A.7) in the Appendix A it becomes
$`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_{,\beta }^\gamma `$ $`=`$ $`g^{\gamma \overline{\sigma }}\{(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R^{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$ (4.37)
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}^A(R_{\overline{\sigma }}\mathrm{\Delta })^BR_\beta ^C\}.`$
On the other hand the second piece of (4.26) is calculated as
$`(R_\alpha \mathrm{\Delta }_{,\beta })_{,\overline{\delta }}R^\gamma `$ $`=`$ $`g^{\gamma \overline{\sigma }}\{f^{ABC}(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}^AR_\beta ^B(R_{\overline{\sigma }}\mathrm{\Delta })^C`$
$`+f^{ABC}(R_\alpha \mathrm{\Delta })^AR_{\beta ,\overline{\delta }}^B(R_{\overline{\sigma }}\mathrm{\Delta })^C`$
$`+f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }})^C\},`$
by means of (A.10). Putting (4.27) and (4.28) together into (4.26) we have
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R^{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_{\beta ,\overline{\delta }}^B(R_{\overline{\sigma }}\mathrm{\Delta })^C`$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }})^C.`$
Calculate the first piece further as
$`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R^{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$ $`=`$ $`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R_\rho g^{\rho \overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$ (4.40)
$`=`$ $`R_\alpha \mathrm{\Delta }R_{\overline{\delta },\rho }g^{\rho \overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$
$`=`$ $`R_\alpha \mathrm{\Delta }_{,\rho }R_{\overline{\delta }}g^{\rho \overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta ,`$
in which the second equality is obtained by (4.22) and (2.15), while the third one by (4.24). Rewrite the last piece of (4.29) as
$`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }})^C`$ $`=`$ $`R_\alpha \mathrm{\Delta }_{,\beta }R_{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }}R^{\overline{\eta }},`$ (4.41)
by (A.8). By means of (A.7), (A.9) and (A.10) the Riemann curvature (4.29) turns out to be
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}^F`$ (4.42)
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\delta }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}^F`$
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^Bf^{CDE}R_\beta ^D(R_{\overline{\delta }}\mathrm{\Delta })^E`$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\eta }}\mathrm{\Delta })^Cf^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^DR_{\overline{\delta }}^E(R^{\overline{\eta }}\mathrm{\Delta })^F.`$
This is the generalization of (4.14). If one replaces $`\mathrm{\Delta }^{AB}`$ by $`\frac{1}{y(v)}\delta ^{AB}`$, then (4.32) reduces to (4.14) due to the formula (4.13) which is only valid for the irreducible case. At this final stage it is worth showing the symmetry property
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`R_{\beta \overline{\sigma }\alpha \overline{\delta }}=R_{\alpha \overline{\delta }\beta \overline{\sigma }},`$ (4.43)
as a consistency check of our calculations. The demonstration will be given in Appendix B. There we also show that the formula (4.32) takes an alternative form such that
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}^F`$
$`+`$ $`f^{ABC}(R_\beta \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\alpha \mathrm{\Delta })^ER_{\overline{\eta }}^F`$
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^Bf^{CDE}(R_\beta \mathrm{\Delta })^DR_{\overline{\delta }}^E`$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\eta }}\mathrm{\Delta })^Cf^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^DR_{\overline{\delta }}^E(R^{\overline{\eta }}\mathrm{\Delta })^F.`$
## 5 Applications
In the $`N=1`$ supersymmetric non-linear $`\sigma `$-model on the Kähler coset space $`G/H`$, the four-fermi coupling is the most interesting part. When the coset space is reducible, the Riemann curvature depends on the $`Y`$-charge of the broken generators through the metric (3.15). It takes the form (4.32) which is rather complicated than that of the irreducible coset space. On top of this complication we have another one, if the homogeneous subgroup $`H`$ contains plural $`U(1)`$s as $`H=H^{}\{U(1)\}^k`$. Namely, the splitting of the broken generators $`X^{\overline{i}}`$ and $`X^i`$ depends on the constants $`v^\mu `$ of the $`Y`$-charge (3.2), so that we may have different sets of the NG bosons . Of course a phenomelogically interesting set should be chosen. Then the four-fermi coupling depends on the $`Y`$-charges of the broken generators $`X^{\overline{i}}`$ in a piculiar way to the choice. It is very interesting from the phenomelogical point of view.
As has been explained in the introduction, the most important part of the four-fermi coupling is
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}(\overline{\psi }^{\overline{\sigma }}\psi ^\alpha )(\overline{\psi }^{\overline{\delta }}\psi ^\beta )`$ (5.1)
in the non-linear $`\sigma `$-model as a low-energy effective theory. We shall present a systematic method to evaluate the four-fermi coupling constants $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}|_{\varphi =\overline{\varphi }=0}`$ by means of (4.32). The method will enable us to fully controle the $`Y`$-charge dependence of $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}|_{\varphi =\overline{\varphi }=0}`$.
By (4.5) and (4.21) we note at first that
$`R_\alpha ^A\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`iy_\alpha (v)\delta _\alpha ^A,R_{\overline{\alpha }}^A\underset{\varphi =\overline{\varphi }=0}{|}=iy_\alpha (v)\delta _{\overline{\alpha }}^A.`$
and
$`(R_\alpha \mathrm{\Delta })^A\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`i\delta _\alpha ^A,(R^{\overline{\alpha }}\mathrm{\Delta })^A={\displaystyle \frac{i}{y_\alpha (v)}}\delta ^{\overline{\alpha }A}`$
$`(R_{\overline{\alpha }}\mathrm{\Delta })^A\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`i\delta _{\overline{\alpha }}^A,(R^\alpha \mathrm{\Delta })^A={\displaystyle \frac{i}{y_\alpha (v)}}\delta ^{\alpha A}`$
By using this (4.32) becomes
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}(f^{\overline{\alpha }\delta \overline{\eta }}f^{\overline{\beta }\sigma \eta }+f^{\overline{\alpha }\sigma \overline{\eta }}f^{\overline{\beta }\delta \eta })y_\eta (v)`$ (5.2)
$`+`$ $`[{\displaystyle \underset{\eta =1}{\overset{N}{}}}(f^{\overline{\alpha }\sigma \eta }f^{\overline{\beta }\delta \overline{\eta }}+f^{\overline{\alpha }\sigma \overline{\eta }}f^{\overline{\beta }\delta \eta })y_\beta (v)+{\displaystyle \underset{C=I,\mu }{\overset{\mathrm{dim}H}{}}}f^{\overline{\alpha }\sigma C}f^{\overline{\beta }\delta C}y_\beta (v)]`$
$``$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}f^{\overline{\alpha }\overline{\beta }\eta }f^{\sigma \delta \overline{\eta }}{\displaystyle \frac{y_\beta (v)y_\delta (v)}{y_\eta (v)}}.`$
Each piece of (5.2) can be computed by means of
$`f^{ABC}f^{CDE}`$ $`=`$ $`{\displaystyle \frac{1}{2𝒩}}\mathrm{tr}([T^A,T^B][T^D,T^E]),`$ (5.3)
with the normalization $`tr(T^AT^B)=2𝒩\delta ^{AB}`$. The Riemann curvature appears with the indices $`\alpha ,\overline{\sigma },\beta ,\overline{\delta }`$ of the three types:
$`y([T^{\overline{\alpha }},T^\sigma ])`$ $`>`$ $`0,y([T^{\overline{\beta }},T^\delta ])<0,`$ (5.4)
$`y([T^{\overline{\alpha }},T^\sigma ])`$ $`<`$ $`0,y([T^{\overline{\beta }},T^\delta ])>0,`$ (5.5)
$`y([T^{\overline{\alpha }},T^\sigma ])`$ $`=`$ $`0,y([T^{\overline{\beta }},T^\delta ])=0,`$ (5.6)
in which $`y([,])`$ is the $`Y`$-charge of the commutator. (4.32) reads
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}f^{\overline{\alpha }\sigma \eta }f^{\overline{\beta }\delta \overline{\eta }}y_\beta (v)`$ (5.7)
$``$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}[f^{\overline{\alpha }\delta \overline{\eta }}f^{\overline{\beta }\sigma \eta }y_\eta (v)+f^{\overline{\alpha }\overline{\beta }\eta }f^{\sigma \delta \overline{\eta }}{\displaystyle \frac{y_\beta (v)y_\delta (v)}{y_\eta (v)}}]`$
for the case (5.4),
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}f^{\overline{\alpha }\sigma \overline{\eta }}f^{\overline{\beta }\delta \eta }(y_\beta (v)y_\eta (v))`$ (5.8)
$``$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}[f^{\overline{\alpha }\delta \overline{\eta }}f^{\overline{\beta }\sigma \eta }y_\eta (v)+f^{\overline{\alpha }\overline{\beta }\eta }f^{\sigma \delta \overline{\eta }}{\displaystyle \frac{y_\beta (v)y_\delta (v)}{y_\eta (v)}}]`$
for the case (5.5), and
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}`$ $`=`$ $`{\displaystyle \underset{C=I,\mu }{\overset{\mathrm{dim}H}{}}}f^{\overline{\alpha }\sigma C}f^{\overline{\beta }\delta C}y_\beta (v)`$ (5.9)
$``$ $`{\displaystyle \underset{\eta =1}{\overset{N}{}}}[f^{\overline{\alpha }\delta \overline{\eta }}f^{\overline{\beta }\sigma \eta }y_\eta (v)+f^{\overline{\alpha }\overline{\beta }\eta }f^{\sigma \delta \overline{\eta }}{\displaystyle \frac{y_\beta (v)y_\delta (v)}{y_\eta (v)}}]`$
for the case (5.6). (5.7) can be alternatively obtained by applying the symmetry property (4.33) to (5.8).
### 5.1 $`SU(3)/\{U(1)\}^2`$
We start with the most simplest reducible coset space $`SU(3)/\{U(1)\}^2`$ to illustrate our basic strategy. The generators of $`SU(3)`$ are
$`\{T^A\}`$ $`=`$ $`\{T_2^1,T_3^1,T_3^2,T_1^2,T_1^3,T_2^3,Q,Q^{}\},`$
with
$`Q=\frac{1}{\sqrt{2}}(T_1^1T_2^2),Q^{}=\sqrt{\frac{3}{2}}T_3^3,`$ (5.10)
and the hermitian condition $`(T_j^i)^{}=T_j^i`$. They satisfy the Lie-algebra
$`[T_i^j,T_k^l]=\delta _k^jT_i^l\delta _i^lT_k^j.`$
The quadratic Casimir takes the form
$`\{T_2^1,T_1^2\}+\{T_3^1,T_1^3\}+\{T_3^2,T_2^3\}+Q^2+Q^2,`$
from which we read the Killing metric $`\delta ^{AB}`$ in (2.18). The $`U(1)`$-charges $`Q`$ and $`Q^{}`$ of the broken generators $`T_i^j(ij)`$ as well as their $`Y`$-charges
$`Y`$ $`=`$ $`vQ+v^{}Q^{}`$
are given in Table 1. By means of them the broken generators are splitted in
two parts: the generators $`X^i`$ with positive $`Y`$-charge and their hermitian conjugates $`X^{\overline{i}}`$ with negative charge. For illustration we plot the broken generators in the $`(Q,Q^{})`$-charge plane in Figure 1. There are three possibilities to draw the line : $`Y=0`$,
each of which gives a different splitting:
$`\mathrm{I}\{X^i\}`$ $`=`$ $`\{T_1^2,T_1^3,T_2^3\},`$
$`\mathrm{II}\{X^i\}`$ $`=`$ $`\{T_1^3,T_2^3,T_2^1\},`$
$`\mathrm{III}\{X^i\}`$ $`=`$ $`\{T_2^3,T_2^1,T_3^1\}.`$
Taking the case I we proceed with the argument. With the identification $`X^1=T_1^2,X^2=T_1^3,X^3=T_2^3,Q=X^q,Q^{}=X^q^{}`$, the non-trivial part of the Lie-algebra reads
$`[X^2,X^{\overline{3}}]`$ $`=`$ $`X^1,[X^3,X^{\overline{2}}]=X^{\overline{1}},`$
$`[X^1,X^3]`$ $`=`$ $`X^2,[X^{\overline{3}},X^{\overline{1}}]=X^{\overline{2}},`$
$`[X^{\overline{1}},X^2]`$ $`=`$ $`X^3,[X^{\overline{2}},X^1]=X^{\overline{3}},`$
$`[X^1,X^{\overline{1}}]`$ $`=`$ $`\sqrt{2}Q,`$
$`[X^2,X^{\overline{2}}]`$ $`=`$ $`\frac{1}{\sqrt{2}}Q+\sqrt{\frac{3}{2}}Q^{},`$ (5.11)
$`[X^3,X^{\overline{3}}]`$ $`=`$ $`\frac{1}{\sqrt{2}}Q+\sqrt{\frac{3}{2}}Q^{},`$
$`[X^q,X^1]`$ $`=`$ $`\sqrt{2}X^1,[X^q^{},X^1]=0,`$
$`[X^q,X^2]`$ $`=`$ $`\frac{1}{\sqrt{2}}X^2,[X^q^{},X^2]=\sqrt{\frac{3}{2}}X^2,`$
$`[X^q,X^3]`$ $`=`$ $`\frac{1}{\sqrt{2}}X^3,[X^q^{},X^3]=\sqrt{\frac{3}{2}}X^3,`$
and their hermitian conjugates. The holomorphic Killing vectors $`R^{A\alpha }`$ are easily obtained by studying (3.6) in the fundamental representation:
$$\begin{array}{ccc}R^{\overline{1}1}=i,\hfill & R^{11}=i(\varphi ^1)^2\hfill & R^{q1}=\sqrt{2}i\varphi ^1,\hfill \\ R^{\overline{2}1}=0,\hfill & R^{21}=i\varphi ^1(\varphi ^2+\frac{1}{2}\varphi ^1\varphi ^3),\hfill & R^{q^{}1}=0,\hfill \\ R^{\overline{3}1}=0,\hfill & R^{31}=i(\varphi ^2+\frac{1}{2}\varphi ^1\varphi ^3),\hfill & \\ R^{\overline{1}2}=\frac{i}{2}\varphi ^3,\hfill & R^{12}=\frac{i}{2}\varphi ^1(\varphi ^2+\frac{1}{2}\varphi ^1\varphi ^3),\hfill & R^{q2}=\frac{i}{\sqrt{2}}\varphi ^2,\hfill \\ R^{\overline{2}2}=i,\hfill & R^{22}=i[(\varphi ^2)^2+\frac{1}{4}(\varphi ^1\varphi ^3)^2],\hfill & R^{q^{}2}=\sqrt{\frac{3}{2}}i\varphi ^2,\hfill \\ R^{\overline{3}2}=\frac{i}{\sqrt{2}}\varphi ^1,\hfill & R^{32}=\frac{i}{2}\varphi ^1(\varphi ^2\frac{1}{2}\varphi ^1\varphi ^3),\hfill & \\ R^{\overline{1}3}=0,\hfill & R^{13}=i(\varphi ^2\frac{1}{2}\varphi ^1\varphi ^3),\hfill & R^{q3}=\frac{i}{\sqrt{2}}\varphi ^3,\hfill \\ R^{\overline{2}3}=0,\hfill & R^{23}=i\varphi ^1(\varphi ^2\frac{1}{2}\varphi ^1\varphi ^3),\hfill & R^{q^{}3}=\sqrt{\frac{3}{2}}i\varphi ^3,\hfill \\ R^{\overline{3}3}=i,\hfill & R^{33}=i(\varphi ^3)^2,\hfill & \end{array}$$
The Riemann curvature $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}|_{\varphi =\overline{\varphi }=0}(G^{\overline{\alpha }\sigma \overline{\beta }\delta }`$), given by (5.2), is calculated by using (5.3) with the commutators (5.11). The Riemann curvature of this type appears in the coset space $`E_7/SU(5)\{U(1)\}^3`$ which we will study in the next subsection. The result is given in Tables 8 and 9 there.
### 5.2 $`E_7/SU(5)\{U(1)\}^3`$
The generators of $`E_7`$ are decomposed as
$`\{T^A\}=\{E_i^{ab},T_a^i,E^a,E_{ab}^i,T_i^a,E_a,T_a^b,T_i^j,T\}`$ (5.12)
in the basis of the subgroup $`SU(5)SU(3)U(1)`$. Here $`a,b,\mathrm{}`$ and $`i,j,\mathrm{}`$ are indices of $`SU(5)`$ and $`SU(3)`$ running over $`15`$ and $`13`$ respectively. They have $`SU(5)SU(3)`$ quantum numbers
$`(\mathrm{𝟏𝟎},\mathrm{𝟑}^{}),(\mathrm{𝟓}^{},\mathrm{𝟑}),(\mathrm{𝟓},\mathrm{𝟏}),(\mathrm{𝟏𝟎}^{},\mathrm{𝟑}),(\mathrm{𝟓},\mathrm{𝟑}^{}),(\mathrm{𝟓}^{},\mathrm{𝟏}),(\mathrm{𝟐𝟒},\mathrm{𝟏}),(\mathrm{𝟏},\mathrm{𝟖}),(\mathrm{𝟏},\mathrm{𝟏}),`$
in the order of (5.12). The non-trivial part of the $`E_7`$ algebra takes the form
$`[E_i^{ab},E_j^{cd}]=\epsilon _{ijk}\epsilon ^{abcde}T_e^k,[E^a,E_j^{bc}]=0,`$
$`[T_a^i,E_j^{bc}]=\delta _j^i(\delta _a^bE^c\delta _a^cE^b),[T_a^i,E^b]=0,`$
$`[E^a,E^b]=0,`$
$`[E_{ab}^i,E_j^{cd}]=\delta _j^i(\delta _a^cT_b^d+\delta _b^dT_a^c\delta _b^cT_a^d\delta _a^dT_b^c)\delta _{ab}^{cd}(T_j^i+\sqrt{\frac{2}{15}}\delta _j^iT),`$
$`[E^a,E_b]=\sqrt{\frac{6}{5}}\delta _b^aTT_b^a,`$ (5.13)
$`[E^a,E_{bc}^i]=\delta _c^aT_b^i\delta _b^aT_c^i,`$
$`[T_a^i,T_j^b]=\delta _a^bT_j^i+\delta _j^iT_a^b+2\sqrt{\frac{2}{15}}\delta _j^i\delta _a^bT,`$
$`[T_a^i,E_{bc}^j]=\frac{1}{2}\epsilon ^{ijk}\epsilon _{abcde}E_k^{de},`$
$`[T_a^i,E_b]=E_{ab}^i,\mathrm{h}.\mathrm{c}.,`$
with
$$\epsilon _{ijk}^{}=\epsilon ^{ijk},\epsilon _{abcde}^{}=\epsilon ^{abcde},\delta _{ab}^{cd}=\delta _a^c\delta _b^d\delta _a^d\delta _b^c.$$
In the basis of the smaller subgroup $`SU(5)\{U(1)\}^3`$ the generators are further decomposed as
$`\{T^A\}=\{X^{\overline{i}},X^i,S^I,Q^\mu \},`$ (5.14)
with
$`\{X^{\overline{i}}\}`$ $`=`$ $`\{E_i^{ab},E^a,T_a^i,T_i^j(i>j)\},`$
$`\{X^i\}`$ $`=`$ $`\{E_{ab}^i,E_a,T_i^a,T_i^j(i<j)\},`$
$`\{S^I\}`$ $`=`$ $`\{T_a^b(_aT_a^a=0)\},`$ (5.15)
$`\{Q^\mu \}`$ $`=`$ $`\{T,Q,Q^{}\}.`$
Here $`Q`$ and$`Q^{}`$ are the $`U(1)`$s contained in $`SU(3)`$ which was given by (5.10). The quadratic Casimir takes the form
$`C`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{E_i^{ab},E_{ab}^i\}+\{T_i^a,T_a^i\}+\{E^a,E_a\}`$
$`+`$ $`\{T_2^1,T_1^2\}+\{T_3^1,T_1^3\}+\{T_3^2,T_2^3\}`$
$`+`$ $`\{T_a^b,T_b^a\}+T^2+Q^2+Q^2.`$
$`\{X^{\overline{i}}\}`$ and $`\{X^i\}`$ are broken generators of $`E_7/SU(5)\{U(1)\}^3`$. In the supersymmetric model on $`E_7/SU(5)\{U(1)\}^3`$ there are pseud NG fermions corresponding to them. Among of them the pseud NG fermions having the same $`SU(5)SU(3)`$ quantum numbers as $`E_i^{ab},T_a^i`$ and $`T_i^j(i>j)`$ are identified with the three families of quarks and leptons $`\psi ^{\genfrac{}{}{0pt}{}{ab}{i}},\psi ^{\genfrac{}{}{0pt}{}{i}{a}}`$ and $`\psi ^{\genfrac{}{}{0pt}{}{j}{i}}`$. The $`Y`$-charge is made of the $`U(1)`$ charges in $`\{Q^\mu \}`$:
$`Y=\alpha T+\beta Q+\gamma Q^{}.`$ (5.16)
These $`U(1)`$-charges are given in Table 3.
The splitting of the broken generators $`\{X^{\overline{i}}\}`$ and $`\{X^i\}`$ changes according to the orientation of the plane $`Y=0`$ in the $`(T,Q,Q^{})`$-charge space. The splitting (5.15) is valid only when the vector coefficients $`(\alpha ,\beta ,\gamma )`$ are chosen such that
$`y_i(v)Y(X^{\overline{i}})>0,\mathrm{for}\mathrm{all}\overline{i},`$
for instance,
$`\alpha =1,\beta <0,\gamma <0,|\beta |=|\gamma |<<1.`$
We proceed with the argument in this special splitting, since the pseud NG fermions are then neatly identified with the three families of quarks and leptons.
In the supersymmetric $`\sigma `$-model on $`E_7/SU(5)\{U(1)\}^3`$ the Riemann curvature $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ is a $`SU(5)`$-covariant tensor. We shall be interested in the four-fermi coupling of the three families of quarks and leptons alone. Then the relevant part of the Riemann curvature appears with the $`SU(5)`$-content of the following types:
$$R_{\alpha \overline{\sigma }\beta \overline{\delta }}\underset{\varphi =\overline{\varphi }=0}{|}G^{\overline{\alpha }\sigma \overline{\beta }\delta }\{\begin{array}{cc}(\mathrm{𝟓}^{},\mathbf{5},\mathbf{5}^{},\mathbf{5}),\hfill & \\ (\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}),\hfill & \\ (\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟏}),\hfill & \\ (\mathrm{𝟓}^{},\mathbf{5},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}),\hfill & (\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{},\mathrm{𝟓}^{},\mathbf{5}),\hfill \\ (\mathrm{𝟓},\mathrm{𝟏𝟎},\mathrm{𝟓}^{},\mathbf{10}^{}),\hfill & (\mathrm{𝟓}^{},\mathbf{10}^{},\mathrm{𝟓},\mathrm{𝟏𝟎}),\hfill \\ (\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟓}^{},\mathrm{𝟓}),\hfill & (\mathrm{𝟓}^{},\mathrm{𝟓},\mathrm{𝟏},\mathrm{𝟏}),\hfill \\ (\mathrm{𝟏},\mathrm{𝟓},\mathrm{𝟓}^{},\mathrm{𝟏}),\hfill & (\mathrm{𝟓}^{},\mathrm{𝟏},\mathrm{𝟓},\mathrm{𝟏}),\hfill \\ (\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}),\hfill & (\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{},\mathrm{𝟏},\mathrm{𝟏}),\hfill \\ (\mathrm{𝟏},\mathrm{𝟏𝟎}^{},\mathrm{𝟏𝟎},\mathrm{𝟏}),\hfill & (\mathrm{𝟏𝟎},\mathrm{𝟏},\mathrm{𝟏},\mathrm{𝟏𝟎}^{}).\hfill \end{array}$$
We evaluate $`G^{\overline{\alpha }\sigma \overline{\beta }\delta }`$ in components by means of (5.2) with (5.3) and (5.13). (See Appendix C.) The results are summarized in Tables 3$``$16 for six types of the Riemann curvature
$$\begin{array}{ccc}(\mathrm{𝟓}^{},\mathbf{5},\mathbf{5}^{},\mathbf{5}),\hfill & (\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}),\hfill & (\mathrm{𝟏},\mathbf{1},\mathbf{1},\mathbf{1}),\hfill \\ (\mathrm{𝟓}^{},\mathbf{5},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}),\hfill & (\mathbf{1},\mathbf{1},\mathbf{5}^{},\mathbf{5}),\hfill & (\mathrm{𝟏},\mathbf{1},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{}).\hfill \end{array}$$
Other types of the Riemann curvature are obtained by applying the symmetry property (4.33) to these results.
## 6 Conclusions
In this paper we have discussed the reducible Kähler coset space $`G/S\{U(1)\}^k`$ in the geometrical approach generalizing the arguments in ref.. We have expressed the Riemann curvature of the coset space in terms of the Killing vectors as (4.32). It is the most important formula in this paper. We have been then interested in the four-fermi coupling of the supersymmetric non-linear $`\sigma `$-model on $`G/S\{U(1)\}^k`$, to the leading order of $`\frac{1}{f}`$. It is given by evaluating the Riemann curvature at the origin of the coset space. We have established the group theoretical method to do this by using the formula (4.32). Otherwise the calculation would be too complicated. Concrete calculations have been done for $`SU(3)/\{U(1)\}^2`$ and $`E_7/SU(5)\{U(1)\}^3`$. The results of the last Kähler coset space is phenomelogically interesting, since they give four-fermi coupling constants among the three families of $`\mathrm{𝟏𝟎}+\mathrm{𝟓}^{}+\mathrm{𝟏}`$ of $`SU(5)`$ in the supersymmetric non-linear $`\sigma `$-model on $`E_7/SU(5)\{U(1)\}^3`$. Among them those involving the three families of right-handed neutrinos are particularly interesting and have been given in Tables 8,9,13$``$16. The dependence of the three $`U(1)`$-charges of the NG pseudo fermions are explicit in these results. Of course, we may take another set of $`U(1)`$s, say $`Q^1,Q^2,Q^3`$, than $`T,Q,Q^{}`$, for instance, those which remain unbroken in the breaking process
$$E_7\stackrel{U(1)}{}E_6\stackrel{U(1)}{}SO(10)\stackrel{U(1)}{}SU(5)$$
as in ref.. The results given by Tables 3$``$16 are still valid if one defines the the Y-charge as
$$Y=\alpha Q^1+\beta Q^2+\gamma Q^3.$$
and replaces Table 2 for $`y(X^{\overline{i}})`$ by a new table with $`Q^1,Q^2,Q^3`$. It is desired to carry out a phenomelogical study by tuning the three arbitrary parameters $`\alpha ,\beta ,\gamma `$.
Acknowledgements
The author would like to thank T. Yanagida for reviving the interest in the subject. He is grateful to T. Kugo for the valuable discussions on the reducible Kähler coset space and reading the manuscript.
## Appendix A
We derive the useful formulae for the calculation in Subsection 4.2. We start with covariantization of the Lie-algebra (2.14):
$$R^{A\alpha }R_{;\alpha }^{B\beta }R^{B\alpha }R_{;\alpha }^{A\beta }=f^{ABC}R^{C\beta }.$$
Multiplying both sides by $`(R_\gamma \mathrm{\Delta })^B`$ or $`(R_{\overline{\gamma }}\mathrm{\Delta })^B`$ and lower the index $`\beta `$ to get
$`R^{A\alpha }R_\gamma \mathrm{\Delta }R_{\overline{\beta },\alpha }R_{\overline{\beta },\gamma }^A`$ $`=`$ $`f^{ABC}(R_\gamma \mathrm{\Delta })^BR_{\overline{\beta }}^C,`$ (A.1)
or
$`R^{A\alpha }R_{\overline{\gamma }}\mathrm{\Delta }R_{\overline{\beta },\alpha }`$ $`=`$ $`f^{ABC}(R_{\overline{\gamma }}\mathrm{\Delta })^BR_{\overline{\beta }}^C,`$ (A.2)
by using (4.23). Noting that
$`R_\gamma \mathrm{\Delta }R_{\overline{\beta },\alpha }`$ $`=`$ $`R_\gamma \mathrm{\Delta }_{,\alpha }R_{\overline{\beta }},`$ (A.3)
and
$`R_{\overline{\gamma }}\mathrm{\Delta }R_{\overline{\beta },\alpha }=R_{\overline{\gamma }}\mathrm{\Delta }R_{\alpha ,\overline{\beta }}=R_{\overline{\gamma }}\mathrm{\Delta }_{,\overline{\beta }}R_\alpha ,`$ (A.4)
by (4.24) and (2.15), we write (A.1) and (A.2) respectively as
$`R_{\overline{\beta },\gamma }^A`$ $`=`$ $`R^{A\alpha }R_\gamma \mathrm{\Delta }_{,\alpha }R_{\overline{\beta }}f^{ABC}(R_\gamma \mathrm{\Delta })^BR_{\overline{\beta }}^C`$ (A.5)
and
$`R^{A\alpha }R_{\overline{\gamma }}\mathrm{\Delta }_{,\overline{\beta }}R_\alpha `$ $`=`$ $`f^{ABC}(R_{\overline{\gamma }}\mathrm{\Delta })^BR_{\overline{\beta }}^C.`$ (A.6)
Taking the complex conjugation of them gives
$`R_{\gamma ,\overline{\beta }}^A`$ $`=`$ $`R^{A\overline{\alpha }}R_{\overline{\beta }}\mathrm{\Delta }_{,\overline{\alpha }}R_\gamma f^{ABC}(R_{\overline{\beta }}\mathrm{\Delta })^BR_\gamma ^C`$ (A.7)
and
$`R^{A\overline{\alpha }}R_\gamma \mathrm{\Delta }_{,\beta }R_{\overline{\alpha }}`$ $`=`$ $`f^{ABC}(R_\gamma \mathrm{\Delta })^BR_\beta ^C.`$ (A.8)
That (A.5) and (A.7) satisfy the Killing condition (2.15) can be easily checked by (4.17), i.e.,
$$R^{A\alpha }\mathrm{\Delta }_{,\alpha }+R^{A\overline{\alpha }}\mathrm{\Delta }_{,\overline{\alpha }}=i[T^A,\mathrm{\Delta }],$$
with $`(T^A)^{BC}=if^{ABC}`$. Multiplying both sides of (A.6) and (A.8) respectively by $`(R_\eta \mathrm{\Delta })^A`$ and $`(R_{\overline{\eta }}\mathrm{\Delta })^A`$ we get
$`R_{\overline{\gamma }}\mathrm{\Delta }_{,\overline{\beta }}R_\eta `$ $`=`$ $`f^{ABC}(R_{\overline{\gamma }}\mathrm{\Delta })^AR_{\overline{\beta }}^B(R_\eta \mathrm{\Delta })^C`$ (A.9)
and
$`R_\gamma \mathrm{\Delta }_{,\beta }R_{\overline{\eta }}`$ $`=`$ $`f^{ABC}(R_\gamma \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\eta }}\mathrm{\Delta })^C`$ (A.10)
owing to (4.23).
## Appendix B
We will check the symmetry property $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}=R_{\alpha \overline{\delta }\beta \overline{\sigma }}`$ of the r.h.s. of (4.32). Put it in the form
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`[f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}`$ (B.1)
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\delta }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}^F]`$
$`+`$ $`[f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^Bf^{CDE}R_\beta ^D(R_{\overline{\delta }}\mathrm{\Delta })^E`$
$``$ $`R_\alpha \mathrm{\Delta }_{,\beta }R_{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }}R^{\overline{\eta }}],`$
remembering (4.31). The first bracket is already symmetric under the interchange of $`\overline{\sigma }`$ and $`\overline{\delta }`$. Therefore we are left with the second bracket to exsamine. The anti-symmetric sum of its first piece by interchanging $`\overline{\sigma }`$ and $`\overline{\delta }`$ becomes
$`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^Bf^{CDE}R_\beta ^D(R_{\overline{\delta }}\mathrm{\Delta })^E`$
$`f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^Bf^{CDE}R_\beta ^D(R_{\overline{\sigma }}\mathrm{\Delta })^E`$
$`=f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^Bf^{CDE}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_{\overline{\delta }}\mathrm{\Delta })^E,`$ (B.2)
by using the Jacobi identity of the structure constants (4.15). On the other hand that of the second piece is given by
$`(R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }}R_\eta R_{\overline{\delta }}\mathrm{\Delta }_{,\overline{\sigma }}R_\eta )R_\alpha \mathrm{\Delta }_{,\beta }R^\eta `$
$`=f^{ABC}(R_{\overline{\sigma }}\mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^Bf^{CDE}(R_\alpha \mathrm{\Delta })^DR_\beta ^E.`$ (B.3)
This is easily checked as follows. Note at first that
$$R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }}R_\eta +R_{\overline{\sigma }}\mathrm{\Delta }R_{\eta ,\overline{\delta }}=0,$$
from (4.24). Then plug (A.7) in this to find
$$R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }}R_\eta R_{\overline{\delta }}\mathrm{\Delta }_{,\overline{\sigma }}R_\eta =f^{ABC}(R_{\overline{\sigma }}\mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR_\eta ^C.$$
Multiplying both sides by $`R_\alpha \mathrm{\Delta }_{,\beta }R^\eta `$ and using (A.8) yields (B.3). From (B.2) and (B.3) the second bracket of (B.1) is also symmetric under the interchange of $`\overline{\sigma }`$ and $`\overline{\delta }`$. Thus we have $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}=R_{\alpha \overline{\delta }\beta \overline{\sigma }}`$.
Next we examine the symmetry property $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}=R_{\beta \overline{\sigma }\alpha \overline{\delta }}`$. We rewrite the Riemann curvature (4.29) as
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`(R_\alpha \mathrm{\Delta })_{,\overline{\delta }}R^{\overline{\eta }}R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\eta }}R_\beta `$
$`+`$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_{\overline{\delta },\beta }^B(R_{\overline{\sigma }}\mathrm{\Delta })^C`$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\sigma }}\mathrm{\Delta }_{,\overline{\delta }})^C,`$
by using the Killing condition (2.15). Then with (A.5),(A.9),(A.10),(4.30) and (4.31) it becomes
$`R_{\alpha \overline{\sigma }\beta \overline{\delta }}`$ $`=`$ $`[f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\beta \mathrm{\Delta })^ER_{\overline{\eta }}^F`$ (B.5)
$`+`$ $`f^{ABC}(R_\beta \mathrm{\Delta })^A(R_{\overline{\delta }}\mathrm{\Delta })^BR^{C\overline{\eta }}f^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^D(R_\alpha \mathrm{\Delta })^ER_{\overline{\eta }}^F]`$
$`+`$ $`[f^{ABC}(R_\alpha \mathrm{\Delta })^A(R_{\overline{\sigma }}\mathrm{\Delta })^Bf^{CDE}(R_\beta ^D\mathrm{\Delta })R_{\overline{\delta }}^E`$
$``$ $`f^{ABC}(R_\alpha \mathrm{\Delta })^AR_\beta ^B(R_{\overline{\eta }}\mathrm{\Delta })^Cf^{DEF}(R_{\overline{\sigma }}\mathrm{\Delta })^DR_{\overline{\delta }}^E(R^{\overline{\eta }}\mathrm{\Delta })^F].`$
The first bracket is symmetric under the interchange of $`\alpha `$ and $`\beta `$. The symmetry of the second bracket can be shown similarly to the previous demonstraration of $`R_{\alpha \overline{\sigma }\beta \overline{\delta }}=R_{\alpha \overline{\delta }\beta \overline{\sigma }}`$.
## Appendix C
We show how to evaluate the Riemann curvature $`G^{(\genfrac{}{}{0pt}{}{ab}{j})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{l})(\genfrac{}{}{0pt}{}{k}{gh})}`$ of the type $`(\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{},\mathrm{𝟏𝟎},\mathrm{𝟏𝟎}^{})`$ by (5.3) and (5.13). For $`i>j`$ we have $`y([E_j^{ab},E_{cd}^i])<0`$. By (5.8) non-trivial components of the Riemann curvature are
$`G^{(\genfrac{}{}{0pt}{}{ab}{j})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{i})(\genfrac{}{}{0pt}{}{j}{gh})}`$ $`=`$ $`tr([E_j^{ab},E_{cd}^i][E_i^{ef},E_{gh}^j])(y(E_i^{ef})y(T_i^j))`$
$``$ $`tr([E_j^{ab},E_i^{ef}][E_{cd}^i,E_{gh}^j]){\displaystyle \frac{y(E_i^{ef})y(E_j^{gh})}{y([E_j^{ab},E_i^{ef}])}}`$
$`=`$ $`\delta _{cd}^{ab}\delta _{gh}^{ef}y(E_j^{ef})\delta _{cdgh}^{abef}{\displaystyle \frac{y(E_i^{ef})y(E_j^{gh})}{y([E_j^{ab},E_i^{ef}])}}.`$
For $`i<j`$ we have $`y([E_j^{ab},E_{cd}^i])>0`$. By (5.7) non-trivial components of the Riemann curvature are
$`G^{(\genfrac{}{}{0pt}{}{ab}{j})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{i})(\genfrac{}{}{0pt}{}{j}{gh})}`$ $`=`$ $`tr([E_j^{ab},E_{cd}^i][E_i^{ef},E_{gh}^j])y(E_i^{ef})`$
$``$ $`tr([E_j^{ab},E_i^{ef}][E_{cd}^i,E_{gh}^j]){\displaystyle \frac{y(E_i^{ef})y(E_j^{gh})}{y([E_j^{ab},E_i^{ef}])}}`$
$`=`$ $`\delta _{cd}^{ab}\delta _{gh}^{ef}y(E_i^{ef})\delta _{cdgh}^{abef}{\displaystyle \frac{y(E_i^{ef})y(E_j^{gh})}{y([E_j^{ab},E_i^{ef}])}}.`$
The same result is also obtained by interchanging the indices as
$$G^{(\genfrac{}{}{0pt}{}{ab}{j})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{i})(\genfrac{}{}{0pt}{}{j}{gh})}=G^{(\genfrac{}{}{0pt}{}{ef}{i})(\genfrac{}{}{0pt}{}{j}{gh})\genfrac{}{}{0pt}{}{ab}{j})(\genfrac{}{}{0pt}{}{i}{cd})}$$
and calculating it by (5.8). For $`i=j`$ we have $`y([E_j^{ab},E_{cd}^i])=0`$. Non-trivial components of the Riemann curvature $`G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{gh})}`$ are evaluated by (5.9). For $`k<i`$
$`G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{gh})}`$ $`=`$ $`tr([E_i^{ab},E_{cd}^i][E_k^{ef},E_{gh}^k])y(E_k^{ef})`$
$``$ $`tr([E_i^{ab},E_k^{ef}][E_{cd}^i,E_{gh}^k]){\displaystyle \frac{y(E_k^{ef})y(E_k^{gh})}{y([E_i^{ab},E_k^{ef}])}}.`$
By the formula
$$tr([E_i^{ab},E_{cd}^i][E_k^{ef},E_{gh}^k])=\delta _{cdgh}^{abef}+\delta _{gh}^{ab}\delta _{cd}^{ef},$$
it becomes
$`G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{ef})}`$ $`=`$ $`\delta _{gh}^{ab}\delta _{cd}^{ef}y(E_k^{ef})\delta _{cdgh}^{abef}{\displaystyle \frac{y(E_i^{ab})y(E_k^{ef})}{y([E_i^{ab},E_k^{ef}])}}.`$
For $`k>i`$
$`G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{gh})}`$ $`=`$ $`\delta _{gh}^{ab}\delta _{cd}^{ef}y(E_i^{ef})\delta _{cdgh}^{abef}{\displaystyle \frac{y(E_i^{ab})y(E_k^{ef})}{y([E_i^{ab},E_k^{ef}])}},`$
by applying the symmetry property
$$G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{gh})}=G^{(\genfrac{}{}{0pt}{}{ef}{k})(\genfrac{}{}{0pt}{}{k}{gh})(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})}$$
to the above result. Of course this can be obtained by a direct calculation. Finally for $`k=i`$
$`G^{(\genfrac{}{}{0pt}{}{ab}{i})(\genfrac{}{}{0pt}{}{i}{cd})(\genfrac{}{}{0pt}{}{ef}{i})(\genfrac{}{}{0pt}{}{i}{gh})}`$ $`=`$ $`tr([E_i^{ab},E_{cd}^i][E_i^{ef},E_{gh}^i])y(E_i^{ef})`$
$`=`$ $`(\delta _{ch}^{ab}\delta _{gd}^{ef}+\delta _{cg}^{ab}\delta _{dh}^{ef}+\delta _{dh}^{ab}\delta _{cg}^{ef}+\delta _{gd}^{ab}\delta _{ch}^{ef})y(E_i^{ef}).`$
Other types of the Riemann curvature are obtained similarly. |
warning/0001/gr-qc0001012.html | ar5iv | text | # Gravitational Solitons and Monodromy Transform Approach to Solution of Integrable Reductions of Einstein Equations
## 1 Introduction
The existence of very rich, integrable structure of Einstein equations, at least for space-times with two-dimensional Abelian isometry groups, have been conjectured by different authors many years ago, but the real discoveries of beautiful integrability properties of these equations and effective procedures for construction of their solutions have been started more than twenty years ago in the papers of Belinskii and Zakharov . In these papers the inverse scattering methods have been developed for Einstein equations for vacuum gravitational fields. In particular, the solution of the entire problem had been reduced to some matrix Riemann - Hilbert problem and soliton generating (dressing) technique had been suggested for calculation of vacuum gravitational solitons on an arbitrary chosen (vacuum) background. Later numerous investigations of integrable reductions of Einstein equations (for vacuum or in the presence of electromagnetic and some other matter fields) have been made by many authors, using different powerful ideas of the modern theory of completely integrable systems. A number of different more or less general approaches were developed and many other interesting results were obtained.<sup>1</sup><sup>1</sup>1Avoiding a detail citation, we refer the readers to the references in a few papers cited here, but mainly – to a large and very useful F.J.Ernst’s collection of related references and abstracts, accessible throw http://pages.slic.com/gravity.
Here we consider Belinskii and Zakharov vacuum soliton generating transformations in the context of so called ”monodromy transform” approach. This approach, developed by the author in \- , provides some general base for the analysis of all known integrable reductions of Einstein equations. In this approach any local solution of reduced Einstein equations is characterized by a set of monodromy data of the fundamental solution of some associated spectral problem.<sup>2</sup><sup>2</sup>2Though the similar constructions surely can be realized for any of the gauge equivalent null curvature representations of the field equations, there were used in this approach a spectral plane $`w`$, which is covered twice by the Belinskii and Zakharov spectral plane $`\lambda `$, and different spectral problem (first constructed by Kinnersley and Chitre in the context of their group - theoretic approach), which is gauge equivalent to the Belinskii and Zakharov one. The reason for the first difference is the absence of differentiation in the transformed linear system with respect to a spectral parameter $`w`$, and for the second one - the simplest monodromy properties of solutions on this spectral plane. The direct and inverse problems of such monodromy transform possess unambiguous solutions.<sup>3</sup><sup>3</sup>3The linear singular integral equation with a scalar kernel, which solves the inverse problem of this monodromy transform, as well as its equivalent regularization – a (quasi-) Fredholm equation of the second kind have been presented in and respectively. It is remarkable, that many physical and geometrical properties of the solutions can be expressed directly in terms of the analytical structure of these monodromy data on the spectral plane. A lot of the known physically interesting solutions possess very simple, rational structures of these monodromy data functions. As it will be shown below, it is also convenient to characterize the Belinskii and Zakharov vacuum soliton generating procedure in terms of the corresponding transformation of the monodromy data of the background solution into the monodromy data of the (multi-) soliton solution. These transformations possess an explicit and remarkably simple, linear-fractional form with coefficients, polynomial in the spectral parameter of the orders less or equal to a number of solitons. This allows to calculate various physical parameters of the generating soliton solutions without (or before) their complete calculation. A generalization of these transformations for the case of electrovacuum solitons, found in , is also presented.
## 2 Belinskii and Zakharov gravitational solitons
The dynamical part of vacuum Einstein equations for the space-time metrics
$$ds^2=f(x^1,x^2)\eta _{\mu \nu }dx^\mu dx^\nu +g_{ab}(x^1,x^2)dx^adx^b$$
(1)
where $`\mu ,\nu =1,2`$; $`a,b=3,4`$; $`f>0`$ and $`\eta _{\mu \nu }=\text{diag }\{ϵ_1,ϵ_2\}`$ with $`ϵ_1=\pm 1`$, $`ϵ_2=\pm 1`$, can be reduced to the form (used in in a bit different notation):
$$\{\begin{array}{c}\eta ^{\mu \nu }_\mu (\alpha _\nu 𝐠𝐠^1)=0\hfill \\ 𝐠^T=𝐠,det𝐠=ϵ\alpha ^2\hfill \end{array}\begin{array}{c}ϵϵ_1ϵ_2\hfill \\ _\mu _\nu \alpha =0\hfill \end{array}\begin{array}{c}\beta :\hfill \end{array}\begin{array}{c}_1\beta =ϵ_1_2\alpha ,\hfill \\ _2\beta =ϵ_2_1\alpha \hfill \end{array}$$
(2)
where $`𝐠=g_{ab}`$ is a real symmetric $`2\times 2`$ \- matrix, and the given above definition of $`ϵ`$ is implied by the Lorentz signature of the metric (1). Therefore, $`ϵ=1`$ corresponds to a hyperbolic case and $`ϵ=1`$ is the elliptic case. Note, that the linear ”harmonic” equation for $`\alpha `$, which follows immediately from the trace of the equation for $`𝐠`$ in (2), provides the existence of the defined above function $`\beta `$, ”harmonically” conjugated to $`\alpha `$. It is convenient to use farther these geometrically defined functions as new coordinates in the linear combinations $`\xi =\beta +j\alpha `$, $`\eta =\beta j\alpha `$, where $`j=1`$ for $`ϵ=1`$ and $`j=i`$ for $`ϵ=1`$.
In these notation the Belinskii and Zakharov spectral problem reads as
$$\{\begin{array}{c}D_\xi 𝚿_{BZ}=\frac{𝕍_\xi }{\lambda j\alpha }𝚿_{BZ}\hfill \\ D_\eta 𝚿_{BZ}=\frac{𝕍_\eta }{\lambda +j\alpha }𝚿_{BZ}\hfill \end{array}\left|\begin{array}{c}D_\xi =_\xi \frac{\lambda }{\lambda j\alpha }\frac{}{\lambda }\hfill \\ D_\eta =_\eta \frac{\lambda }{\lambda +j\alpha }\frac{}{\lambda }\hfill \end{array}\right|\begin{array}{c}𝕍_\xi =j\alpha _\xi 𝐠𝐠^1\hfill \\ 𝕍_\eta =j\alpha _\eta 𝐠𝐠^1\hfill \end{array}$$
(3)
with additional ”reduction” conditions imposed on the solutions of (3):
$$𝚿_{BZ}^T(\lambda \frac{ϵ\alpha ^2}{\lambda })𝐠^1𝚿_{BZ}=𝐊_{BZ}(w),𝐊_{BZ}^T=𝐊_{BZ},D_\xi w=D_\eta w=0$$
(4)
For the construction of vacuum solitons, based on the form (3), (4) of the spectral problem, the dressing transformation $`𝚿_{BZ}(\xi ,\eta ,\lambda )=\chi _{BZ}(\xi ,\eta ,\lambda )\underset{BZ}{\overset{o}{𝚿}}(\xi ,\eta ,\lambda )`$ was used in with the soliton ansatz<sup>4</sup><sup>4</sup>4In the case $`ϵ=1`$ the number of solitons (poles) can be odd as well, but we shall not consider this case here.
$$\chi _{BZ}=𝕀+\underset{k=1}{\overset{2N}{}}\frac{_k(\xi ,\eta )}{\lambda \mu _k(\xi ,\eta )},\chi _{BZ}^1=𝕀+\underset{k=1}{\overset{2N}{}}\frac{𝕊_k(\xi ,\eta )}{\lambda \nu _k(\xi ,\eta )}$$
(5)
where $`𝕀`$ is the identity matrix, the $`2\times 2`$-matrix residues at the poles as well as the pole trajectories $`\mu _k`$ and $`\nu _k`$ are real or constitute complex conjugated pairs. Then in all constraints on these functions, which follow from (3), had been successfully solved and all metric components of the generating soliton solutions have been explicitly expressed in terms of a number of integration constants and the matrix $`\underset{BZ}{\overset{o}{𝚿}}(\xi ,\eta ,\lambda )`$, characterizing arbitrarily chosen ”background” solution. (See also for more compact, determinant form of these soliton solutions).
## 3 Vacuum $`w`$ \- solitons.
The function $`w`$, which was used as the new spectral parameter in the mentioned above monodromy transform approach, have been introduced also in , as a solution of the equations $`D_\xi w=D_\eta w=0`$, such that $`2w=\lambda +2\beta +ϵ\alpha ^2/\lambda .`$ Thus, the spectral plane $`\lambda `$ covers twice the spectral plane $`w`$. The Kinnersley - Chitre-like linear system which substitute in our approach the linear system (3), and the gauge transformation between these systems in our notation are
$$\{\begin{array}{c}_\xi 𝚿=\frac{𝐔(\xi ,\eta )}{2i(w\xi )}𝚿\hfill \\ _\eta 𝚿=\frac{𝐕(\xi ,\eta )}{2i(w\eta )}𝚿\hfill \end{array}\begin{array}{c}𝚿_{BZ}(\xi ,\eta ,\lambda )=𝐀_{}𝚿(\xi ,\eta ,w)\hfill \\ 𝐀_{}=𝐠\epsilon +i\lambda 𝕀,\epsilon =\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)\hfill \end{array}$$
(6)
The following two groups of conditions constitute a spectral problem, based on the linear system (6) and equivalent to the reduced vacuum equations (2) :
$`\{\begin{array}{c}2i(w\xi )_\xi 𝚿=𝐔(\xi ,\eta )𝚿\hfill \\ 2i(w\eta )_\eta 𝚿=𝐕(\xi ,\eta )𝚿\hfill \end{array}\begin{array}{c}\text{rank }𝐔=1,\text{tr }𝐔=i,\hfill \\ \text{rank }𝐕=1,\text{tr }𝐕=i,\hfill \end{array}`$ (11)
$`\{\begin{array}{c}𝚿^{}𝐖𝚿=𝐊(w)\hfill \\ 𝐊^{}(w)=𝐊(w)\hfill \end{array}{\displaystyle \frac{𝐖}{w}}=4i\epsilon ,\epsilon =\left(\begin{array}{cc}\hfill 0& \hfill 1\\ \hfill 1& \hfill 0\end{array}\right)`$ (16)
For the construction of the Belinskii and Zakharov vacuum solitons, based on the spectral problem (11), (16), we use similar dressing transformation $`𝚿(\xi ,\eta ,w)=\chi (\xi ,\eta ,w)\stackrel{o}{𝚿}(\xi ,\eta ,w)`$ with a bit different, than in (5), soliton ansatz <sup>5</sup><sup>5</sup>5For the case of real poles an alternative construction of the Belinskii and Zakharov solitons in the context of the spectral problem (11), (16), which corresponds to the choice $`X(w)=1`$ in (17) and $`𝐊(w)=_{k=1}^N\left(\frac{w\stackrel{~}{w}_k}{ww_k}\right)\stackrel{o}{𝐊}(w)`$, had been presented in .
$$\chi =X(w)\left(𝐈+\underset{k=1}{\overset{N}{}}\frac{𝐑_k(\xi ,\eta )}{ww_k}\right),\chi ^1=\frac{1}{X(w)}\left(𝐈+\underset{k=1}{\overset{N}{}}\frac{𝐒_k(\xi ,\eta )}{w\stackrel{~}{w}_k}\right),$$
(17)
where $`X(w)=_{k=1}^N\left(\frac{ww_k}{w\stackrel{~}{w}_k}\right)^{\frac{1}{2}}`$, the gauge condition $`𝐊(w)=\stackrel{o}{𝐊}(w)`$ is used, the constant pole locations $`w_k`$ and $`\stackrel{~}{w}_k`$ for each $`k=1,2,\mathrm{},N`$ are the pairs of different real or complex conjugated constants, <sup>6</sup><sup>6</sup>6It is useful to note, that the number $`N`$ of $`w`$-poles (solitons) in (17) corresponds to $`2N`$ solitons ($`\lambda `$-poles) in (5); each real pole with $`w_k\stackrel{~}{w}_k`$ is equivalent to a pair of real $`\lambda `$ \- poles in (5), while each of the complex $`w`$-poles with $`\stackrel{~}{w}_k=\overline{w_k}`$ correspond to a pair of complex conjugated to each other $`\lambda `$-poles in (5). and the values, having ”o” overhead, correspond to the background solution. Again, similarly to the original Belinskii and Zakharov construction of $`\lambda `$-solitons, all constraints, which follow from (11), (16) for the matrix residues in (17) can be solved , and for any choice of the background solution the generating $`N`$-soliton solution can be expressed in terms of the background $`\stackrel{o}{𝚿}(\xi ,\eta ,w)`$ and a set of integration constants (four real constants per each of the poles $`w_k`$).
## 4 Electrovacuum $`w`$ \- solitons
One of the important features of the spectral problem (11), (16) is that very small changes of its structure lead to a matrix problem, equivalent to the space-time symmetry reduced electrovacuum Einstein - Maxwell equations :
$`\{\begin{array}{c}2i(w\xi )_\xi 𝚿=𝐔(\xi ,\eta )𝚿\hfill \\ 2i(w\eta )_\eta 𝚿=𝐕(\xi ,\eta )𝚿\hfill \end{array}\begin{array}{c}\text{rank }𝐔=1,\text{tr }𝐔=i,\hfill \\ \text{rank }𝐕=1,\text{tr }𝐕=i,\hfill \end{array}`$ (22)
$`\{\begin{array}{c}𝚿^{}𝐖𝚿=𝐊(w)\hfill \\ 𝐊^{}(w)=𝐊(w)\hfill \end{array}{\displaystyle \frac{𝐖}{w}}=4i\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0\end{array}\right),𝐖^{55}=1`$ (28)
where the unknown matrix variables $`𝚿(\xi ,\eta ,w)`$, $`𝐔(\xi ,\eta )`$, $`𝐕(\xi ,\eta )`$ and $`𝐖(\xi ,\eta ,w)`$ are now $`3\times 3`$ \- matrices and $`𝐖^{55}`$ denotes the lower right component of $`3\times 3`$-matrix $`𝐖`$ and $`𝚿^{}(\xi ,\eta ,w)\overline{𝚿^T(\xi ,\eta ,\overline{w})}`$. With the use of (22), (28) the electrovacuum soliton solutions have been constructed using the same dressing transformation $`𝚿=\chi \stackrel{o}{𝚿}`$ with a more particular, than (17), soliton ansatz :
$$\chi =𝐈+\underset{k=1}{\overset{N}{}}\frac{𝐑_k(\xi ,\eta )}{ww_k},\chi ^1=𝐈+\underset{k=1}{\overset{N}{}}\frac{𝐒_k(\xi ,\eta )}{w\stackrel{~}{w}_k}$$
(29)
where the poles of $`\chi `$ are complex conjugated to the poles of $`\chi ^1`$, i.e. for each $`k=1,2,\mathrm{},N`$ we have $`\stackrel{~}{w}_k=\overline{w_k}`$. This leads to electrovacuum generalization (which includes six real parameters per each $`w`$-pole of $`\chi `$) of the Belinskii and Zakharov vacuum solitons with complex conjugated poles, while a generalization of vacuum solitons with real poles does not arise in this way. <sup>7</sup><sup>7</sup>7Many electrovacuum solutions which generalize vacuum solitons with real poles and on some specially chosen backgrounds (e.g., on the Minkowski background) can be constructed as the analytical continuations of soliton solutions with complex poles in the space of their constant parameters. This complex analytical continuation is quite similar to the known one, which relates the ”underextreme” and ”overextreme” parts of the Kerr - Newman family of solutions. Another way for construction of such solutions is a direct integration of the integral equation, which the spectral problem (22), (28) can be reduced to. This integral equation, derived in , can be solved and the solutions with arbitrarily large number of free constant parameters can be constructed explicitly for rational values of some special functional parameters in the kernel of this equation, called as ”monodromy data”, and expressed in terms of a number of arbitrary rational functions of spectral parameter $`w`$ .
## 5 Monodromy data of the solutions
The analysis, suggested below, is based on the monodromy transform approach, developed in \- . In this approach any local solutions of reduced vacuum or electrovacuum Einstein equations is characterized unambiguously by a finite set of functional parameters, which are functions of the spectral parameter $`w`$ only and which admit a simple interpretation as the monodromy data on the spectral plane for the fundamental solution $`𝚿(\xi ,\eta ,w)`$ of the spectral problem (22), (28), corresponding to a given solution of the field equations under consideration. In order to define the monodromy data, we need, first of all, to fix some gauge freedom, existing in (22), (28). For this we impose at some chosen ”initial” space-time point $`(\xi _0,\eta _0)`$<sup>8</sup><sup>8</sup>8Actually, this is a point in the orbit space of the space-time isometry group. the universal ”normalization” conditions for the metric components (say, $`𝐠(\xi _0,\eta _0)=ϵ_0\text{diag }\{1,ϵ\alpha _0^2\}`$, where $`ϵ_0=\pm 1`$), which determine the value of $`𝐖_o(w)𝐖(\xi _0,\eta _0,w)`$ as
$$𝐖_o(w)=4i(w\beta _0)\left(\begin{array}{ccc}\hfill 0& \hfill 1& \hfill 0\\ \hfill 1& \hfill 0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 0\end{array}\right)+\left(\begin{array}{ccc}\hfill 4ϵ_0ϵ\alpha _0^2& \hfill 0& \hfill 0\\ \hfill 0& \hfill 4ϵ_0& \hfill 0\\ \hfill 0& \hfill 0& \hfill 1\end{array}\right)$$
(30)
where $`\alpha _0=(\xi _0\eta _0)/2j`$, $`\beta _0=(\xi _0+\eta _0)/2`$. The normalization condition, used for the value of $`𝚿(\xi _0,\eta _0,w)`$, which determines then the ”normalized” value of $`𝐊(w)`$, and the corresponding gauge transformations can be chosen in the form
$$\begin{array}{c}𝚿(\xi _0,\eta _0,w)=𝐈\hfill \\ 𝐊(w)=𝐖_o(w)\hfill \end{array}\begin{array}{c}𝚿𝚿𝐂(w),𝐂(w)=𝚿^1(\xi _0,\eta _0,w)\hfill \\ \begin{array}{c}𝚿𝐀𝚿𝐀^1,\hfill \\ 𝐖(𝐀^{})^1𝐖𝐀^1,\hfill \end{array}𝐀(w)=\left(\begin{array}{cc}SL(2,R)\hfill & \begin{array}{c}0\\ 0\end{array}\hfill \\ \begin{array}{cc}a^3& a^4\end{array}\hfill & 1\hfill \end{array}\right)\hfill \end{array}$$
(31)
where the $`3\times 3`$ \- matrix $`𝐀`$ is constant. The complex constants $`a^3`$ and $`a^4`$ change the additive constants in the definitions of the components of a complex electromagnetic potentials and $`SL(2,R)`$ part of $`𝐀`$ corresponds to a linear transformation of the coordinates $`x^3`$, $`x^4`$ in (1).
A detail analysis of the analytical structure on the spectral plane of the solutions of (22), (28) shows the existence of some universal properties of $`𝚿(\xi ,\eta ,w)`$.<sup>9</sup><sup>9</sup>9From now, the functions $`𝚿`$, $`𝐖`$, $`𝐊`$ are considered in the gauges, fixed as in (30), (31). In particular, it is holomorphic function of $`w`$ everywhere outside four algebraic branchpoints and the cut $`L=L_++L_{}`$ joining these points, as it is shown on Fig. 1. It turns out, that the behaviour of $`𝚿`$ near the branchpoints can be described by the monodromy matrices $`𝐓_\pm (w)`$, which characterize the linear transformations of $`𝚿`$, continued analytically along the paths $`T_\pm `$, rounding one of the branchpoints and joining different edges of $`L_+`$ or $`L_{}`$ respectively:
$$𝚿\stackrel{T_\pm }{}\stackrel{~}{𝚿}=𝚿𝐓_\pm (w),𝐓_\pm (w)=𝐈2\frac{𝐥_\pm (w)𝐤_\pm (w)}{(𝐥_\pm (w)𝐤_\pm (w))}.$$
(32)
It is remarkable, that these matrices, satisfying the identities $`𝐓_\pm ^2(w)𝐈`$, are independent of the space-time coordinates $`\xi `$, $`\eta `$. The structure (32) allows to express $`𝐓_\pm `$ in terms of the four complex projective vectors $`𝐤_\pm (w)`$ and $`𝐥_\pm (w)`$, but it was found in , that (28) relate unambiguously $`𝐥_\pm (w)`$ and $`𝐤_\pm ^{}(w)`$ with the same suffices. Therefore, both $`𝐓_\pm `$ are determined completely by four scalar functions, which parametrize the components of two projective vectors $`𝐤_\pm (w)`$:
$$𝐤_\pm (w)=\{1,𝐮_\pm (w),𝐯_\pm (w)\}$$
(33)
The functions $`𝐮_\pm (w)`$, $`𝐯_\pm (w)`$, which domains of holomorphicity are shown on Fig. 1, are called as monodromy data. These data characterize unambiguously any local solution of the field equations near some chosen ”initial” or ”reference” point $`(\xi _0,\eta _0)`$, however, it is useful to keep in mind, that for a given solution these data are dependent upon the choice of this point. Note, that for pure vacuum case $`𝐯(w)0`$.
## 6 Soliton generating transformations in terms of the monodromy data
Besides the definition of the monodromy data for any local vacuum or electrovacuum solution, described in the previous section, we recall here that these data enter explicitly into the local structure of $`𝚿`$ near the cuts $`L_\pm `$. This structure can be described by the expressions :
$$𝚿=\lambda _\pm ^1\psi _\pm (\xi ,\eta ,w)𝐤_\pm (w)+𝐌_\pm (\xi ,\eta ,w)$$
(34)
where $`\lambda _+=\sqrt{(w\xi )/(w\xi _0)}`$, $`\lambda _{}=\sqrt{(w\eta )/(w\eta _0)}`$ and the column-vectors $`\psi _\pm `$, the row-vectors $`𝐤_\pm `$ and the matrices $`𝐌_\pm `$ are holomorphic at the points of the cut $`L_+`$ or $`L_{}`$, respectively to their suffices. In accordance with these expressions, the monodromy data can be calculated from the branching parts of the components of $`𝚿`$ on the cuts $`L_\pm `$. We use now the expressions (34) for calculation of the transformations of the vectors $`𝐤_\pm `$ and therefore, of the monodromy data, induced by the soliton generating transformations.
It is easy to see, that any dressing transformation $`𝚿=\chi \stackrel{o}{𝚿}(\xi ,\eta ,w)`$ for the normalized $`𝚿`$ \- functions reads
$$𝚿=𝐀\chi \stackrel{o}{𝚿}\chi _0^1𝐀^1$$
(35)
where $`\stackrel{o}{𝚿}`$ is the normalized background solution, $`\chi _0(w)=\chi (\xi _0,\eta _0,w)`$ and the constant matrix $`𝐀`$ should be chosen afterwards for normalization of the metric components and the matrix $`𝐖`$. Let us consider now the expression (35) near the cuts $`L_\pm `$, using there the local representations (34) for $`𝚿`$ and $`\stackrel{o}{𝚿}`$. These expressions imply the following transformations of the projective vectors $`𝐤(w)`$:
$$𝐤_\pm (w)=\stackrel{o}{𝐤}_\pm (w)\chi _0^1(w)𝐀^1.$$
(36)
Now it is easy to calculate the monodromy data $`𝐮_\pm (w)`$, $`𝐯_\pm (w)`$ for solitons, as the ratios of the components of the projective vectors (36), using one of the soliton ansatz (17) or (29). The result can be presented in the linear-fractional form (for vacuum background $`\stackrel{o}{𝐯}_\pm (w)0`$ and for vacuum solitons $`𝐯_\pm (w)0`$):
$$𝐮_\pm (w)=\frac{𝒰_0+𝒰_1\stackrel{o}{𝐮}_\pm (w)+𝒰_2\stackrel{o}{𝐯}_\pm (w)}{𝒬_0+𝒬_1\stackrel{o}{𝐮}_\pm (w)+𝒬_2\stackrel{o}{𝐯}_\pm (w)},𝐯_\pm (w)=\frac{𝒱_0+𝒱_1\stackrel{o}{𝐮}_\pm (w)+𝒱_2\stackrel{o}{𝐯}_\pm (w)}{𝒬_0+𝒬_1\stackrel{o}{𝐮}_\pm (w)+𝒬_2\stackrel{o}{𝐯}_\pm (w)},$$
where all coefficients are polynomials in $`w`$, which orders do not exceed the number of w-solitons $`N`$ (or $`2N`$ in the Belinskii and Zakharov formalism).
## Acknowledgments
This work was partly supported by the British Engineering and Physical Sciences Research Council and the Russian Foundation for Basic Research Grants 99-01-01150, 99-02-18415. |
warning/0001/hep-ex0001044.html | ar5iv | text | # Low energy atmospheric muon neutrinos in MACRO
## Abstract
We present the measurement of two event samples induced by atmospheric $`\nu _\mu `$ of average energy $`\overline{E}_\nu 4GeV`$. In the first sample, the neutrino interacts inside the MACRO detector producing an upward-going muon leaving the apparatus. The ratio of the number of observed to expected events is $`0.57\pm 0.05_{stat}\pm 0.06_{syst}\pm 0.14_{theor}`$ with an angular distribution similar to that expected from the Bartol atmospheric neutrino flux. The second is a mixed sample of internally produced downward-going muons and externally produced upward-going muons stopping inside the detector. These two subsamples are selected by topological criteria; the lack of timing information makes it impossible to distinguish stopping from downgoing muons. The ratio of the number of observed to expected events is $`0.71\pm 0.05_{stat}\pm 0.07_{syst}\pm 0.18_{theor}`$ . Using the ratio of the two subsamples (for which most theoretical uncertainties cancel) we can test the pathlength dependence of the oscillation hypothesis. The probability of agreement with the no-oscillation hypothesis is $`5\%`$ .
The deviations of our observations from the expectations has a preferred interpretation in terms of $`\nu _\mu `$ oscillations with maximal mixing and $`\mathrm{\Delta }m^210^3÷10^2eV^2`$. These parameters are in agreement with our results from upward throughgoing muons, induced by $`\nu _\mu `$ of much higher energies.
The results from several underground detectors which measure the flux of atmospheric neutrinos give strong indication that $`\nu _\mu `$’s oscillate into neutrinos of another type . Fully-contained and partially-contained neutrino-induced events observed in underground detectors come from neutrinos of energy $`1GeV`$. The flux of atmospheric neutrinos of several tens of $`GeV`$ can be inferred from the measurement of neutrino-induced upward-going muons that traverse the entire detector (up-throughgoing muons). The hypothesis of neutrino oscillations, with best-fit parameters $`\mathrm{sin}^22\theta _{mix}1`$ and $`\mathrm{\Delta }m^2`$ in the range of a few times 10<sup>-3</sup> eV<sup>2</sup>, can explain the observed anomalies both in the ratio of contained $`\nu _\mu `$ to $`\nu _e`$ events (Super-K, Soudan 2) and in the zenith angle distribution of up-throughgoing muons (MACRO, Super-K).
The MACRO detector measures both the high energy (median energy $`50GeV`$) and few $`GeV`$ energy atmospheric neutrino fluxes. In Ref. the interpretation of the data in terms of $`\nu _\mu `$ oscillations came from a deficit and from an anomalous zenith angle distribution of the observed up-throughgoing muons originating from $`\nu _\mu `$ interactions in the rock below the detector. Here we report on the measurement of the flux of lower energy $`(\overline{E}_\nu 4GeV`$) atmospheric neutrinos through the detection of $`\nu _\mu `$ interactions inside the apparatus (yielding partially contained upgoing and downgoing muons) and by the detection of externally produced upward-going muons stopping inside the detector .
MACRO is a large area, modular tracking detector located in Hall B of the Gran Sasso Underground Laboratory in Italy, with an average rock overburden of 3700 hg/cm<sup>2</sup>. It is a rectangular box, 76.6 m $`\times `$ 12 m $`\times `$ 9.3 m, divided longitudinally into six supermodules and vertically into a lower part (4.8 m high) and an upper part (4.5 m high). The active detection elements are planes of streamer tubes for tracking and liquid scintillation counters for fast timing. The lower half of the detector is filled with streamer tube planes alternating with trays of crushed rock, which provide most of the 5.3 $`kton`$ target mass for partially-contained neutrino interactions. The upper part is hollow and contains the electronics racks and work areas. There are 10 horizontal streamer tube planes in the bottom half of the detector, and 4 planes on the top, all with wire and 27 stereo strip readouts. Six vertical planes of streamer tubes cover each side of the detector. The intrinsic angular resolution for muons is between 0.1 and 1.0 depending on the track length. The scintillator system consists of three widely-separated layers of horizontal boxes, and on each vertical side of the detector a layer of vertical boxes inserted between the streamer tubes. The time (position) resolution for muons in a scintillator box is about 500 ps ($`11`$ cm). The direction of the muons passing across MACRO is determined by the time-of-flight between two layers of scintillation counters.
The results presented in this letter come from 4.1 live years of data taking with the full detector, from April 1994 to February 1999.
About $`33\times 10^6`$ downgoing muons were collected, and were used to monitor the detector efficiency, the running conditions and the acceptance. The trigger rate due to downgoing muons is $`0.3Hz`$. The trigger efficiency for each scintillation counter and for the streamer tubes was monitored over the data taking period using the downgoing muons.
Two samples of atmospheric muon neutrinos in the few-$`GeV`$ energy range are measured. In the first sample (up partially-contained or $`IU`$=Internal Upgoing $`\mu `$ events) there are (mainly) events induced by charged current (CC) interactions of upgoing $`\nu _\mu `$ inside the lower part of MACRO. An upgoing muon is produced, which crosses two scintillation layers (Fig. 1), so that the measurement of the direction is made through time-of-flight.
The second sample is a mix of upgoing and downgoing events. The partially contained downgoing events (down partially-contained or $`ID`$=Internal Downgoing $`\mu `$) are induced by downgoing $`\nu _\mu `$, interacting in the lower part of MACRO. The upgoing stopping events ( $`UGS`$ = Upward Going Stopping muons) are induced by interactions of upgoing $`\nu _\mu `$ below the detector yielding an upgoing muon which stops inside the detector. Both the down partially-contained and the upgoing stop events cross only the bottom layer of liquid scintillation counters (see Fig. 1) and are identified by means of topological criteria. The lack of timing information makes it impossible to distinguish between the two subsamples. Fig. 2 shows the parent neutrino energy distribution from a Monte Carlo calculation for the three event topologies detectable in MACRO. The energy spectrum and the median energy of the two samples presented in this letter are almost the same.
The identification of $`IU`$ events is based on topological criteria and time-of-flight measurements. The main requirement is the presence of at least two hit scintillator clusters, respectively in the center layer and in the upper part of the apparatus (see Fig. 1). This is the expected topology for a neutrino interacting in the lower detector and producing an upward-going muon with enough energy to exit the apparatus. It is also the topology of the much more numerous downgoing muons stopping in the lower detector. Scintillation timing allows the separation of the two classes of events. Moreover, the scintillators are required to match a streamer tube track reconstructed in space by our standard track-finding algorithms . For $`IU`$ candidates, the lowest point of a track (the starting point) must be inside the apparatus as a condition for the containment of the $`\nu _\mu `$ interaction vertex. To reject fake partially-contained events entering from a detector insensitive zone, the extrapolation of the track in the lower part of the detector must geometrically cross and not fire at least one scintillator layer and one streamer tube plane, or at least three planes of streamer tubes. These conditions were tuned on Monte Carlo simulated events, including evaluation of detector inefficiencies. Other cuts are applied to reject background events from downgoing atmospheric muons. They are related to the goodness of the geometrical agreement between scintillator hits and the streamer track, to the proper operation of the scintillation counters and to the quality of the time measurement. The measured $`1/\beta `$ distribution after all analysis cuts (including the requirements of vertex containment) is shown in Fig. 3. The measured muon velocity $`\beta c`$ is evaluated with the sign convention that upgoing muons have $`1/\beta 1`$. A total of 121 events survive in the range $`1.3<1/\beta <0.7`$, which is taken as the range of $`IU`$ signal.
From the time distribution of Fig. 3 one expects some background events; they are mostly due to wrong time measurements or secondary particle hits, yielding an almost flat $`1/\beta `$ distribution. The fit of the distribution in the range $`6.0<1/\beta <0.3`$ to a gaussian plus a straight line gives an estimated background of 5 events in the signal region. After background subtraction, we have 116 up partially-contained events.
The identification of $`ID+UGS`$ events is based on topological criteria. The candidates have a track starting (ending) in the lower apparatus, and crossing the bottom detector face. The track must also be located or oriented in such a way that it could not have entered or exited undetected through insensitive zones in the apparatus. Events with scintillator hits outside the bottom layer, or with the reconstructed track pointing to a detector insensitive zone between modules, are rejected. The standard muon tracking procedure is based on at least four aligned hits in the streamer tubes. This corresponds to a minimum traversed detector thickness $`t_v200g/cm^2`$ (standard sample). A dedicated retracking procedure was applied to all remaining events. The retracking procedure requires at least three streamer tube hits ($`t_v100g/cm^2`$), aligned with respect to a fired scintillation counter in the bottom layer. The number of retracked events is $`5\times 10^5`$ (to be compared with the $`32\times 10^6`$ standard muon tracks); these events follow the same analysis as the standard sample.
After the software cuts, 879 events survive. Some of them are tracked incorrectly (mostly by the retracking), or are bending downgoing muons, entering from a detector insensitive zone. Due to the bending, only a fraction of the streamer tube hits are used by the tracking algorithms. In order to reject these fake candidates, we made a double scan with the MACRO Event Display. To eliminate any bias from the scan procedure, and to evaluate the absolute and relative scanning efficiencies, Monte Carlo (MC) simulated events (described below) were randomly injected into the data sample before the scanning stage. Two physicists independently scanned the merged sample. At the end of the scan, 200 events in the real data (106 $`ID+UGS`$ candidates are retracked events) are accepted as upgoing stopping or partially contained downgoing muons. 97% of the real events selected by one physicist were also selected by the other.
Downgoing muons which pass near or through the detector may produce low-energy, upgoing particles, which could simulate neutrino-induced upgoing muons if the downgoing muon misses the detector. This background has been evaluated using a full simulation, based on our measurements . The background is $`7\pm 2`$ events. A second background source could arise from atmospheric muons and detector inefficiency. Using a simulated sample of $`10^7`$ atmospheric muons, which includes measured detector inefficiencies, no events were selected by the above described procedure. After background subtraction, 193 events represent the down partially-contained plus upgoing stopping signal.
The expected number of neutrino-induced events was estimated from GMACRO , a GEANT-based full MC detector simulation. The $`\nu _e`$ and $`\nu _\mu `$ interaction rates have been computed using the atmospheric neutrino flux of the Bartol group and the neutrino cross sections of Ref. . In this cross-section model, the contributions of the exclusive channels of lowest multiplicity (quasi-elastic and single pion production) are calculated separately from deep inelastic scattering (DIS). The DIS contribution to the $`\nu `$N cross section was computed using the GRV-LO-94 parton distribution functions. Using these neutrino fluxes and cross sections for $`E_\nu 300MeV`$, we expect a total ($`CC+NC`$) interaction rate of $`71.5/ktony(\nu _e+\overline{\nu }_e)`$ and $`148.1/ktony(\nu _\mu +\overline{\nu }_\mu )`$. Two simulated samples have been generated, because of the different vertex locations for the $`IU`$ and the $`ID+UGS`$ events; the simulated events were processed with the same analysis chain as the data.
For the $`IU`$ events, a sample of $`10^5`$ interactions inside the apparatus was generated (equivalent to $`85.9`$ years live time). The simulation indicates (see Table 1) that $`87\%`$ of detected $`IU`$ events come from $`\nu _\mu `$ charged current (CC) interactions, $`9\%`$ from $`\nu _e`$ CC and the remaining fraction from neutral current (NC) interactions. Due to detector inefficiencies and analysis algorithm failures, some neutrino-induced events originating in the rock surrounding the detector are expected to contribute to the selected sample of up partially-contained events (upward-throughgoing $`\mu `$’s appearing as partially-contained). The vertex containment requirements reduce this background to about $`1\%`$, evaluated using a simulated sample of up throughgoing muons . The fully-automated selection gives a total number of 202 expected events in the $`IU`$ event signal region, $`1.3<1/\beta <0.7`$ for $`4.1y`$ of live-time .
For the $`ID+UGS`$ events, $`1.16\times 10^6`$ neutrino interactions were simulated in a larger volume (including the experimental hall, the detector and the surrounding rock). The generated events correspond to a live time of $`31.1y`$. The dimensions of the interaction volume (with $`13m`$ of rock below the detector, and a total rock mass of $`165kton`$ plus $`5.3kton`$ of the apparatus itself) were chosen to reduce to less than $`1\%`$ the number of $`\nu _\mu `$-induced stopping muons produced outside that volume. The 2199 events which survived the software selection for the $`ID+UGS`$ were merged with the real events which passed the same software selection, and visually scanned. After the scan procedure 2074 (=94.3%) of the simulated and reconstructed events were accepted, together with the 200 real events. The expected rate is 273 events in $`4.1y`$ live-time. In Table I we give the main features of the $`IU`$ and $`ID+UGS`$ simulated samples (expected rate, percentage of CC $`\nu _\mu `$ interactions, median parent neutrino energy, and fraction of events induced by upgoing neutrinos).
The number of detected $`IU`$ events in the 4.1 $`y`$ live-time is 116, while the expected number is 202. The ratio of the measured to the expected events is $`R_{IU}=(\frac{Data}{MC})_{IU}=0.57\pm 0.05_{stat}\pm 0.06_{syst}\pm 0.14_{theor}`$. For the $`ID+UGS`$ sample, 193 events are detected while 273 are expected. The ratio is $`R_{ID+UGS}=0.71\pm 0.05_{stat}\pm 0.07_{syst}\pm 0.18_{theor}`$. Each data set is (within errors) consistent with a constant deficit (43% for the $`IU`$ sample, 29% for the ID+UGS) in all zenith bins compared to the Monte Carlo expectations assuming no oscillations. Fig. 4 shows the zenith angle distributions for the two measured data sets and for the Monte Carlo simulations. Due to the analisys cuts and to the apparatus acceptance, there are no events in the last bin of the two distributions. The expectations are affected by a systematic theoretical error due to the uncertainties regarding the atmospheric neutrino flux and the neutrino cross sections. At present there is no unique and reliable estimate of the total theoretical uncertainty for the rate calculations. Each experimental group, and for each event category, has its own way to estimate it. For this analysis we conservatively estimate 20% for the flux and 15% for the cross section, which add in quadrature to an error of 25%. For our high energy events we quoted 17%, while in the recent SuperKamiokande analysis of neutrino-induced stopping muons 22% was quoted .
Our measured value of $`R_{IU}=0.57`$ is quite far from its expected value of unity. If we ignored theoretical errors (i.e. if we assumed the flux and cross section as we modeled them were accurate), an experiment with our statistical and experimental uncertainties would only fluctuate so far from unity with probability $`2.5\times 10^4`$. However, if we take the $`25\%`$ theoretical error into account, the probability becomes $`6.5\%`$. $`R_{ID+UGS}=0.71`$ also differs from unity, though not as significantly as $`R_{IU}`$.
If the observed deficit were due only to an overall theoretical overestimate of the neutrino flux or cross sections, one would expect $`R_{IU}R_{ID+UGS}`$ (small differences would remain due to residual geomagnetic effects). Furthermore, the theoretical uncertainties largely cancel if the ratio $`IU/ID+UGS`$ between the measured number of events is compared with the expectation. The partial error cancellation arises from the almost identical energy spectra of parent neutrinos for the two samples of events; we evaluated the remaining error as 5%. The experimental systematic uncertainty for the ratio is estimated at 6%. The measured ratio is $`\frac{IU}{ID+UGS}=0.60\pm 0.07_{stat}`$, while the expectation without oscillations is $`0.74\pm 0.04_{sys}\pm 0.04_{theo}`$. The probability to obtain a ratio so far from the expected one is 5%, almost independent of the neutrino flux and neutrino cross sections used for the predictions.
We investigated if the observed discrepancies between data and expectations could be explained by possible systematic effects. The detector mass is known to $`\pm 5\%`$. The uncertainty for the detector acceptance was estimated by comparing the shape of the zenith distribution of downward-going muons stopping inside the detector with a MC expectation based on the known rock overburden: the two distributions agree within $`6\%`$. Other uncertainties arise from the live-time estimate ($`3\%`$), the effective containment of the interaction vertex depending on the simulation of the detector response to internal neutrino interactions ($`4\%`$) and the background subtraction ($`4\%`$). Adding all these contributions in quadrature yields our quoted experimental systematic uncertainty of $`10\%`$, too small to account for the observed discrepancy.
The number of expected events was also evaluated using the NEUGEN neutrino event generator (developed by the Soudan and MINOS collaborations) as input to our MC simulations. The NEUGEN generator predicts $`6\%(5\%)`$ fewer $`IU`$ (ID+UGS) events than our default generator , well within the estimated systematic uncertainty for neutrino cross sections ($`15\%`$).
Our data disfavor the no-oscillations hypothesis regardless of overall normalization; they are consistent with neutrino oscillations ($`\nu _\mu `$ disappearance) with maximal mixing and $`\mathrm{\Delta }m^2(1÷20)\times 10^3eV^2`$. As a “test point”, we use the best-fit parameters from our high-energy analyses , $`\mathrm{\Delta }m^2=2.5\times 10^3eV^2`$ and $`sin^22\theta _{mix}=1`$. The predicted numbers of events and the angular distributions are indicated by the dashed histograms in Fig. 4; they are in good agreement with the measured data.
For $`\mathrm{\Delta }m^2(1÷20)\times 10^3eV^2`$, upgoing neutrinos (which induce $`IU`$ and $`UGS`$ events), which travel thousands of kilometers through the Earth, are reduced by $`50\%`$. Almost no reduction is expected for downgoing partially-contained muons. In this scenario, and for a pure $`\nu _\mu `$ CC interaction sample, the expected event rate is 1/2 of the $`IU`$ and 3/4 of the $`ID+UGS`$ predictions without oscillations. The predicted reduction for upgoing $`\nu _\mu `$ is less than 1/2 because of the $`\nu _e`$ and NC event contaminations. Our data disfavor $`\mathrm{\Delta }m^2>10^2eV^2`$, for which the $`ID`$ events are also reduced; both the $`ID+UGS`$ and $`IU`$ event rates are $`1/2`$ of the no-oscillations expectation. We also disfavor $`\mathrm{\Delta }m^2<10^3eV^2`$, for which the shape of the angular distributions (Fig. 4) is modified.
Assuming oscillations (with the “test point” parameters) 115 up partially-contained and 202 down partially-contained plus upgoing stopping events are expected. For the $`IU`$ events, the reduction from the no oscillations hypothesis is 0.57, to be compared with the measured value of $`R_{IU}=0.57`$. For the $`ID+UGS`$ events, it is 0.76, to be compared with $`R_{ID+UGS}=0.71`$. The quoted numbers use our default normalization. Recent flux calculations suggest that the Bartol flux which we use may be too high (though within the quoted theoretical error). As far as the event rates are concerned, a lower normalization of the flux can still be partially compensated at low energies by different interaction cross sections for neutrinos.
The event distributions as a function of the muon pathlength inside the detector have also been studied, as an independent consistency check. In Fig. 5 the data (black points) are compared with the MC expectation (solid lines for no neutrino oscillations; dashed lines for oscillations at our test point). The shapes are similar, and the data prefer the reduced normalization of the oscillation prediction.
We estimated the most likely values of $`\mathrm{\Delta }m^2`$ and $`sin^22\theta _{mix}`$ using a $`\chi ^2`$ method for data and Monte Carlo for the data of Fig. 4. The $`\chi ^2`$ was computed with ten degrees of freedom: the histograms ($`2\times 4`$ bins, normalized so as to contain only distribution shape information), the $`\frac{IU}{ID+UGS}`$ ratio and the overall normalization. The statistical and systematic errors are added in quadrature; the systematic uncertainty is 10% in each bin of the angular distributions, 5% for the ratio, and 25% for the normalization. Fig. 6 shows the 90% confidence level region, based on the application of the MC prescriptions of Ref. on a $`(sin^22\theta _{mix},\mathrm{\Delta }m^2)`$ grid. The expected flux $`\mathrm{\Phi }_{osc}^i`$ for a given point of $`(sin^22\theta _{mix},\mathrm{\Delta }m^2)`$ in the grid is obtained by weighting each simulated event with its surviving probability $`P(\nu _\mu \nu _\mu )=1sin^22\theta _{mix}sin^2(1.27\mathrm{\Delta }m^2L/E_\nu )`$ in that bin. The maximum of the $`\chi ^2`$ probability (97%) occurs at $`sin^22\theta _{mix}=1.0`$; this $`\chi ^2`$ probability is almost constant in the interval $`\mathrm{\Delta }m^2=(1÷20)\times 10^3eV^2`$. In the region of the maximum, the flux normalization factor is $`1.02`$ in both data sets ( i.e. the data are 2% higher than the oscillated predictions with our normalization).
In conclusion, we presented measurements of two samples of events induced by relatively low-energy neutrinos ($`\overline{E}_\nu 4GeV`$) interacting in MACRO or in the surrounding rock. The neutrinos originate from cosmic ray interactions in the upper atmosphere above the detector (downgoing events) or up to $`13000km`$ away on the opposite side of the Earth (upgoing events). The ratio of the number of observed to expected events (no oscillations) is $`0.57\pm 0.05_{stat}\pm 0.06_{syst}\pm 0.14_{theor}`$ for the $`IU`$ events and $`0.71\pm 0.05_{stat}\pm 0.07_{sys}\pm 0.18_{theo}`$ for the $`ID+UGS`$. Within statistics, the observed deficits are uniform over the zenith angle. From the ratio of $`IU`$ to $`ID+UGS`$, the probability that there is an overall reduction in the number of neutrino-induced muons is 5%. The hypothesis of muon neutrino oscillations explains the different deficits in $`IU`$ and $`ID+UGS`$ events with higher probability. The large theoretical uncertainties for the neutrino flux and cross sections is dominant in each data set; the ratio of the two low energy samples is dominated by statistical uncertainties. The regions with $`\mathrm{\Delta }m^2>3\times 10^4eV^2`$ and $`sin^22\theta _{mix}>0.25`$ are allowed at $`90\%`$ C.L. The best region corresponds to maximal mixing and $`\mathrm{\Delta }m^2=(1÷20)\times 10^3eV^2`$. This result confirms the scenario proposed by the measurement of higher-energy neutrino-induced muons by MACRO , as well as by other experiments , all of which favor the $`\nu _\mu `$ oscillation hypothesis with maximal mixing and $`\mathrm{\Delta }m^2`$ of a few times $`10^3eV^2`$.
We acknowledge the support of the staff of the Gran Sasso Laboratory and of the Institutions participating in the experiment. We thank the Istituto Nazionale di Fisica Nucleare (INFN), the U.S. Department of Energy and the U.S. National Science Foundation for their support. We thank INFN, FAI, ICTP (Trieste), NATO and WorldLab for providing fellowships and grants for non-Italian citizens. |
warning/0001/nlin0001017.html | ar5iv | text | # Superconvergence of period doubling cascade in trapezoid maps
## 1 Introduction
It is well known that in a class of one-dimensional map, as a parameter is changed, the period doubling bifurcation cascades, and there exist several universal properties . One of the universal quantities is the so called Feigenbaum constant $`\delta `$ and it depends only on the exponent $`z`$ characterizing the behavior of the map around the critical point.
About a decade ago, there had been controversy on the $`z\mathrm{}`$ limit of $`\delta (z)`$. As a typical example, the map $`x_{n+1}=f(x_n)=1a|x_n|^z`$ has been extensively studied and two different values were conjectured for this limit, one is finite and the other is infinity . This problem was solved by J. P. Eckmann and H. Epstein and they proved that the limit is finite .
In the previous paper, instead of considering a map with finite $`z`$ and taking $`z\mathrm{}`$, we investigated a map with $`z=\mathrm{}`$. As such maps, we treated the symmetric and asymmetric trapezoid maps.
The symmetric trapezoid map $`x_{n+1}=T_{(a,b)}(x_n)`$ defined in $`[0,1]`$ is given by
$$T_{(a,b)}(x)=\{\begin{array}{cc}ax\hfill & \text{for }0x(1b)/2,\hfill \\ a(1b)/2\hfill & \text{ for }(1b)/2x(1+b)/2\text{ },\hfill \\ a(1x)\hfill & \text{ for }(1+b)/2x1\text{ },\hfill \end{array}$$
(1.1)
where $`0<b<1`$ and $`a>0`$. See Fig.1.
As $`a`$ is increased from 1 with fixed $`b`$, the successive period doubling bifurcations take place. We assign the symbol $`L,C`$ and $`R`$ to an orbit $`x_i`$ according to the following rule; $`L\mathrm{for}x_i[0,(1b)/2]I_L,C\mathrm{for}x_i((1b)/2,(1+b)/2)I_C,R\mathrm{for}x_i[(1+b)/2,1]I_R`$. We denote this correspondence by $`H(x_i)`$. Further, we define the conjugate of $`R`$ or $`L`$ as $`\overline{R}=L`$ or $`\overline{L}=R`$, respectively. Then, in the previous paper we proved that as a slope $`a`$ of the trapezoid map is increased, the period doubling bifurcation cascades and the symbolic sequence is the Metropolis-Stein-Stein sequence $`R^m`$ as usual. And, defining $`a_m`$ as the onset of a $`2^m`$-cycle, we calculated $`\delta _m\frac{a_{m+1}a_m}{a_{m+2}a_{m+1}}`$ and proved that $`\delta _m`$ scales as
$$\delta _ma_c^{2^m}$$
where $`a_c`$ is the accumulation point of the period doubling cascade. We called this phenomena the super-convergent period doubling cascade.
Further, we investigated the asymmetric trapezoid map $`x_{n+1}=A_{(a,b,\gamma )}(x_n)`$ defined in ,
$$A_{(a,b,\gamma )}(x)=\{\begin{array}{cc}ax\hfill & \text{for }0x\alpha ,\hfill \\ \alpha a\hfill & \text{ for }\alpha x\beta \text{ },\hfill \\ \gamma a(1x)\hfill & \text{ for }\beta x1\text{ },\hfill \end{array}$$
(1.2)
where $`\gamma `$ is the ratio of two slopes of the trapezoid, $`\alpha =\frac{\gamma }{1+\gamma }(1b)`$ and $`\beta =\alpha +b=\frac{b+\gamma }{1+\gamma }`$. When $`\gamma =1`$, $`A_{(a,b,\gamma )}(x)`$ is reduced to $`T_{(a,b)}(x)`$. In the asymmetric map, when $`\gamma `$ and $`b`$ are fixed and $`a`$ is increased, we proved similar results as in the symmetric map, in particular, as the scaling of $`\delta _m`$ we obtained
$$\delta _m\gamma ^{(1)^m/3}(a_c\gamma ^{2/3})^{2^m}.$$
Finally, we gave approximate expressions for the accumulation point as functions of $`b`$ which is the length of the smaller side of the trapezoid in both symmetric and asymmetric cases.
In this paper, we give the detailed and complete description of the proofs given in the previous paper. New results in this paper are the scaling relations for the period doubling bifurcation starting with a period $`p(3)`$ solution which appears by a tangent bifurcation.
In the following section, we treat the symmetric trapezoid map, and then in the section 3, we treat the asymmetric map. In sections 4 and 5, we study the period doubling bifurcation for period $`p(3)`$ solution in symmetric and asymmetric cases, respectively. We give summary and discussion in the last section.
## 2 The symmetric case
We consider $`T_𝒂(x)T_{(a,b)}(x)`$. For brevity we define $`𝒂=(a,b)`$. Let $`x_M`$ be the maximum value of $`T_𝒂(x)`$, i.e., $`x_Ma\frac{1b}{2}`$. Let $`x_{0,1}`$ be the non-zero fixed point of $`T_𝒂(x)`$. Then the following lemma is easily proved.
lemma 1
For $`1<a<a_1`$, $`T_𝒂(x)`$ has the stable fixed point $`x_{0,1}=x_M`$, which satisfies $`\frac{1b}{2}<x_M<\frac{1+b}{2}`$. $`a_1`$ is defined by the equation $`x_M=\frac{1+b}{2}`$, that is $`a_1=\frac{1+b}{1b}`$.
For $`a>a_1`$, $`T_𝒂(x)`$ has the unstable fixed point $`x_{0,1}=x^{}\frac{a}{a+1}`$ and $`x^{}>\frac{1+b}{2}`$. At $`a=a_1,x_M=x^{}.`$
See Fig.2.
In the region $`a>a_1`$, by iterating $`T_𝒂(x)`$ twice and rescaling $`x`$ such that the domain becomes $`[0,1]`$, we obtain a trapezoid map with different parameter $`𝒂^{(1)},𝒂^{(1)}=(a^{(1)},b^{(1)})`$.
See Fig.3 and 4. In fact, we obtain the following lemma.
lemma 2
When $`T_𝒂(x)`$ has the non-zero unstable fixed point $`x^{}=\frac{a}{a+1}(>\frac{1}{2})`$, that is, for $`a>a_1`$, by rescaling the coordinate $`x`$ as $`x^{(1)}=\frac{x^{}x}{2x^{}1}`$ and restricting $`x`$ to the interval $`[1x^{},x^{}]`$, $`T_𝒂^2(x)`$ is transformed to $`T_{𝒂^{(1)}}(x^{(1)})`$ which is defined for $`x^{(1)}`$ in $`[0,1]`$. Here $`𝒂^{(1)}𝝋(𝒂)(a^2,u(a)b)`$ and $`u(a)=\frac{a+1}{a1}`$.
The proof is straightforward. $`𝝋(𝒂)`$ is defined as long as $`a1`$. In this paper, we restrict ourselves to the case of $`a>1`$.
Now, we define the $`m`$-th iteration of $`𝝋`$. That is,
$`𝒂^{(m)}`$ $``$ $`(a^{(m)},b^{(m)})𝝋^m(𝒂)=(a^{2^m},u_m(a)b),`$ (2.1)
$`u_m(a)`$ $``$ $`{\displaystyle \underset{l=0}{\overset{m1}{}}}u(a^{2^l})={\displaystyle \underset{l=0}{\overset{m1}{}}}{\displaystyle \frac{a^{2^l}+1}{a^{2^l}1}}.`$
These are defined as long as $`a1`$. For $`a>1`$, as functions of $`a`$, $`a^{(m)}`$ is a continuous strictly increasing function and $`b^{(m)}`$ is a continuous strictly decreasing function. $`lim_a\mathrm{}b^{(m)}=b<1`$ for any $`m>0`$. Let us prove the following lemma.
lemma 3
For any positive integer $`m`$, there exists the unique value of $`a=a_m`$ such that
$$b^{(m)}(a_m)=1.$$
(2.2)
$`\{a_m\}_{m=1}^{\mathrm{}}`$ is an increasing sequence, i.e., $`1<a_1<a_2<\mathrm{}`$. Further, the relation $`b<b^{(m)}<1`$ holds for $`a_m<a`$ for $`m1`$.
Proof
Let us consider the case of $`m=1`$.
$`b^{(1)}(a)=1`$ has the unique solution $`a=\frac{1+b}{1b}>1`$, and this is $`a_1`$ defined in lemma 1.
Next, let us assume that $`b^{(m)}(a_m)=1`$ and $`a_m>1`$ for $`m1`$. Since $`b^{(m+1)}(a)=\frac{a^{(m)}+1}{a^{(m)}1}b^{(m)}(a)`$, $`b^{(m+1)}(a_m)>1`$ follows. Thus, there exists the unique value of $`a=a_{m+1}(>a_m)`$ such that $`b^{(m+1)}(a_{m+1})=1`$.
Therefore, from the mathematical induction, $`b^{(m)}(a)=1`$ has a unique solution $`a_m`$ for any positive integer $`m`$ and $`\{a_m\}_{m=1}^{\mathrm{}}`$ is an increasing sequence. The inequality $`b<b^{(m)}<1`$ for $`a_m<a`$ is immediately follows from the fact that the function $`b^{(m)}(a)`$ is strictly decreasing. Q.E.D.
For positive integer $`m`$, we define $`T_{𝝋^m(𝒂)}(x^{(m)})`$ for $`a>a_m`$ from $`T_{𝝋^{m1}(𝒂)}(x^{(m1)})`$ successively by the same procedure as in the lemma2.
We define $`x_M^{(m)}a^{(m)}\frac{1b^{(m)}}{2}`$, which is the maximum value of $`T_{𝝋^m(𝒂)}(x^{(m)})`$. Further, we define $`x^{(m)}\frac{a^{(m)}}{a^{(m)}+1}`$. $`x_M^{(0)}=x_M`$ and $`x^{(0)}=x^{}`$.
lemma 4
For any non-negative integer $`m`$ and for $`a_m<a`$ there exists a unique non-zero fixed point for $`T_{𝝋^m(𝒂)}(x^{(m)})`$ which is defined in $`[0,1]`$. For $`a_m<a<a_{m+1}`$, $`\frac{1b^{(m)}}{2}<x_M^{(m)}<\frac{1+b^{(m)}}{2}`$ and $`x_M^{(m)}`$ is the stable fixed point of $`T_{𝝋^m(𝒂)}(x^{(m)})`$. At $`a=a_{m+1}`$, $`x_M^{(m)}=x^{(m)}`$. For $`a_{m+1}<a`$, the non-zero fixed point for $`T_{𝝋^m(𝒂)}(x^{(m)})`$ is $`x^{(m)}`$ and unstable, and $`x^{(m)}>\frac{1+b^{(m)}}{2}`$. Here, we define $`x^{(0)}x,a_01,a^{(0)}a`$ and $`b^{(0)}b`$.
Proof
First of all, we notice that for $`m1`$ $`T_{𝝋^m(𝒂)}(x^{(m)})`$ is really a trapezoid map because from lemma 3 $`b<b^{(m)}<1`$ for $`a_m<a`$ and for $`m1`$.
For the case $`m=0`$, the statement follows from lemma 1.
Let us assume that the statement holds for $`m0`$. Since $`x^{(m)}=\frac{a^{(m)}}{a^{(m)}+1}`$ is the unstable fixed point of $`T_{𝒂^{(m)}}(x^{(m)})`$ for $`a_{m+1}<a`$, from lemma 2, taking $`T_{𝒂^{(m)}}^2(x^{(m)})`$ and rescaling $`x^{(m)}`$ as $`x^{(m+1)}=\frac{x^{(m)}x^{(m)}}{2x^{(m)}1}`$, we obtain $`T_{𝒂^{(m+1)}}(x^{(m+1)})`$ which is defined in $`[0,1]`$. The following equivalence relations are easily proved for $`m0`$,
$$b^{(m+1)}(a)=1x_M^{(m)}=\frac{1+b^{(m)}}{2}a^{(m)}=\frac{1+b^{(m)}}{1b^{(m)}}.$$
(2.3)
Let us consider $`x_M^{(m+1)}=a^{(m+1)}\frac{1b^{(m+1)}}{2}`$. From the relation (5), since $`b^{(m+2)}(a_{m+2})=1`$ it follows that $`x_M^{(m+1)}=\frac{1+b^{(m+1)}}{2}`$ at $`a=a_{m+2}`$. Since $`x_M^{(m+1)}`$ is the strictly increasing function w.r.t. $`a`$, we obtain $`0<\frac{1b^{(m+1)}}{2}<x_M^{(m+1)}<\frac{1+b^{(m+1)}}{2}`$ for $`a_{m+1}<a<a_{m+2}`$. This implies $`x_M^{(m+1)}`$ is the unique non-zero fixed point for $`T_{𝝋^{m+1}(𝒂)}(x^{(m+1)})`$ and stable. It is easily shown that for $`a_{m+2}<a`$, $`x_M^{(m+1)}`$ is no more fixed point but $`x^{(m+1)}=\frac{a^{(m+1)}}{a^{(m+1)}+1}`$ becomes unstable fixed point and $`x^{(m+1)}>\frac{1+b^{(m+1)}}{2}`$. At $`a=a_{m+2},a_M^{(m+1)}=x^{(m+1)}`$ holds. Thus, by the mathematical induction, the proof completes.
lemma 5
The unique non-zero fixed point of $`T_{𝝋^m(𝒂)}(x^{(m)})`$ in lemma 4 is a member of a periodic cycle with prime period $`2^m`$ of $`T_𝒂(x)`$.
Proof
The case of $`m=0`$, it is trivial. Now, let us fix $`m(0)`$ and assume that for $`0km`$ the unique non-zero fixed point of $`T_{𝝋^k(𝒂)}(x^{(k)})`$ which exists for $`a_k<a`$ is the member of the cycle with prime period $`2^k`$. Thus, from lemma 4 for $`a_{m+1}<a`$ the unique non-zero fixed point of $`T_{𝝋^{m+1}(𝒂)}(x^{(m+1)})`$ does not correspond to any periodic points with period less than $`2^{m+1}`$ because it is stable for $`a_{m+1}<a<a_{m+2}`$. Thus, the non-zero fixed point of $`T_{𝝋^{m+1}(𝒂)}(x^{(m+1)})`$ is the member of the cycle with prime period $`2^{m+1}`$. This completes the proof.
Thus, we obtain the following theorem.
Theorem 1
$`x=0`$ is period 1 solution of $`T_𝒂(x)`$ for any $`a>0`$. Except for this, for any non-negative integer $`m`$, for $`a_m<a<a_{m+1}`$ there is the unique cycle with prime period $`2^k`$ for $`k=0,1,\mathrm{},m`$ and no other periodic points exist. The periodic solution with prime period $`2^m`$ is stable for $`a_m<a<a_{m+1}`$ and unstable for $`a_{m+1}<a`$. At $`a=a_{m+1}`$, the period $`2^m`$ cycle and the period $`2^{m+1}`$ cycle coincide.
Proof
For the case of $`m=0`$, the statement is trivial. Let us assume that the statement holds for $`0,1,\mathrm{},m`$. Then we only have to prove that for $`a_{m+1}<a<a_{m+2}`$, except for the period $`2^k`$ cycles for $`k=1m+1`$, there is no other cycle. Let us consider any initial point $`x`$ in (0, 1) which does not corresponds to the unstable cycles of period $`2^k,(k=0m).`$ Then, the orbit $`T_𝒂^i(x)`$ finally enters the domain of the map $`T_{𝝋^m(𝒂)}(x^{(m)})`$ if $`i`$ is appropriately chosen. Thus, since any point except for $`x^{(m)}=0`$ and 1 converges to the fixed point of $`T_{𝝋^m(𝒂)}(x^{(m)})`$ for $`a_{m+1}<a<a_{m+2}`$, the nonexistence of other periodic points for $`a_{m+1}<a<a_{m+2}`$ follows. The latter part of the theorem follows from lemma 4 immediately.
From the above arguments we conclude that as $`a`$ is increased from 1, the period doubling bifurcation cascades, and at $`a=a_m`$ the periodic solution with prime period $`2^m`$ appears for $`m0`$.
Now, let us investigate the upper bound of the sequence $`\{a_m\}`$. $`x_M^{(m)}(a)=a^{(m)}\frac{1b^{(m)}}{2}`$ is defined for $`a>1`$ and the strictly increasing function w.r.t $`a`$. For any $`m>0,x_M^{(m)}(a_m)=0`$ and $`lim_a\mathrm{}x_M^{(m)}(a)=\mathrm{}`$ follows. Thus, for $`m>0`$ we define $`a_M(m)`$ as the unique solution of the equation
$$x_M^{(m)}(a)=1.$$
(2.4)
For $`m=0`$, we define $`a_M(0)=\frac{2}{1b}`$. We abbreviate this as $`a_M`$. Then, we obtain the following lemma.
lemma 6
$`1<a_1<a_2<\mathrm{}<a_m<a_{m+1}<\mathrm{}`$ (2.5)
$`<a_M(m+1)<a_M(m)\mathrm{}<a_M(2)<a_M(1)<a_M.`$
Proof
For $`m>0`$, since $`x_M^{(m)}(a_m)=0,a_m<a_M(m)`$ follows. For $`m0`$, $`x_M^{(m+1)}(a)`$ is rewritten as
$$x_M^{(m+1)}(a)=\{a^{(m)}\}^2\frac{1}{a^{(m)}1}(x_M^{(m)}(a)\frac{1+b^{(m)}}{2}).$$
(2.6)
Then, $`x_M^{(m+1)}(a_M(m))=\frac{a^{(m)}}{a^{(m)}1}>1.`$ Therefore, $`a_M(m+1)<a_M(m)`$ for $`m0`$. Thus, we obtain $`a_m<a_{m+1}<a_M(m+1)<a_M(m)`$ for $`m0`$. Q.E.D.
From this,
$$a_c\underset{m\mathrm{}}{lim}a_m\underset{m\mathrm{}}{lim}a_M(m)<a_M$$
(2.7)
follows. Therefore $`a_c`$ is finite.
Now, let us investigate the symbolic sequence. The sequence of $`R`$ and $`L`$ for the period $`2^m`$ solution becomes Metropolis-Stein-Stein(MSS) sequence. Let us prove this.
For $`m0`$, we define the onset of the $`2^n`$-cycle for $`T_{𝝋^m(𝒂)}(x^{(m)})`$ as $`a_n^{(m)}`$ and the $`i`$-th orbit of the period $`2^n`$ cycle as $`x_{n,i}^{(m)}`$. $`a_n^{(0)}=a_n`$. For $`m=0`$, we often omit the superscript $`(0)`$. As $`x_{m,1}`$ we set the largest $`x`$ value among $`2^m`$ members of the periodic cycle. That is, $`x_{m,1}=x_M=a(1b)/2`$ when the cycle is stable. Further, for $`m0`$ we divide the coordinate space $`[0,1]`$ of $`x^{(m)}`$ into the three intervals as $`I_L^{(m)}[0,\frac{1b^{(m)}}{2}],I_C^{(m)}(\frac{1b^{(m)}}{2},\frac{1+b^{(m)}}{2})`$ and $`I_R^{(m)}[\frac{1+b^{(m)}}{2},1].`$
lemma 7
For any positive integer $`m`$, when $`a_m<a`$,
1. $`I_C^{(m)}x^{(m)}`$ is equivalent to $`I_C^{(m1)}x^{(m1)}`$,
2. if $`I_L^{(m)}x^{(m)}`$ then $`I_R^{(m1)}x^{(m1)}`$,
3. if $`I_R^{(m)}x^{(m)}`$ then $`I_L^{(m1)}x^{(m1)}`$,
where $`x^{(m)}`$ and $`x^{(m1)}`$ are related by the coordination transformation used to define $`T_{𝝋^m(𝒂)}(x^{(m)})`$ from $`T_{𝝋^{m1}(𝒂)}(x^{(m1)})`$.
proof
For $`a_m<a`$, $`x^{(m)}=0,\frac{1b^{(m)}}{2},\frac{1+b^{(m)}}{2}`$ and 1 correspond to $`x^{(m1)}=x^{(m1)},\frac{1+b^{(m1)}}{2},\frac{1b^{(m1)}}{2}`$ and $`1x^{(m1)}`$, respectively. Since this correspondence is linear, the statements hold.
lemma 8
The fixed point $`x_{0,1}^{(m)}`$ of $`T_{𝝋^m(𝒂)}(x^{(m)})`$ corresponds to $`x_{m,2^m}`$.
proof
Let $`x_{m,i}`$ be the orbit of $`2^m`$ solution corresponding to $`x_{0,1}^{(m)}`$. From lemma 4, for $`a_m<a<a_{m+1},x_{0,1}^{(m)}=x_M^{(m)}`$ and $`I_C^{(m)}x_{0,1}^{(m)}`$. Then, from lemma 7, $`I_Cx_{m,i}`$ follows. Therefore, $`x_{m,i}`$ is mapped to $`x_M`$ by $`T_𝒂`$, which is $`x_{m,1}`$. Thus, $`i=2^m`$.
lemma 9
For any positive integer $`m`$, symbols $`H(x_{m,i})`$ for the number of $`2^m`$ cycle satisfy the followings.
1. $`H(x_{m,i})`$ does not change for $`a_ma`$ and for $`i=1,\mathrm{},2^m1`$.
2. For even $`m`$,
$$H(x_{m,2^m})=\{\begin{array}{cc}L\hfill & \text{for }a=a_m,\hfill \\ C\hfill & \text{ for }a_m<a<a_{m+1}\text{ },\hfill \\ R\hfill & \text{ for }a_{m+1}a\text{ }.\hfill \end{array}$$
and for odd $`m`$,
$$H(x_{m,2^m})=\{\begin{array}{cc}R\hfill & \text{for }a=a_m,\hfill \\ C\hfill & \text{ for }a_m<a<a_{m+1}\text{ },\hfill \\ L\hfill & \text{ for }a_{m+1}a\text{ }.\hfill \end{array}$$
3. $`\{H(x_{m,i})\}`$ is the MSS sequence $`R^m`$ <sup>)</sup><sup>)</sup>) Strictly speaking, the last symbol of the sequence in our definition is $`R`$ or $`L`$ and is different from that in the MSS sequence, $`C`$ for $`a_{m+1}a`$.
proof
Let us consider the case of $`m=1`$. At $`a=a_1`$, $`x_{0,1}=\frac{1+b}{2}I_R`$. Then, at $`a=a_1,H(x_{0,1})=H(x_{1,1})=H(x_{1,2})=R`$. For $`a_1<a<a_2`$, $`x_{0,1}^{(1)}=x_M^{(1)}=\frac{a^{(1)}(1b^{(1)})}{2}`$ and for $`a_2a,x_{0,1}^{(1)}=x^{(1)}`$. From lemma 8, $`x_{0,1}^{(1)}`$ corresponds to $`x_{1,2}`$. Thus, from lemma 7, $`H(x_{1,2})`$ is $`C`$ for $`a_1<a<a_2`$, and is $`L`$ for $`a_2a`$. On the other hand, for $`a_1a`$, since $`x_{1,1}x_{0,1}=\frac{a}{a+1}`$ then $`H(x_{1,1})=R`$. Therefore, $`H(x_{1,1})H(x_{1,2})=RL`$ for $`a_2a`$. Therefore, for $`m=1`$, the lemma holds.
Next, we assume that the lemma holds for the case of $`m(1)`$. At $`a=a_{m+1}`$ , $`x_{m+1,i}`$ and $`x_{m+1,i+2^m}`$ emerge from $`x_{m,i}(i=12^m)`$. Then, at $`a=a_{m+1}`$, $`H(x_{m,i})=H(x_{m+1,i})=H(x_{m+1,i+2^m})`$ for $`i=1,\mathrm{},2^m`$. These are $`L`$ or $`R`$. Let us assume that at some value of $`a(>a_{m+1})`$, $`H(x_{m+1,i})`$ becomes $`C`$. Then, $`x_{m+1,i}`$ is mapped to $`x_M`$ by $`T_𝒂`$. Thus, in this case, $`a<a_{m+2}`$ and $`x_{m+1,i+1}=x_M=x_{m+1,1}`$ follow. So, $`i`$ should be $`2^{m+1}`$. Thus, $`x_{m+1,i}(i=12^{m+1}1)`$ does not change its symbol for $`a_{m+1}a`$. From lemma 8 $`x_{m+1,2^{m+1}}`$ corresponds to $`x_{0,1}^{(m+1)}=x_M^{(m+1)}`$ for $`a_{m+1}aa_{m+2}`$. Thus, $`H(x_{m+1,2^{m+1}})=C`$ for $`a_{m+1}<a<a_{m+2}`$. For $`a_{m+2}a`$, $`x_{0,1}^{(m+1)}=x^{(m+1)}I_R^{(m+1)}`$. Thus, from lemma 7, for $`a_{m+2}a`$,
$$H(x_{m+1,2^{m+1}})=\{\begin{array}{cc}L\hfill & \text{ for odd }m+1,\hfill \\ R\hfill & \text{ for even }m+1.\hfill \end{array}$$
Therefore, the statements 1 and 2 hold for $`m+1`$.
At $`a=a_{m+1},H(x_{m+1,2^{m+1}})=H(x_{m,2^m})`$. Then, from the assumption,
$$H(x_{m+1,2^{m+1}})=\{\begin{array}{cc}R\hfill & \text{ for odd }m+1,\hfill \\ L\hfill & \text{ for even }m+1.\hfill \end{array}$$
Let $`H(x_{m,1})H(x_{m,2})\mathrm{}H(x_{m,2^m})`$ be the MSS sequence for $`aa_{m+1}`$. From the above argument, for $`aa_{m+2}`$, the sequence for $`2^{m+1}`$ cycle is
$`H(x_{m,1})H(x_{m,2})\mathrm{}H(x_{m,2^m1})H(x_{m,2^m})`$
$`\times H(x_{m,1})H(x_{m,2})\mathrm{}H(x_{m,2^m1})\overline{H(x_{m,2^m})}.`$
This is the MSS sequence. Thus, the statement 3 is proved. This completes the proof.
Now, we derive the equations for which $`a_m`$ and $`a_c`$ should satisfy, respectively, and obtain the asymptotic expression for $`\delta _m`$. First, we assign 0 or 1 to any orbit $`x_i`$ with the symbol $`L`$ or $`R`$, respectively. We denote this correspondence as $`s_i=s(x_i)`$. Further, we define the function $`f_s(x)`$ for $`xI_LI_R`$ as
$$f_s(x)=sa+(12s)ax,$$
(2.8)
where $`s=s(x)`$. Starting from the maximum value of $`T_{a,b}(x)`$, $`x_1=x_M=a(1b)/2`$, if $`x_1,x_2,\mathrm{},x_{n1}I_LI_R`$, $`x_n`$ is expressed as
$`x_n`$ $`=`$ $`f_{s_{n1}}(x_{n1})=f_{s_{n1}}f_{s_{n2}}\mathrm{}f_{s_1}(x_1)={\displaystyle \underset{l=1}{\overset{n}{}}}\xi _la^l,`$ (2.9)
$`\xi _l`$ $`=`$ $`s_{nl}{\displaystyle \underset{j=nl+1}{\overset{n1}{}}}(12s_j)\mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}2}ln1,`$
$`\xi _1`$ $`=`$ $`s_{n1},\xi _n={\displaystyle \frac{1b}{2}}{\displaystyle \underset{j=1}{\overset{n1}{}}}(12s_j).`$
For $`m1`$ let us consider the period $`2^m`$ solution $`\{x_{m,i}\}`$ for $`a_maa_{m+1}`$.
$`x_{m,2^m}`$ $`=`$ $`f_{s_{2^m1}}(x_{m,2^m1}){\displaystyle \underset{l=0}{\overset{2^m1}{}}}c_l^{(m)}a^{2^ml}F_m(a),`$ (2.10)
$`c_l^{(m)}`$ $`=`$ $`s_l{\displaystyle \underset{j=l+1}{\overset{2^m1}{}}}(12s_j)=\xi _{2^ml},\mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}l2^m2,`$
$`c_0^{(m)}`$ $`=`$ $`{\displaystyle \frac{1b}{2}}{\displaystyle \underset{j=1}{\overset{2^m1}{}}}(12s_j)=s_0{\displaystyle \underset{j=1}{\overset{2^m1}{}}}(12s_j),`$
$`c_{2^m1}^{(m)}`$ $`=`$ $`s_{2^m1},`$
where $`s_0\frac{1b}{2}`$. Note that $`F_m(a)`$ is determined by the sequence $`(s_1,\mathrm{},s_{2^m1})`$ and is independent of $`s_{2^m}`$. From lemma 9, the symbols $`H(x_{m,i})(i=12^m1)`$ do not change for $`a_ma`$. Further, for any positive integer $`j`$ these symbols are equal to the first $`2^m1`$ symbols of the period $`2^{m+j}`$ solution when $`a_{m+j}a`$. Therefore, $`F_m(a)`$ expresses $`x_{2^{m+j},2^m}`$ in the region $`[a_{m+j},a_{m+j+1}]`$. From the statement 2 in lemma 9, the following relations follow
$`x_{m,2^m}(a_m)`$ $`=`$ $`{\displaystyle \frac{1+(1)^{m1}b}{2}},`$ (2.11)
$`x_{m,2^m}(a_{m+1})`$ $`=`$ $`{\displaystyle \frac{1+(1)^mb}{2}}.`$ (2.12)
Thus, we obtain the following conditions for $`a_m`$,
$$F_m(a_m)=\frac{1+(1)^{m1}b}{2},$$
(2.13)
or
$$F_{m1}(a_m)=\frac{1+(1)^{m1}b}{2}.$$
(2.14)
Since the symbolic sequence is the Metropolis-Stein-Stein sequence $`R^m`$,
$`\mathrm{\Pi }_{j=1}^{2^m1}(12s_j)`$ $`=`$ $`(1)^m,`$ (2.15)
$`s_{2^m}`$ $`=`$ $`{\displaystyle \frac{1+(1)^m}{2}}`$ (2.16)
follow. From these relations, for $`l0`$ we obtain
$$r_l\{\underset{j=1}{\overset{2^m1}{}}(12s_j)\}^1c_l^{(m)}=(1)^mc_l^{(m)}=s_l\underset{j=1}{\overset{l}{}}(12s_j).$$
That is, $`r_l`$ is $`m`$-independent as long as it is defined. Thus, we get
$`r_l`$ $`=`$ $`s_l{\displaystyle \underset{j=1}{\overset{l}{}}}(12s_j)\text{ for any }l(0),`$ (2.17)
$`r_0`$ $`=`$ $`(1b)/2,r_{2^m}=(1+(1)^m)/2\text{ for any }m(0),`$ (2.18)
and for $`l>0`$ the successive values $`(r_l,r_{l+1})`$ take the following six sets of values, $`(0,\pm 1),(\pm 1,0),(1,1),(1,1)`$. <sup>)</sup><sup>)</sup>)In paper , two cases $`(0,1)`$ and $`(1,0)`$ are missing. Let us define
$$G_m(a)(1)^ma^{2^m}F_m(a)=\underset{l=0}{\overset{2^m1}{}}r_la^l.$$
(2.19)
$$G_{\mathrm{}}(z)\underset{m\mathrm{}}{lim}G_m(z).$$
(2.20)
Then, $`G_{\mathrm{}}(z)`$ is the analytic function for $`|z|>1`$. The equation (2$``$13) becomes
$$G_m(a_m)=a_m^{2^m}\frac{(1)^mb}{2},$$
(2.21)
and the accumulation point $`a_c`$ satisfies the equation,
$$G_{\mathrm{}}(a_c)=0.$$
(2.22)
Let us estimate $`a_m`$ for large $`m`$. Putting $`a_m=a_cϵ_m`$ and using the mean value theorem, the right hand side of eq.(2$``$21) is rewritten as
$$a_m^{2^m}\frac{(1)^mb}{2}=(a_c^{2^m}+2^m\widehat{a}_m^{2^m1}ϵ_m)\frac{(1)^mb}{2},$$
where $`a_m<\widehat{a}_m<a_c`$. On the other hand, the left hand side of eq.(2$``$21) is expressed as
$$G_m(a_m)=G_m(a_c)G_m^{}(\overline{a}_m)ϵ_m,$$
where $`a_m<\overline{a}_m<a_c`$. Thus, we obtain from eq.(2$``$22)
$$ϵ_m=\{G_m(a_c)+\frac{b(1)^m}{2}a_c^{2^m}\}(1+\overline{h}_m)/G_{\mathrm{}}^{}(a_c)$$
(2.23)
where $`\overline{h}_m=G_{\mathrm{}}^{}(a_c)/(G_m^{}(\overline{a}_m)+2^m(\widehat{a}_m)^{2^m1}\frac{(1)^mb}{2})1`$ and $`lim_m\mathrm{}\overline{h}_m=0`$. Here, we assume $`G_{\mathrm{}}^{}(a_c)0`$, which is proved later. Equation ($`222`$) is rewritten as
$$G_{\mathrm{}}(a_c)=G_m(a_c)+\underset{l=0}{\overset{\mathrm{}}{}}r_{2^m+l}a_c^{l2^m}=0.$$
(2.24)
Using the relation
$$r_{2^m+l}=r_l\text{for }1l2^m1,$$
(2.25)
eq.($`224`$) is further changed to the following.
$$G_{\mathrm{}}(a_c)=G_m(a_c)(1a_c^{2^m})+a_c^{2^m}(r_0+r_{2^m}+a_c^{2^m}\underset{l=0}{\overset{\mathrm{}}{}}r_{2^{m+1}+l}a_c^l)=0.$$
(2.26)
Thus, we get
$$G_m(a_c)=\frac{b+(1)^m}{2}a_c^{2^m}+(a_c^{2^m})^2q_m,$$
(2.27)
where $`q_m=\frac{1}{1a_c^{2^m}}(\frac{b+(1)^m}{2}_{l=0}^{\mathrm{}}r_{2^{m+1}+l}a_c^l)`$ and $`|q_m|<\frac{a_c(2a_c1)}{(a_c1)^2}`$. Substituting eq.(2$``$27) into eq.(2$``$23) we obtain,
$$ϵ_m=\frac{ba_c^{2^m}}{G_{\mathrm{}}^{}(a_c)}(1+h_m),$$
(2.28)
where $`h_m=\overline{h}_m(1+\frac{q_m}{b}a_c^{2^m})+\frac{q_m}{b}a_c^{2^m}`$ and $`lim_m\mathrm{}h_m=0`$. Thus,
$`\delta _m`$ $`=`$ $`{\displaystyle \frac{ϵ_mϵ_{m+1}}{ϵ_{m+1}ϵ_{m+2}}}=a_c^{2^m}(1+l_m),`$
$`l_m=[h_mh_{m+1}a_c^{2^m}\{1+h_{m+1}a_c^{2^m}(1+h_{m+2})\}]`$
$`/[1+h_{m+1}a_c^{2^{m+1}}(1+h_{m+2})],`$
and $`lim_m\mathrm{}l_m=0`$.
Next, we give the alternative relation for the onset point $`a_m`$. As is shown in lemma 3, at $`a=a_m`$ we have the following equation (2$``$2) for $`m>0`$,
$$b^{(m)}(a_m)=u_m(a_m)b=1.$$
Defining $`v(a)\frac{a1}{a+1}=1/u(a)`$ and $`v_m(a)1/u_m(a)=_{l=0}^{m1}v(a^{2^l})`$, we obtain for $`m>0`$
$$v_m(a_m)=b.$$
(2.30)
For $`m1`$, $`v_m(a)`$ is rewritten as
$`v_m(a)`$ $`=`$ $`{\displaystyle \frac{1a^1}{1+a^{2^{m1}}}}{\displaystyle \underset{l=0}{\overset{m2}{}}}(1a^{2^l}),`$
$`{\displaystyle \underset{l=0}{\overset{1}{}}}(1a^{2^l})1,v_1(a)=v(a).`$
As is easily shown, for $`m2`$, $`G_{m1}(a)`$ is expressed by $`v_m(a)`$ as follows,
$$G_{m1}(a)=\frac{1}{2}\{(1+a^{2^{m1}})v_m(a)b(1)^ma^{2^{m1}}\}.$$
(2.32)
See Appendix A. Then the equation obtained from the equation ($`214`$)
$$G_{m1}(a_m)=(1)^{m1}a_m^{2^{m1}}\frac{1+(1)^{m1}b}{2}$$
(2.33)
is equivalent to the eq.(2$``$30). From the relation (2$``$32),
$$G_{\mathrm{}}(a)=\frac{1}{2}(v_{\mathrm{}}(a)b)$$
(2.34)
follows for $`a>1`$. Then $`G_{\mathrm{}}(a_c)=0`$ is equivalent to $`v_{\mathrm{}}(a_c)=b`$ which follows from eq.(2$``$30) immediately. Since
$$v_{\mathrm{}}^{}(a)=2v_{\mathrm{}}(a)\underset{l=0}{\overset{\mathrm{}}{}}\frac{2^la^{2^l1}}{a^{2^{l+1}}1},$$
(2.35)
we obtain
$$v_{\mathrm{}}^{}(a)>0,G_{\mathrm{}}^{}(a)>0\text{for }a>1.$$
(2.36)
Putting $`a=a_c`$ in the eq.(2$``$35), we get
$`v_{\mathrm{}}^{}(a_c)`$ $`=`$ $`2b\tau (a_c),`$ (2.38)
$`\tau (a_c){\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{2^la_c^{2^l1}}{a_c^{2^{l+1}}1}}>0.`$
Thus,
$$G_{\mathrm{}}^{}(a_c)=v_{\mathrm{}}^{}(a_c)/2=b\tau (a_c)>0.$$
(2.39)
Therefore, substituting this expression(2$``$38) into eq.(2$``$28) we obtain
$$ϵ_m=\frac{a_c^{2^m}}{\tau (a_c)}(1+h_m).$$
(2.40)
## 3 The asymmetric case
As in the symmetric case, we can discuss the period doubling cascade when $`a`$ is increased with fixed $`b`$ and $`\gamma `$ in the asymmetric case. For brevity we define $`𝒂=(a,b,\gamma )`$. $`\alpha `$ and $`\beta `$ are expressed as
$$\alpha =\frac{\gamma (1b)}{1+\gamma },\beta =\alpha +b=\frac{b+\gamma }{1+\gamma }.$$
$`\alpha ,\beta `$ and $`\gamma `$ are related by $`\alpha =\gamma (1\beta )`$. As long as $`0<b<1`$, it holds that $`0<\alpha <1`$ and $`0<\beta <1`$. Let us define $`a_La`$ and $`a_R\gamma a`$. Let $`x_M`$ be the maximum value of $`A_𝒂(x)`$, i.e., $`x_Ma\alpha `$. We define $`x^{}\frac{\gamma a}{1+\gamma a}`$. And let $`x_{0,1}`$ be the non-zero fixed point of $`A_𝒂(x)`$. The following lemma is easily proved.
lemma 1’
1. The case of $`\beta 1/2`$. For $`1<a<a_1`$, $`A_𝒂(x)`$ has the stable fixed point $`x_{0,1}=x_M`$, which satisfies $`\alpha <x_M<\beta `$. $`a_1`$ is defined by the equation $`x_1=\beta `$, that is $`a_1=\frac{\beta }{\alpha }`$.
2. The case of $`\beta <1/2`$. For $`1<a<\frac{\beta }{\alpha }`$, $`A_𝒂(x)`$ has the stable fixed point $`x_{0,1}=x_M`$, which satisfies $`\alpha <x_M<\beta `$. For $`a\frac{\beta }{\alpha }`$ the non-zero fixed point becomes $`x^{}`$ and continues to be stable until $`\frac{1\beta }{\alpha }`$. Thus, in this case we define $`a_1=\frac{1\beta }{\alpha }=\frac{1}{\gamma }`$.
In both cases, for $`a>a_1(>1)`$, $`A_𝒂(x)`$ has the unstable fixed point $`x_{0,1}=x^{}>\beta `$.
In the region $`a>a_1`$, by iterating $`A_𝒂(x)`$ twice and rescaling $`x`$, we obtain a trapezoid map with different parameter $`𝒂^{(1)},𝒂^{(1)}=(a^{(1)},b^{(1)},\gamma ^{(1)})`$. We obtain the following lemma.
lemma 2’
When $`A_𝒂(x)`$ has the non-zero unstable fixed point of $`x^{}=\frac{\gamma a}{\gamma a+1}(>\beta )`$, by rescaling the coordinate $`x`$ as $`x^{(1)}=\frac{x^{}x}{\mathrm{\Delta }x}`$, $`A_𝒂^2(x)`$ is transformed to $`A_{𝒂^{(1)}}(x^{(1)})`$ which is defined for $`x^{(1)}`$ in $`[0,1]`$, where several parameters and variables are defined as
$`\mathrm{\Delta }x`$ $`=`$ $`\gamma (a1)/(1+\gamma a),𝒂^{(1)}=𝝋(𝒂),`$
$`𝝋`$ $`:`$ $`(a,b,\gamma )(a^{(1)},b^{(1)},\gamma ^{(1)}),`$
$`a^{(1)}=a_R^2=(\gamma a)^2,a_R^{(1)}=a_La_R=\gamma a^2,`$
$`b^{(1)}`$ $`=`$ $`u(a,\gamma )b,u(a,\gamma ){\displaystyle \frac{\gamma a+1}{\gamma (a1)}},\gamma ^{(1)}=1/\gamma ,\alpha ^{(1)}={\displaystyle \frac{\alpha a\beta }{\gamma (a1)}}.`$
The proof is straightforward. Regardless of the value of $`\beta `$, $`\beta ^{(1)}>1/2`$ follows. Now, we define the $`m`$-th iteration of $`𝝋`$ for $`m0`$. That is,
$`𝒂^{(m)}`$ $``$ $`𝝋^m(𝒂)(a^{(m)},b^{(m)},\gamma ^{(m)}),`$ (3.2)
$`a^{(m)}`$ $`=`$ $`\gamma ^{2(2^m(1)^m)/3}a^{2^m},a_R^{(m)}=\gamma ^{(2^{m+1}+(1)^m)/3}a^{2^m},`$
$`b^{(m)}`$ $`=`$ $`u_m(a,\gamma )b,u_m(a,\gamma ){\displaystyle \underset{l=0}{\overset{m1}{}}}u(a^{(l)},\gamma ^{(l)}),\gamma ^{(m)}=\gamma ^{(1)^m},`$
$`\alpha ^{(m)}`$ $`=`$ $`{\displaystyle \frac{\gamma ^{(m)}(1b^{(m)})}{1+\gamma ^{(m)}}},\beta ^{(m)}={\displaystyle \frac{b^{(m)}+\gamma ^{(m)}}{1+\gamma ^{(m)}}},`$
where $`u_0(a,\gamma )1,a^{(0)}a,b^{(0)}b`$ and $`\gamma ^{(0)}\gamma `$.
Since $`a^{(m)}=(\gamma a)^{2^m}\gamma ^{2[2^{m1}+(1)^m]/3}`$ and $`\gamma a_11`$, it is easily shown that $`a^{(m)}>1`$ holds for $`a>a_1`$ and $`m0`$. Thus, for $`m1`$, $`u_m(a,\gamma )`$ and $`b^{(m)}(a,\gamma )`$ are defined for $`a>a_1`$. In the below, we assume $`a>1`$. For $`m1`$, $`b^{(m)}`$ is a continuous strictly decreasing function w.r.t. $`a`$ for $`a>a_1`$, and $`lim_a\mathrm{}b^{(m)}(a,\gamma )=b`$.
Let us prove the following lemma.
lemma 3’
For integer $`m2`$, there exists the unique value of $`a=a_m`$ greater than 1 such that
$$a^{(m1)}=\frac{\beta ^{(m1)}}{\alpha ^{(m1)}},$$
(3.3)
which is equivalent to
$$b^{(m)}(a_m)=1.$$
(3.4)
$`\{a_m\}_{m=1}^{\mathrm{}}`$ is the increasing sequence, $`1<a_1<a_2<\mathrm{}`$, and for $`a_m<a,b<b^{(m)}<1`$ for $`m1`$.
Proof
The following equivalence relations are easily proved if these quantities are defined.
$$b^{(m)}(a)=1x_M^{(m1)}=\beta ^{(m1)}a^{(m1)}=\frac{\beta ^{(m1)}}{\alpha ^{(m1)}}>1.$$
(3.5)
Let us consider the $`m=2`$ case. For $`a>1`$, $`b^{(1)}`$ is defined and $`a_2`$ is the solution of the equation $`a^{(1)}=\frac{\beta ^{(1)}}{\alpha ^{(1)}}`$. There is the unique solution greater than 1 of this equation, $`a_2=\frac{1}{2\alpha }(1+\sqrt{1\frac{4\alpha ^2}{\gamma }})`$. $`a_2>a_1`$ is easily proved. From the relation (3$``$4), $`b^{(2)}=1`$ at $`a=a_2`$.
Next, let us assume that $`b^{(m)}(a_m)=1`$ for $`m(2)`$ and $`a=a_m>1`$. Then for $`a>a_m`$, $`b<b^{(m)}<1`$ and the function $`\beta ^{(m)}/\alpha ^{(m)}=1+\frac{1+\gamma ^{(m)}}{\gamma ^{(m)}}\frac{b^{(m)}}{1b^{(m)}}`$ is continuous and decreasing w.r.t. $`a`$. On the other hand, $`a^{(m)}`$ is the continuous increasing function w.r.t. $`a`$. Further, we obtain following limits.
$`\underset{aa_m+0}{lim}{\displaystyle \frac{\beta ^{(m)}}{\alpha ^{(m)}}}`$ $`=`$ $`\mathrm{},\underset{a\mathrm{}}{lim}{\displaystyle \frac{\beta ^{(m)}}{\alpha ^{(m)}}}=1+{\displaystyle \frac{1+\gamma ^{(m)}}{\gamma ^{(m)}}}{\displaystyle \frac{b}{1b}}=\text{ finite},`$
$`\underset{aa_m+0}{lim}a^{(m)}`$ $`=`$ $`\text{finite},\underset{a\mathrm{}}{lim}a^{(m)}=\mathrm{}.`$
Thus, there exists the unique value of $`a_{m+1}(>a_m)`$ such that $`a^{(m)}=\beta ^{(m)}/\alpha ^{(m)}>1`$. Thus, $`b^{(m+1)}(a_{m+1})=1.`$ Therefore, the first half of the lemma is proved. The second half immediately follows from the decreasing property of $`b^{(m)}`$ and the facts $`b^{(1)}(a=\beta /\alpha )=1,a_1\beta /\alpha `$ and $`1<a_1<a_2`$. Q.E.D.
For positive integer $`m`$, we define $`A_{𝝋^m(𝒂)}(x^{(m)})`$ for $`a>a_m`$ from $`A_{𝝋^{m1}(𝒂)}(x^{(m1)})`$ successively by the same procedure as in the lemma2’.
We define $`x_M^{(m)}a^{(m)}\alpha ^{(m)}`$, which is the maximum value of $`A_{𝝋^m(𝒂)}(x^{(m)})`$. Further, we define $`x^{(m)}\frac{\gamma ^{(m)}a^{(m)}}{\gamma ^{(m)}a^{(m)}+1}`$. $`x_M^{(0)}=x_M`$ and $`x^{(0)}=x^{}`$.
lemma 4’
For non-negative integer $`m`$ and for $`a_m<a`$ there exists unique non-zero fixed point for $`A_{𝝋^m(𝒂)}(x^{(m)})`$. For $`a_m<a<a_{m+1}`$, the fixed point is stable. For $`m1`$, the stable fixed point is $`x_M^{(m)}`$ and $`\alpha ^{(m)}<x_M^{(m)}<\beta ^{(m)}`$.<sup>)</sup><sup>)</sup>)For $`m=0`$, see lemma 1’ At $`a=a_{m+1}`$, $`x_M^{(m)}=x^{(m)}=\beta ^{(m)}`$. For $`a_{m+1}<a`$, the non-zero unstable fixed point is $`x^{(m)}`$ and unstable and $`x^{(m)}>\beta ^{(m)}`$. Here, we define $`a_01,\alpha ^{(0)}\alpha ,\beta ^{(0)}\beta `$.
Proof
First of all, we notice that $`A_{𝝋^m(𝒂)}(x^{(m)})`$ is really a trapezoid map for $`m1`$ because from lemma 3’ $`b<b^{(m)}<1`$ for $`a_m<a`$ and $`m1`$.
For the case $`m=0`$, the statement follows from lemma 1’.
Let us assume that the statement holds for $`m0`$. Since $`x^{(m)}=\frac{\gamma ^{(m)}a^{(m)}}{\gamma ^{(m)}a^{(m)}+1}`$ is the unstable fixed point of $`A_{𝒂^{(m)}}(x^{(m)})`$ for $`a_{m+1}<a`$, from lemma 2’, taking $`A_{𝒂^{(m)}}^2(x^{(m)})`$ and rescaling $`x^{(m)}`$ as $`x^{(m+1)}=\frac{x^{(m)}x^{(m)}}{\mathrm{\Delta }x^{(m)}}`$, we obtain $`A_{𝒂^{(m+1)}}(x^{(m+1)})`$ which is defined in $`[0,1],`$ where $`\mathrm{\Delta }x^{(m)}=\gamma ^{(m)}(a^{(m)}1)/(1+\gamma ^{(m)}a^{(m)})`$. From the relation (2$``$14) it follows that $`x_M^{(m)}=\beta ^{(m)}`$ at $`a=a_{m+1}`$. Let us consider $`x_M^{(m+1)}=a^{(m+1)}\alpha ^{(m+1)}`$. Since $`x_M^{(m+1)}`$ is the continuous increasing function w.r.t. $`a`$, we obtain $`\alpha ^{(m+1)}<x_M^{(m+1)}<\beta ^{(m+1)}`$ for $`a_{m+1}<a<a_{m+2}`$. This implies $`x_M^{(m+1)}`$ is the unique fixed point for $`A_{𝝋^{m+1}(𝒂)}(x^{(m+1)})`$ and stable. It is easily shown that for $`a_{m+2}<a`$, $`x_M^{(m+1)}`$ is no more fixed point but $`x^{(m+1)}=\frac{\gamma ^{(m+1)}a^{(m+1)}}{\gamma ^{(m+1)}a^{(m+1)}+1}`$ becomes unstable fixed point. Further, it is easily shown that $`x^{(m+1)}>\beta ^{(m+1)}`$ for $`a_{m+1}<a`$, and $`x_M^{(m+1)}=x^{(m+1)}`$ at $`a=a_{m+2}`$. Q.E.D.
lemma 5’
The unique non-zero fixed point of $`A_{𝝋^m(𝒂)}(x^{(m)})`$ in lemma 4’ is the periodic cycle with prime period $`2^m`$.
The proof is the same as that in the symmetric case. Thus, we obtain the following theorem.
Theorem 2
$`x=0`$ is period 1 solution of $`A_𝒂(x)`$ for any $`a>0`$. Except for this, for any non-negative integer $`m`$, for $`a_m<a<a_{m+1}`$ there is the unique cycle with prime period $`2^k`$ for $`k=0,1,\mathrm{},m`$ and no other periodic points exist. The periodic solution with prime period $`2^m`$ is stable for $`a_m<a<a_{m+1}`$ and unstable for $`a_{m+1}<a`$. At $`a=a_{m+1}`$, the period $`2^m`$ cycle and the period $`2^{m+1}`$ cycle coincide.
The uniqueness of the cycle with prime period $`2^k`$ $`k=0,1,2,\mathrm{},m`$ for $`a_m<a<a_{m+1}`$ is proved by the same argument as that in the symmetric case. Thus, likewise the symmetric case, we conclude that as $`a`$ is increased the period doubling bifurcation cascades, and at $`a=a_m`$ the periodic solution with prime period $`2^m`$ appears.
Now, let us investigate the upper bound of the sequence $`\{a_m\}`$. For $`m1`$ $`x_M^{(m)}(a)`$ is defined at least for $`a_1<a`$. These are the strictly increasing functions w.r.t $`a`$. For any $`m1`$, $`x_M^{(m)}(a_m)=0`$ and $`lim_a\mathrm{}x_M^{(m)}(a)=\mathrm{}`$ hold. Therefore, for $`m1`$, we define $`a_M(m)`$ as the unique solution of the equation
$$x_M^{(m)}(a)=1.$$
(3.6)
For $`m=0`$, we define $`a_M(0)=\frac{1}{\alpha }`$ and $`a_Ma_M(0)`$. Then, we obtain the following lemma.
lemma 6’
$`1<a_1<a_2<\mathrm{}<a_m<a_{m+1}<\mathrm{}`$ (3.7)
$`<a_M(m+1)<a_M(m)\mathrm{}<a_M(2)<a_M(1)<a_M`$
Proof
For $`m1`$, since $`x_M^{(m)}(a_m)=0,a_m<a_M(m)`$ follows. For $`m0`$, $`x_M^{(m+1)}`$ is rewritten as
$$x_M^{(m+1)}=\gamma ^{(m)}\{a^{(m)}\}^2\frac{1}{a^{(m)}1}(x_M^{(m)}\beta ^{(m)}).$$
(3.8)
Then, we obtain $`x_M^{(m+1)}(a_M(m))=\frac{a^{(m)}}{a^{(m)}1}>1`$ for $`m0`$. Therefore, $`a_M(m+1)<a_M(m)`$ for $`m0`$. Thus, we obtain $`a_m<a_{m+1}<a_M(m+1)<a_M(m).`$ Q.E.D.
From this,
$$a_c\underset{m\mathrm{}}{lim}a_m\underset{m\mathrm{}}{lim}a_M(m)<a_M$$
(3.9)
follows. Therefore $`a_c`$ is finite.
Now, let us investigate the symbolic sequence. The sequence of $`R`$ and $`L`$ for the period $`2^m`$ solution becomes Metropolis-Stein-Stein(MSS) sequence. Let us prove this.
As in the case of the symmetric map, for $`m0`$ we define the onset of the $`2^n`$-cycle for $`A_{𝝋^m(𝒂)}(x^{(m)})`$ as $`a_n^{(m)}`$ and the $`i`$-th orbit of the period $`2^n`$ cycle as $`x_{n,i}^{(m)}`$. $`a_n^{(0)}`$ is equal to previously defined $`a_n`$. For $`m=0`$, we often omit the superscript $`(0)`$. As $`x_{m,1}`$ we set the largest $`x`$ value among $`2^m`$ members of the periodic cycle. Then, $`x_{m,1}=x_M=a\alpha `$ when the cycle is stable. Further, the coordinate space $`[0,1]`$ of $`x^{(m)}`$ into the three intervals as $`I_L^{(m)}[0,\alpha ^{(m)}],I_C^{(m)}(\alpha ^{(m)},\beta ^{(m)})`$ and $`I_R^{(m)}[\beta ^{(m)},1].`$
Then, we obtain the corresponding lemma 7’, 8’ and 9’ for $`A_𝒂(x)`$ to the lemma 7, 8, and 9 for $`T_𝒂(x)`$, respectively. We omit the statements and proofs of these lemmas, since they are almost the same as those for $`T_𝒂(x)`$. Thus, we obtain the MSS sequence for the period $`2^m`$ solution.
Now, as in the symmetric case, we derive the equations for $`a_m`$ and $`a_c`$, and obtain the asymptotic expression for $`\delta _m`$. Defining $`s(x)`$ as before, $`f_s(x)`$ for $`xI_LI_R`$ is defined as
$$f_s(x)=a(\eta +\nu x),\eta =\gamma s(x),\nu =1(1+\gamma )s(x).$$
(3.10)
In the below, $`s_i,\eta _i`$ and $`\nu _i`$ are values evaluated at $`x=x_{m,i}`$. Starting from $`x_{m,1}=\alpha a`$, we obtain $`x_{m,2^m}`$ for $`a_maa_{m+1}`$,
$`x_{m,2^m}`$ $`=`$ $`f_{s_{2^m1}}(x_{m,2^m1}){\displaystyle \underset{l=0}{\overset{2^m1}{}}}c_l^{(m)}a^{2^ml}F_m(a),`$ (3.11)
$`c_l^{(m)}`$ $`=`$ $`\eta _l{\displaystyle \underset{j=l+1}{\overset{2^m1}{}}}\nu _j\mathrm{for}\mathrm{\hspace{0.33em}\hspace{0.33em}1}l2^m2,`$
$`c_0^{(m)}`$ $`=`$ $`\alpha {\displaystyle \underset{j=1}{\overset{2^m1}{}}}\nu _j=\eta _0{\displaystyle \underset{j=1}{\overset{2^m1}{}}}\nu _j,`$
$`c_{2^m1}^{(m)}`$ $`=`$ $`\eta _{2^m1},`$
where $`\eta _0\alpha `$. Thus,
$`F_m(a_m)`$ $`=`$ $`{\displaystyle \frac{\alpha +\beta +(1)^{m1}b}{2}},`$ (3.12)
$`F_{m1}(a_m)`$ $`=`$ $`{\displaystyle \frac{\alpha +\beta +(1)^{m1}b}{2}}.`$ (3.13)
We define $`\zeta _l`$ and $`\widehat{\zeta }_l`$ as the numbers of 1 in $`s_1,s_2,\mathrm{},s_{2^l1}`$ and in $`s_1,s_2,\mathrm{},s_{2^l}`$, respectively. Then, their expressions are
$`\widehat{\zeta }_l=(2^{l+1}+(1)^l)/3\text{for}l0,`$ (3.14)
$`\zeta _l=\widehat{\zeta }_l{\displaystyle \frac{1+(1)^l}{2}}(l0),`$
$`\zeta _l=\zeta _{l1}+\widehat{\zeta }_{l1}\mathrm{for}l1,`$
where $`\zeta _00`$. The following relations hold,
$`{\displaystyle \underset{l=0}{\overset{m2}{}}}\widehat{\zeta }_l=\zeta _{m1},`$ (3.15)
$`\zeta _{2n+1}={\displaystyle \frac{4^{n+1}1}{3}}=\widehat{\zeta }_{2n+1}(n0)\text{: odd number},`$
$`\zeta _{2n}={\displaystyle \frac{2(4^n1)}{3}}=\widehat{\zeta }_{2n}1(n1)\text{: even number},`$
$`{\displaystyle \underset{j=1}{\overset{2^m1}{}}}\nu _j=(\gamma )^{\zeta _m}\text{ for }m1.`$
From these it follows that $`\widehat{\zeta }_l`$ is odd number for $`l1`$. Then, from the last relation in eq.(3$``$15), we obtain for $`l0`$
$$r_l(\gamma )^{\zeta _m}c_l^{(m)}=\{\underset{j=1}{\overset{2^m1}{}}\nu _j\}^1\eta _l\underset{j=l+1}{\overset{2^m1}{}}\nu _j=\eta _l\{\underset{j=1}{\overset{l}{}}\nu _j\}^1.$$
That is, $`r_l`$ is $`m`$-independent as long as it is defined. Thus, we get
$`r_l`$ $`=`$ $`\eta _l\{{\displaystyle \underset{j=1}{\overset{l}{}}}\nu _j\}^1\text{ for any }l(0),`$ (3.16)
$`r_0`$ $`=`$ $`\alpha ,r_{2^m}=s_{2^m}\gamma ^{1\widehat{\zeta }_m}(m0).`$ (3.17)
Thus, we define
$$G_m(a)(\gamma )^{\zeta _m}a^{2^m}F_m(a)=\underset{l=0}{\overset{2^m1}{}}r_la^l.$$
(3.18)
Let us consider the condition for the convergence of $`G_m(a)`$. Using the relations (3$``$19) and
$$r_l=\gamma ^{\widehat{\zeta }_m}r_{2^m+l}\text{ for }1l2^m1(m1),$$
(3.19)
we obtain the following recursive relation for $`m1`$,
$$G_{m+1}(a)=(1\stackrel{~}{a}_m)G_m(a)+\stackrel{~}{a}_m(\alpha \gamma s_{2^m}),$$
(3.20)
and then we obtain for $`m1`$
$$G_m(a)=\alpha \underset{l=0}{\overset{m1}{}}(1\stackrel{~}{a_l})+\underset{k=0}{\overset{m1}{}}[\underset{l=k+1}{\overset{m1}{}}(1\stackrel{~}{a_l})]\stackrel{~}{a_k}(\alpha \gamma s_{2^k}),$$
(3.21)
where $`\stackrel{~}{a_l}=a^{2^l}\gamma ^{\widehat{\zeta _l}}=\frac{1}{\sqrt{a^{(l1)}}}`$. Since $`\stackrel{~}{a_l}=(a\gamma ^{2/3})^{2^l}\gamma ^{(1)^l/3}`$, putting $`\widehat{\gamma }=`$ max $`(\gamma ,\gamma ^1)`$, we get
$$\stackrel{~}{a_l}(a\gamma ^{2/3})^{2^l}\widehat{\gamma }^{1/3}.$$
(3.22)
Thus, if $`a\gamma ^{2/3}>1,lim_m\mathrm{}G_m(a)`$ converges. That is, $`G_{\mathrm{}}(z)lim_m\mathrm{}G_m(z)`$ is the analytic function for $`|z|>\gamma ^{2/3}`$. <sup>)</sup><sup>)</sup>)In the previous paper, as the condition for the convergence of $`G_{\mathrm{}}(z)`$ we gave $`|z|>\text{max }(1,\gamma ^1)`$. This is a sufficient condition. From eq.($`312`$) we obtain
$$G_m(a_m)=(\gamma )^{\zeta _m}a_m^{2^m}\frac{\alpha +\beta +(1)^{m1}b}{2},$$
(3.23)
and then
$$G_{\mathrm{}}(a_c)=0.$$
(3.24)
Similar to the symmetric case, we obtain from eq.(3$``$23)
$`ϵ_m`$ $`=`$ $`\{G_m(a_c)\mathrm{\Gamma }a_c^{2^m}\}(1+\overline{h}_m)/G_{\mathrm{}}^{}(a_c),`$
$`\mathrm{\Gamma }=(1)^{\zeta _m}{\displaystyle \frac{\alpha +\beta +(1)^{m1}b}{2}},\overline{h}_m={\displaystyle \frac{G_{\mathrm{}}^{}(a_c)}{G_m^{}(\overline{a}_m)+\mathrm{\Gamma }2^m(\widehat{a}_m)^{2^m1}}}1,`$
where $`a_m<\overline{a}_m<a_c,a_m<\widehat{a}_m<a_c`$ and $`lim_m\mathrm{}\overline{h}_m=0`$. Further, using the relation (3$``$17) we get
$`G_m(a_c)`$ $`=`$ $`a_c^{2^m}\gamma ^{\widehat{\zeta }_m}(\alpha \gamma s_{2^m})+(a_c^{2^m}\gamma ^{\widehat{\zeta }_m})^2q_m,`$
$`q_m={\displaystyle \frac{1}{1a_c^{2^m}\gamma ^{\widehat{\zeta }_m}}}(\alpha \gamma s_{2^m}+\gamma ^{2\widehat{\zeta }_m}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}r_{2^{m+1}+l}a_c^l).`$
Substituting (3$``$26) into (3$``$25) , we obtain
$`ϵ_m`$ $`=`$ $`a_ca_m={\displaystyle \frac{ba_c^{2^m}\gamma ^{\zeta _m}}{G_{\mathrm{}}^{}(a_c)}}(1+h_m),`$ (3.28)
$`h_m=\overline{h}_m(1+a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}{\displaystyle \frac{q_m}{b}})+a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}{\displaystyle \frac{q_m}{b}}.`$
$`lim_m\mathrm{}h_m=0`$ is shown in appendix B. Thus,
$`\delta _m`$ $`=`$ $`\gamma ^{(1)^m/3}(a_c\gamma ^{2/3})^{2^m}(1+l_m),`$
$`l_m=[h_mh_{m+1}a_c^{2^m}\gamma ^{\widehat{\zeta }_m}(1+h_{m+1})+a_c^{2^{m+1}}\gamma ^{\widehat{\zeta }_{m+1}}(1+h_{m+2})]`$
$`/[1+h_{m+1}a_c^{2^{m+1}}\gamma ^{\widehat{\zeta }_{m+1}}(1+h_{m+2})],`$
and $`lim_m\mathrm{}l_m=0`$.
To obtain a similar relation to the relation (2$``$32) in the case of symmetric map, we define $`v_m(a,\gamma )`$ as
$$v_m(a,\gamma )=1/u_m(a,\gamma )=\underset{l=0}{\overset{m1}{}}v(a^{(l)},\gamma ^{(l)}),v(a,\gamma )=\frac{1}{u(a,\gamma )}=\frac{\gamma (a1)}{\gamma a+1}.$$
(3.30)
For $`m1`$, $`v_m(a)`$ is rewritten as
$`v_m(a)`$ $`=`$ $`{\displaystyle \frac{(1a^1)}{g_{m1}(a)}}{\displaystyle \underset{l=0}{\overset{m2}{}}}h_l(a),`$
$`g_l(a)1+\gamma ^{\widehat{\zeta }_l}a^{2^l},h_l(a)1\gamma ^{\widehat{\zeta }_l}a^{2^l},\text{ for }l0,`$
$`{\displaystyle \underset{l=0}{\overset{1}{}}}h_l(a)1,v_1(a)=v(a).`$
Then, for $`m2`$, the following equation is proved by using the above relations
$$G_{m1}(a)=\alpha +\frac{1}{1+\gamma ^1}\{g_{m1}(a)v_m(a)1(1)^ma^{2^{m1}}\gamma ^{\zeta _{m1}}\}.$$
(3.32)
See Appendix A. Thus, the condition for $`a_m`$, $`v_m(a_m)=b`$, is equivalent to the following equation derived from eq.(3$``$13),
$$G_{m1}(a_m)=(\gamma )^{\zeta _{m1}}a_m^{2^{m1}}\frac{\alpha +\beta +(1)^{m1}b}{2}.$$
(3.33)
From relation (3$``$32), we get for $`a>\gamma ^{2/3}`$
$$G_{\mathrm{}}(a)=\alpha +\frac{1}{1+\gamma ^1}(v_{\mathrm{}}(a)1).$$
(3.34)
Thus,
$`G_{\mathrm{}}^{}(a)`$ $`=`$ $`{\displaystyle \frac{1}{1+\gamma ^1}}v_{\mathrm{}}^{}(a)={\displaystyle \frac{2b}{1+\gamma ^1}}\tau (a),`$
$`\tau (a){\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1+\gamma ^{(l)}}{2(a^{(l)}1)(\gamma ^{(l)}a^{(l)}+1)}}\gamma ^{2(2^l(1)^l)/3}2^la^{2^l1}>0.`$
Therefore, we obtain
$$ϵ_m=\frac{1}{2\tau (a_c)}a_c^{2^m}\gamma ^{\zeta _m}(1+\gamma ^1)(1+h_m).$$
(3.36)
## 4 Period doubling bifurcation of the period $`p`$ solution : symmetric case
In this section, we consider the period doubling bifurcation of the periodic solution with the prime period $`p(3)`$ which emerges by a tangent bifurcation.
Let us consider the mapping $`T_𝒂^p(x)`$. At $`a=a_M=\frac{2}{1b}`$, the map looks like Fig.5.
As $`a`$ is decreased from $`a=a_M`$, at some value of $`a=a_{p,0}`$, the map becomes tangent to the line $`y=x`$, as is illustrated in Fig.6. If we put $`x_1=a\frac{1b}{2}`$, at this point the symbols for $`x_1,x_2,\mathrm{},x_{p1}`$ are $`RL\mathrm{}L`$. Then,
$$x_p=a^{p1}(1\frac{a(1b)}{2}).$$
(4.1)
As is shown in Appendix C, $`x_p(a)`$ takes the maximum value at $`a=a_{p,max}\frac{p1}{p}a_M`$ and $`\frac{1+b}{2}<x_p(a_{p,max})`$ for $`p3`$. Thus, the point $`a_{p,0}`$ where the tangent bifurcation takes place satisfies the condition
$$x_p(a_{p,0})=\frac{1+b}{2}.$$
(4.2)
This equation has two solutions, one is $`a=1`$ and the other corresponds to $`a_{p,0}`$. Thus, we obtain
$$1<\frac{p1}{p}a_M<a_{p,0}<a_{p,1},$$
where $`a_{p,1}`$ is defined by
$$x_p(a_{p,1})=\frac{1b}{2},$$
(4.3)
and $`a_{p,1}>1`$. See Fig.7.
$`a_{p,1}`$ is the onset point of the period $`2p`$ point. The symbolic sequence for $`x_1,x_2,\mathrm{},x_{p1}`$ is $`RL\mathrm{}L`$ for $`a_{p,0}<a`$. See Appendix D. For $`a_{p,0}<a<a_{p,1}`$, let us consider the map $`T_𝒂^p(x)`$. The unstable periodic point with period $`p`$ of $`T_𝒂(x)`$ is
$$x^{}=\frac{a^{p1}(a1)}{a^p1}.$$
(4.4)
Rescaling the map $`T_𝒂^p(x)`$ around $`x=1/2`$ for $`a>a_{p,0}`$, we obtain the trapezoid map $`T_{𝒂^{(0)}}(x^{(0)})`$, where
$`𝒂^{(0)}`$ $`=`$ $`(a^{(0)},b^{(0)}),a^{(0)}=a^p,b^{(0)}={\displaystyle \frac{b}{\mathrm{\Delta }x}}=r(a)b,\mathrm{\Delta }x=2x^{}1,`$ (4.5)
$`r(a)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Delta }x}}={\displaystyle \frac{a^p1}{a^p2a^{p1}+1}},x^{(0)}={\displaystyle \frac{x^{}x}{2x^{}1}}.`$
For $`a>1`$, $`r^{}(a)<0`$. As before, we define $`x_M^{(0)}=a^{(0)}\frac{1b^{(0)}}{2}`$. For $`a_{p,0}<a<a_{p,1}`$, $`x_M^{(0)}`$ corresponds to $`x_p`$, that is $`x_M^{(0)}=\frac{x^{}x_p}{2x^{}1}`$. Therefore, at $`a=a_{p,0},x_M^{(0)}=0`$ and at $`a=a_{p,1},x_M^{(0)}=\frac{1+b^{(0)}}{2}`$. We note that $`b^{(0)}(a_{p,0})=1.`$ Since $`r(a)`$ is the strictly decreasing continuous function w.r.t. $`a`$ for $`a>1`$, $`b^{(0)}(a)`$ is the strictly decreasing continuous function and $`x_M^{(0)}(a)`$ is the strictly increasing continuous function w.r.t. $`a`$. Therefore, for $`a_{p,0}<a<a_{p,1},\frac{1b^{(0)}}{2}<x_M^{(0)}<\frac{1+b^{(0)}}{2}`$ , and for $`a_{p,1}<a,\frac{1+b^{(0)}}{2}<x_M^{(0)}`$. That is, $`x_M^{(0)}`$ is the stable fixed point of $`T_{𝒂^{(0)}}(x^{(0)})`$ for $`a_{p,0}<a<a_{p,1}`$, and for $`a_{p,1}<a`$ the fixed point is $`x_u^{}\frac{a^{(0)}}{a^{(0)}+1}`$ and is unstable. For $`a_{p,1}<a`$, rescaling the map $`T_{𝒂^{(0)}}^2(x^{(0)})`$ we obtain the trapezoid map $`T_{𝒂^{(1)}}(x^{(1)})`$, where,
$`𝒂^{(1)}`$ $``$ $`𝝋(𝒂^{(0)})=(a^{(1)},b^{(1)}),a^{(1)}=\{a^{(0)}\}^2=a^{2p},`$ (4.6)
$`b^{(1)}`$ $`=`$ $`u(a^{(0)})b^{(0)}=u(a^p)r(a)b,u(a)={\displaystyle \frac{a+1}{a1}}.`$
As before, we define $`𝒂^{(m)}`$ as
$`𝒂^{(m)}`$ $``$ $`𝝋^m(𝒂^{(0)})=(a^{(m)},b^{(m)}),a^{(m)}=\{a^{(m1)}\}^2=a^{2^mp},`$ (4.7)
$`b^{(m)}`$ $`=`$ $`u(a^{(m1)})b^{(m1)}=u_m(a)r(a)b,u_m(a)={\displaystyle \underset{l=0}{\overset{m1}{}}}u(a^{(l)}).`$
Since $`b^{(m)}(a)`$ is the decreasing function w.r.t. $`a`$, we can repeat the same argument as in section 2. At $`a=a_{p,m}`$, the period $`2^m`$ solution of $`T_{𝒂^{(0)}}(x^{(0)})`$ appears and $`a_{p,m}`$ satisfies the condition
$$b^{(m)}(a_{p,m})=1,\text{for }m0,$$
(4.8)
and
$$1<a_{p,0}<a_{p,1}<a_{p,2}<\mathrm{}<\frac{2}{1b}=a_M.$$
(4.9)
The symbolic sequence for the period $`2^m`$ solution is the MSS sequence for the map $`T_{𝒂^{(0)}}(x^{(0)})`$. For $`a_{p,m}aa_{p,m+1}`$, starting from $`x_1^{(0)}=x_M^{(0)}`$, $`x_{2^m}^{(0)}`$ is expressed by the function $`F_m(a^{(0)},b^{(0)})`$ . Here, $`F_m(a,b)`$ is defined by eq.($`210`$) and we include $`b`$ as an independent variable to $`F_m`$ explicitly. Thus, we obtain the other condition for $`a_{p,m}`$,
$$F_m(a^{(0)}(a_{p,m}),b^{(0)}(a_{p,m}))=\frac{1+(1)^{m1}b^{(0)}(a_{p,m})}{2}.$$
(4.10)
In terms of $`G_m(a,b)`$ defined by eq.(2$``$19), the equation (2$``$21) becomes
$$G_m(a^{(0)}(a_{p,m}),b^{(0)}(a_{p,m}))=a^{(0)}(a_{p,m})^{2^m}\frac{(1)^mb^{(0)}(a_{p,m})}{2},$$
(4.11)
and we obtain
$$G_{\mathrm{}}(a^{(0)}(a_{p,c}),b^{(0)}(a_{p,c}))=0,$$
(4.12)
where $`a_{p,c}=lim_m\mathrm{}a_{p,m}`$. From these equations, we get for $`ϵ_{p,m}=a_{p,c}a_{p,m}`$,
$$ϵ_{p,m}\frac{a_{p,c}^{p2^m}r(a_{p,c})b}{\frac{G_{\mathrm{}}}{a^{(0)}}pa_{p,c}^{p1}+\frac{1}{2}|r^{}(a_{p,c})|b}.$$
(4.13)
Since $`\frac{G_{\mathrm{}}(a,b)}{a}>0`$, the denominator is positive. Therefore, we obtain
$$\delta _ma_{p,c}^{p2^m}.$$
(4.14)
As for the relations among $`a_{p,0}`$ s, the following ordering holds,
$$a_ca_{3,0}<a_{4,0}<\mathrm{}<a_{p,0}<a_{p+1,0}<\mathrm{}<a_M.$$
(4.15)
See Appendix E.
## 5 Period doubling bifurcation of the period $`p3`$ solution : asymmetric case
As in the previous case, we investigate the period doubling bifurcation of the period $`p(3)`$ solution in the asymmetric case.
In this case also, we consider the period $`p`$ solution starting with $`x_1=a\alpha `$ with the sequence of symbols, $`RL\mathrm{}L`$. Then, we obtain
$$x_p=\gamma a^{p1}(1a\alpha ).$$
(5.1)
As is shown in Appendix C, $`x_p(a)`$ takes the maximum value at $`a=a_{p,max}\frac{p1}{p}a_M`$ and $`\beta <x_p(a_{p,max})`$ for $`p3`$, where $`a_M=1/\alpha `$ as before. Then, the onset point $`a_{p,0}(>a_{p,max})`$ of the period $`p`$ solution satisfies
$$x_p(a_{p,0})=\beta .$$
(5.2)
Although this equation has two solutions, we should adopt the larger one as in the symmetric case. The symbolic sequence for $`x_1,x_2,\mathrm{},x_{p1}`$ is $`RL\mathrm{}L`$ for $`a_{p,0}a`$. See Appendix D. On the other hand, the onset point of the period $`2p`$ solution, $`a_{p,1}`$, satisfies
$$x_p(a_{p,1})=\alpha .$$
(5.3)
Then, for $`a_{p,0}<a`$, the unstable periodic point $`x^{}`$ with period $`p`$ for $`A_𝒂(x)`$ is
$$x^{}=\frac{(\gamma a1)\gamma a^{p1}}{\gamma ^2a^p1}.$$
(5.4)
Similar to the symmetric case, we obtain an asymmetric trapezoid map $`A_{𝒂^{(0)}}(x^{(0)})`$ by rescaling the map $`A_𝒂^p(x)`$ in the vicinity of $`I_C`$ for $`a>a_{p,0}`$, where
$`𝒂^{(0)}`$ $`=`$ $`(a^{(0)},b^{(0)},\gamma ^{(0)}),a^{(0)}=\gamma ^2a^p,\gamma ^{(0)}=\gamma ^1,b^{(0)}={\displaystyle \frac{b}{\mathrm{\Delta }x}}=r(a)b,`$ (5.5)
$`\mathrm{\Delta }x`$ $`=`$ $`{\displaystyle \frac{\gamma g(a)}{\gamma ^2a^p1}},g(a)=\gamma a^p(1+\gamma )a^{p1}+1,`$
$`r(a)`$ $`=`$ $`1/\mathrm{\Delta }x={\displaystyle \frac{\gamma ^2a^p1}{\gamma g(a)}},x^{(0)}={\displaystyle \frac{x^{}x}{\mathrm{\Delta }x}},`$
$`\alpha ^{(0)}`$ $``$ $`{\displaystyle \frac{x^{}\beta }{\mathrm{\Delta }x}}={\displaystyle \frac{\gamma ^2a^p(1\beta )\gamma a^{p1}+\beta }{\gamma g(a)}},`$
$`\beta ^{(0)}`$ $``$ $`{\displaystyle \frac{x^{}\alpha }{\mathrm{\Delta }x}}={\displaystyle \frac{\gamma ^2a^p(1\alpha )\gamma a^{p1}+\alpha }{\gamma g(a)}}.`$
$`r^{}(a)`$ becomes
$$r^{}(a)=\frac{(1+\gamma )a^{p2}f(a)}{\gamma g(a)^2},$$
(5.6)
where $`f(a)=\gamma ^2a^pp\gamma a+p1`$. We can prove that at least for $`a_{p,0}a`$, $`f(a)>0`$. Thus, for $`a_{p,0}a`$, $`r^{}(a)<0`$. Further, $`\alpha ^{(0)^{}}(a)=\frac{ba^{p2}f(a)}{\gamma g(a)^2}>0`$ and $`\beta ^{(0)^{}}(a)=\gamma \alpha ^{(0)^{}}(a)<0`$ for $`aa_{p,0}`$. Thus, $`b^{(0)}`$ is the strictly decreasing continuous function and $`x_M^{(0)}=a^{(0)}\alpha ^{(0)}`$ is the strictly increasing continuous function w.r.t. $`a`$. At $`a=a_{p,0}`$, $`\alpha ^{(0)}=0,\beta ^{(0)}=b^{(0)}=1`$. For $`a_{p,0}aa_{p,1}`$ $`x_p(a)`$ corresponds to $`x_M^{(0)}`$, that is, $`x_M^{(0)}=\frac{x^{}x_p}{\mathrm{\Delta }x}`$. Then, at $`a=a_{p,0},x_M^{(0)}=0`$ and at $`a=a_{p,1},x_M^{(0)}=\beta ^{(0)}`$. Thus, $`\alpha ^{(0)}<x_M^{(0)}<\beta ^{(0)}`$ for $`a_{p,0}<a<a_{p,1}`$, and $`x_M^{(0)}`$ is the stable fixed point for $`A_{𝒂^{(0)}}(x^{(0)})`$. Further, for $`a_{p,1}<a`$, $`x_M^{(0)}>\beta ^{(0)}`$ and the fixed point becomes $`x_u^{}\frac{\gamma ^{(0)}a^{(0)}}{1+\gamma ^{(0)}a^{(0)}}`$ and unstable. For $`a_{p,1}<a`$, rescaling the map $`A_{𝒂^{(0)}}^2(x^{(0)})`$, we obtain $`A_{𝒂^{(1)}}(x^{(1)})`$, where
$`𝒂^{(1)}`$ $``$ $`𝝋(𝒂^{(0)})=(a^{(1)},b^{(1)},\gamma ^{(1)}),a^{(1)}=\{\gamma ^{(0)}a^{(0)}\}^2=\gamma ^2a^{2p},`$ (5.7)
$`b^{(1)}`$ $`=`$ $`u(a^{(0)},\gamma ^{(0)})b^{(0)}=u(a^p)r(a)b,\gamma ^{(1)}=\{\gamma ^{(0)}\}^1,u(a,\gamma )={\displaystyle \frac{\gamma a+1}{\gamma (a1)}}.`$
Thus, $`a^{(1)}`$ is increasing, $`b^{(1)}`$ is decreasing and $`\gamma ^{(1)}`$ is constant w.r.t. $`a`$. Further, we define $`𝒂^{(m)}`$ as
$`𝒂^{(m)}`$ $``$ $`𝝋^m(𝒂^{(0)})=(a^{(m)},b^{(m)},\gamma ^{(m)}),a^{(m)}=\{\gamma ^{(m1)}a^{(m1)}\}^2,`$ (5.8)
$`b^{(m)}`$ $`=`$ $`u(a^{(m1)},\gamma ^{(m1)})b^{(m1)}=u_m(a,\gamma )r(a)b,u_m(a,\gamma )={\displaystyle \underset{l=0}{\overset{m1}{}}}u(a^{(l)},\gamma ^{(l)}),`$
$`\gamma ^{(m)}`$ $`=`$ $`\{\gamma ^{(m1)}\}^1.`$
Since $`b^{(m)}`$ is the decreasing function w.r.t. $`a`$, we can make a quite similar argument to that in §4. In particular, at $`a=a_{p,m}`$, the period doubling bifurcation takes place and
$$b^{(m)}(a_{p,m})=1\text{for }m0,$$
(5.9)
and
$$1<a_{p,0}<a_{p,1}<a_{p,2}<\mathrm{}<\frac{1}{\alpha }=a_M.$$
(5.10)
The symbolic sequence for the period $`2^m`$ solution is the MSS sequence for the map $`A_{𝒂^{(0)}}(x^{(0)})`$. For $`a_{p,m}aa_{p,m+1}`$, starting from $`x_1^{(0)}=x_M^{(0)}`$, $`x_{2^m}^{(0)}`$ is expressed by the function $`F_m(a^{(0)},b^{(0)},\gamma ^{(0)})`$ and we get
$$F_m(a^{(0)}(a_{p,m}),b^{(0)}(a_{p,m}),\gamma ^{(0)})=\frac{\alpha ^{(0)}(a_{p,m})+\beta ^{(0)}(a_{p,m})+(1)^{m1}b^{(0)}(a_{p,m})}{2},$$
(5.11)
where $`F_m(a,b,\gamma )`$ is defined by eq.(3$``$11). In terms of $`G_m(a,b,\gamma )`$ defined by eq.(3$``$18), the equation (5$``$11) becomes
$`G_m(a^{(0)}(a_{p,m}),b^{(0)}(a_{p,m}),\gamma ^{(0)})=`$ (5.12)
$`(\gamma ^{(0)})^{\zeta _m}\{a_m^{(0)}(a_{p,m})\}^{2^m}{\displaystyle \frac{\alpha ^{(0)}(a_{p,m})+\beta ^{(0)}(a_{p,m})+(1)^{m1}b^{(0)}(a_{p,m})}{2}}.`$
Let us define $`a_{p,c}=lim_m\mathrm{}a_{p,m}`$. Thus, we obtain
$$G_{\mathrm{}}(a^{(0)}(a_{p,c}),b^{(0)}(a_{p,c}),\gamma ^{(0)})=0,$$
(5.13)
and $`ϵ_{p,m}=a_{p,c}a_{p,m}`$ is given by
$$ϵ_{p,m}\frac{(\gamma ^2a_{p,c}^p)^{2^m}\gamma ^{\zeta _m}r(a_{p,c})b}{\frac{G_{\mathrm{}}}{a^{(0)}}\gamma ^2pa_{p,c}^{p1}+\frac{1}{2}|r^{}(a_{p,c})|b},$$
(5.14)
and then
$$\delta _m\gamma ^{(1)^{m1}/3}(a_{p,c}^p\gamma ^{4/3})^{2^m}.$$
(5.15)
As for the relations among $`a_{p,0}`$ s, the following ordering holds,
$$a_ca_{3,0}<a_{4,0}<\mathrm{}<a_{p,0}<a_{p+1,0}<\mathrm{}<a_M.$$
(5.16)
See Appendix E.
## 6 Summary and discussion
In this paper, we studied the symmetric and the asymmetric trapezoid maps rigorously. We gave the proofs of several results about the period doubling bifurcation which occurs as the parameter $`a`$ is increased. We obtained following scaling results for the period doubling bifurcation starting from a period one solution.
1. Symmetric case
$`ϵ_m`$ $`=`$ $`{\displaystyle \frac{ba_c^{2^m}}{G_{\mathrm{}}^{}(a_c)}}(1+h_m),`$
$`\delta _m`$ $`=`$ $`a_c^{2^m}(1+l_m)`$
where $`lim_m\mathrm{}h_m=0`$ and $`lim_m\mathrm{}l_m=0`$.
2. Asymmetric case
$`ϵ_m`$ $`=`$ $`{\displaystyle \frac{ba_c^{2^m}\gamma ^{\zeta _m}}{G_{\mathrm{}}^{}(a_c)}}(1+h_m),`$
$`\delta _m`$ $`=`$ $`\gamma ^{(1)^m/3}(a_c\gamma ^{2/3})^{2^m}(1+l_m),`$
where $`lim_m\mathrm{}h_m=0`$ and $`lim_m\mathrm{}l_m=0`$.
The new results in this paper are on the period doubling bifurcation which starts from the periodic solution with any prime period $`p(3)`$.
1. Symmetric case
$`ϵ_{p,m}`$ $``$ $`{\displaystyle \frac{a_{p,c}^{p2^m}r(a_{p,c})b}{\frac{G_{\mathrm{}}}{a^{(0)}}pa_{p,c}^{p1}+\frac{1}{2}|r^{}(a_{p,c})|b}}.`$
$`\delta _m`$ $``$ $`a_{p,c}^{p2^m}.`$
2. Asymmetric case
$`ϵ_m`$ $``$ $`{\displaystyle \frac{(\gamma ^2a_{p,c}^p)^{2^m}\gamma ^{\zeta _m}r(a_{p,c})b}{\frac{G_{\mathrm{}}}{a^{(0)}}\gamma ^2pa_{p,c}^{p1}+\frac{1}{2}|r^{}(a_{p,c})|b}},`$
$`\delta _m`$ $``$ $`\gamma ^{(1)^{m1}/3}(a_{p,c}^p\gamma ^{4/3})^{2^m}.`$
These results imply that for any period doubling cascade, the accumulation rate to the accumulation point is extremely fast, in fact, it is exponential.
In a one-dimensional map with one hump, the Feigenbaum constant depends on the power $`z`$ which characterizes the behaviour of the map in the vicinity of the critical point. It has been proved that $`lim_z\mathrm{}\delta (z)=`$ finite. The result in this paper shows that $`\delta (\mathrm{})`$ is infinity. That is, the superconvergence takes place only in the case $`z=\mathrm{}`$.
This feature is considered to be attributed to the flatness of the summit. For example, in a one humped map with a flat summit, superconvergence will take place if it satisfies some conditions, e.g. the absolute value of the derivative is greater than some constant $`\lambda >1`$ outside a region which contains the flat part of the map.
In the other papers , similar results were obtained by a different method. In those studies, the authors estimated quantities in question by inequalities and obtained that $`a_m`$ is quadratically convergent, in the case of period doubling of period one solution. This implies $`lim_m\mathrm{}\frac{\mathrm{ln}\delta _m}{\mathrm{ln}\delta _{m1}}=2.`$ On the other hand, our method is constructive. In fact, we gave the precise equation for the onset point $`a_m`$ of the periodic point with period $`2^m`$ and also gave the equation for the accumulation point $`a_c`$. From these equations, we obtained the above scaling relations. We also extended the argument not only to the asymmetric trapezoid map but also period doubling of the prime period $`p(3)`$ solution which emerges from a tangent bifurcation.
In , the authors mention the problem of convergence of the sequence $`\mathrm{\Delta }ϵ_m/(\mathrm{\Delta }ϵ_{m1})^2`$, and from our result, this is $`ϵ_m/ϵ_{m1}^2G_{\mathrm{}}^{}(a_c)\gamma ^{(1+(1)^m)/2}/b`$ and does not converge for the asymmetric case.
The trapezoid maps studied here are a kind of exactly solvable models of the renomalization group. It is interesting to investigate the further detailed bifurcation structures in these models. This is a future problem.
## A Proof of relations ($`232`$) and ($`332`$)
In this appendix, we prove the relation (3$``$32). The relation (2$``$32) is obtained by putting $`\gamma =1`$.
For $`m2`$, the relation ($`332`$) is
$`G_{m1}(a)`$ $`=`$ $`\alpha +{\displaystyle \frac{1}{1+\gamma ^1}}\{g_{m1}(a)v_m(a)1(1)^ma^{2^{m1}}\gamma ^{\zeta _{m1}}\},`$
$`v_m(a)={\displaystyle \frac{(1a^1)}{g_{m1}(a)}}{\displaystyle \underset{l=0}{\overset{m2}{}}}h_l(a),`$
$`g_l(a)1+\gamma ^{\widehat{\zeta }_l}a^{2^l},h_l(a)1\gamma ^{\widehat{\zeta }_l}a^{2^l},\text{ for }l0.`$
Let us define $`\varphi _m`$ for $`m1`$ as
$$\varphi _m(a)g_m(a)v_{m+1}(a)=(1a^1)\underset{l=0}{\overset{m1}{}}h_l(a).$$
(A.2)
We expand $`\varphi _m(a)`$ as
$$\varphi _m(a)\underset{l=0}{\overset{2^m}{}}w_l^{(m)}a^l,$$
(A.3)
and then using relations (3$``$15) we obtain for $`m1`$,
$$w_0^{(m)}=1,w_{2^m}^{(m)}=(1)^{m1}\gamma ^{\zeta _m}.$$
(A.4)
In particular,
$$w_1^{(1)}=(1+\gamma ^1).$$
(A.5)
Let us derive the recursive relations for $`w_l^{(m)}`$. For $`m1`$,
$$\varphi _{m+1}(a)=g_m(a)v_{m+1}(a)=h_m(a)\varphi _m(a)=\varphi _m(a)\gamma ^{\widehat{\zeta }_m}a^{2^m}\varphi _m(a).$$
(A.6)
Comparing coefficients, we obtain for $`m1`$,
$`𝒪(a^{2^{m+1}}):w_{2^{m+1}}^{(m+1)}=\gamma ^{\widehat{\zeta }_m}w_{2^m}^{(m)},`$ (A.7)
$`𝒪(a^{2^m}):w_{2^m}^{(m+1)}=w_{2^m}^{(m)}\gamma ^{\widehat{\zeta }_m}w_0^{(m)},`$ (A.8)
$`𝒪(a^{(2^m+l)}):w_{2^m+l}^{(m+1)}=\gamma ^{\widehat{\zeta }_m}w_l^{(m)}\text{ for }0<l<2^m,`$ (A.9)
$`𝒪(a^l):w_l^{(m+1)}=w_l^{(m)}\text{ for }0l<2^m.`$ (A.10)
Using the relation (A$``$4), eq.(A$``$7) is automatically satisfied. Form eq.(A$``$8), we get
$$w_{2^m}^{(m+1)}=(1+\gamma ^1)s_{2^m}\gamma ^{1\widehat{\zeta }_m}.$$
(A.11)
Now, let us return to the quantity $`G_{m1}(a)`$. We put the r.h.s. of relation (A$``$1) as $`\widehat{G}_{m1}(a)`$, and expand it by $`a^1`$ for $`m2`$.
$$\widehat{G}_{m1}(a)=\underset{l=0}{\overset{2^{m1}1}{}}\overline{r}_l^{(m1)}a^l.$$
(A.12)
Thus, we obtain the following relations for $`m2`$,
$`\overline{r}_0^{(m1)}`$ $`=`$ $`\alpha ,`$ (A.13)
$`\overline{r}_l^{(m1)}`$ $`=`$ $`{\displaystyle \frac{w_l^{(m1)}}{1+\gamma ^1}},\text{ for }0<l2^{m1}1.`$ (A.14)
Thus, from (A$``$11), (A$``$9) and (A$``$10) we obtain for $`m1`$
$`\overline{r}_{2^m}^{(m+1)}`$ $`=`$ $`s_{2^m}\gamma ^{1\widehat{\zeta }_m},`$ (A.15)
$`\overline{r}_{l+2^m}^{(m+1)},`$ $`=`$ $`\gamma ^{\widehat{\zeta }_m}\overline{r}_l^{(m)}\text{ for }0<l2^m1,`$ (A.16)
$`\overline{r}_l^{(m+1)}`$ $`=`$ $`\overline{r}_l^{(m)}\text{ for }0l2^m1.`$ (A.17)
The relation (A$``$17) implies that $`\overline{r}_l^{(m)}`$ is $`m`$ independent as long as it is defined and so we omit the superscript $`(m)`$. From relations (A$``$15) and (A$``$16) it turns out that $`\overline{r}_l`$ for $`l=2^m+1,\mathrm{},2^{m+1}1`$ are determined by $`\overline{r}_l`$ for $`l=1,\mathrm{},2^m1`$. Thus, $`\overline{r}_1`$ and $`\overline{r}_{2^m}`$ for $`(m1)`$ determines $`\overline{r}_l`$ for $`l>0`$. From eqs.(A$``$5) and (A$``$14) $`\overline{r}_1=1=r_1`$ follows. Also, from eqs.(3$``$17) and (3$``$19) we note $`r_l`$ satisfies the same relations as (A$``$15) and (A$``$16). Then, $`\overline{r}_{2^m}=r_{2^m}`$ for $`(m1)`$. Therefore, $`r_l=\overline{r}_l`$ follows for $`l1`$. Finally, from the relation (A$``$13) we have $`\overline{r}_0=\alpha =r_0`$. Therefore, we obtain $`\widehat{G}_{m1}(a)=G_{m1}(a)`$. That is, the relation (3$``$32) holds.
## B Proof of $`lim_m\mathrm{}h_m=0`$
In this appendix, we prove
$$\underset{m\mathrm{}}{lim}h_m=0,$$
where $`h_m`$ is
$`h_m`$ $`=`$ $`\overline{h}_m(1+a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}{\displaystyle \frac{q_m}{b}})+a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}{\displaystyle \frac{q_m}{b}},`$
$`q_m={\displaystyle \frac{1}{1a_c^{2^m}\gamma ^{\widehat{\zeta }_m}}}(\alpha \gamma s_{2^m}+\gamma ^{2\widehat{\zeta }_m}{\displaystyle \underset{l=0}{\overset{\mathrm{}}{}}}r_{2^{m+1}+l}a_c^l).`$
To show this, we only have to prove that
$$\underset{m\mathrm{}}{lim}a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}q_m=0.$$
First we estimate the series $`S_{l=0}^{\mathrm{}}r_{2^{m+1}+l}a_c^l`$.
Let $`\tau _l`$ be the number of 1 in $`s_1,s_2,\mathrm{},s_l`$. Then, $`r_l`$ is
$$r_l=\gamma s_l[\underset{j=1}{\overset{l}{}}\{1(1+\gamma )s_j\}]^1=\gamma s_l(\gamma )^{\tau _l}.$$
(B.1)
Thus,
$$|r_l|\gamma ^{1\tau _l}.$$
Let $`\stackrel{~}{\tau }_l`$ be the number of 1 in $`s_{2^{m+1}+1},s_{2^{m+1}+2},\mathrm{},s_{2^{m+1}+l}`$. Then, for $`l0`$ we obtain
$$\tau _{2^{m+1}+l}=\widehat{\zeta }_{m+1}+\stackrel{~}{\tau }_l,\stackrel{~}{\tau }_ll,\stackrel{~}{\tau }_00.$$
Then,
$$|S|\underset{l=0}{\overset{\mathrm{}}{}}\gamma ^{1\widehat{\zeta }_{m+1}\stackrel{~}{\tau }_l}a_c^l.$$
Since $`a_c\gamma >1`$, we obtain
$`|S|`$ $``$ $`{\displaystyle \frac{a_c}{a_c1}}\gamma ^{1\widehat{\zeta }_{m+1}},\text{ for }\gamma 1,`$
$`|S|`$ $``$ $`{\displaystyle \frac{a_c\gamma }{a_c\gamma 1}}\gamma ^{1\widehat{\zeta }_{m+1}},\text{ for }\gamma <1.`$
Thus for any $`\gamma `$,
$$|S|\mathrm{const}.\gamma ^{\widehat{\zeta }_{m+1}}.$$
From the following relation
$$\zeta _m\widehat{\zeta }_{m+1}=\widehat{\zeta }_m\frac{1+(1)^{m1}}{2},$$
we get
$$a_c^{2^m}\gamma ^{\zeta _m}|S|\mathrm{const}.a_c^{2^m}\gamma ^{\widehat{\zeta }_m}\gamma ^{\frac{1+(1)^{m1}}{2}}.$$
Further,
$$\zeta _m2\widehat{\zeta }_m=\widehat{\zeta }_m\frac{1+(1)^m}{2},$$
and then
$$a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}=a_c^{2^m}\gamma ^{\widehat{\zeta }_m}\gamma ^{\frac{1+(1)^m}{2}}.$$
Thus, since $`a_c^{2^m}\gamma ^{\widehat{\zeta }_m}`$ tends to 0 as $`m\mathrm{}`$, we obtain
$$\underset{m\mathrm{}}{lim}a_c^{2^m}\gamma ^{\zeta _m2\widehat{\zeta }_m}q_m=0.$$
## C Proofs of $`a_{p,max}=\frac{p1}{p}a_M`$ and $`x_p(a_{p,max})>\beta `$ for $`p3`$
In this appendix, we prove that for the asymmetric case, $`x_p=\gamma a^{p1}(1a\alpha )`$ has the maximum value at $`a=a_{p,max}\frac{p1}{p}a_M`$ and $`x_p(a_{p,max})>\beta `$ for $`p3`$. Similar results for the symmetric case is obtained by putting $`\gamma =1`$.
Differentiating $`x_p(a)`$ w.r.t. $`a`$, we get
$$x_p^{}(a)=\gamma a^{p2}(p1)(1\frac{p}{p1}a\alpha ).$$
(C.1)
Then, at $`a=a_{p,max}\frac{p1}{p}a_M`$, $`x_p(a)`$ has the maximum value
$$x_p(a_{p,max})=\frac{\gamma }{p}(\frac{p1}{p})^{p1}a_M^{p1}=\frac{1}{p(1\beta )}(\frac{p1}{p})^{p1}\alpha ^{2p}.$$
(C.2)
Since $`0<\alpha <\beta <1`$ and $`p3`$,
$$x_p(a_{p,max})>\frac{1}{p(1\beta )}(\frac{p1}{p})^{p1}\beta ^{2p}=\frac{1}{p}(\frac{p1}{p})^{p1}\frac{\beta }{g(\beta )},$$
(C.3)
where $`g(\beta )\beta ^{p1}(1\beta )>0`$. $`g(\beta )`$ has the maximum value at $`\beta =\frac{p1}{p}`$ and
$$g(\frac{p1}{p})=\frac{1}{p}(\frac{p1}{p})^{p1}.$$
(C.4)
Thus,
$$x_p(a_{p,max})>\frac{1}{p}(\frac{p1}{p})^{p1}\frac{\beta }{g(\frac{p1}{p})}=\beta .$$
(C.5)
## D Proof of $`H(x_1)H(x_2)\mathrm{}H(x_{p1})=RL\mathrm{}L`$ for $`a_{p,0}a`$
In this appendix, we prove that for $`a_{p,0}a`$, the symbolic sequence for $`x_1,x_2,\mathrm{},x_{p1}`$ is $`RL\mathrm{}L`$. We prove this in the asymmetric case. We obtain the results for the symmetric case by putting $`\gamma =1`$.
First, we prove the following relation,
$$a_{p,0}>\frac{\beta }{\alpha }.$$
(D.1)
For $`p3`$, the following relation holds,
$$x_p(\frac{\beta }{\alpha })=(\frac{\beta }{\alpha })^{p2}\beta >\beta .$$
(D.2)
Since $`x_p(a)=\gamma a^{p1}(1a\alpha )`$ is a strictly decreasing continuous function for $`a_{p,max}a`$, and $`a_{p,max}<a_{p,0}`$, we obtain $`a_{p,0}>\frac{\beta }{\alpha }`$.
Now, let us fix $`a`$ such as $`aa_{p,0}`$. Then, from the relation (D$``$1)
$$x_1=a\alpha a_{p,0}\alpha >\beta .$$
Thus, $`H(x_1)=R`$. Then,
$$x_2=\gamma a(1x_1)=\frac{\gamma a^{p1}(1a\alpha )}{a^{p2}}=\frac{x_p}{a^{p2}}\frac{\beta }{a^{p2}}\frac{\beta }{a_{p,0}^{p2}}<\frac{\beta }{a_{p,0}}<\alpha .$$
Thus, $`x_2I_L`$. Now, let us assume $`x_n<\beta `$ for $`2np2`$. Then,
$$x_{n+1}=ax_n=a^n\gamma (1a\alpha )=\frac{a^{p1}\gamma (1a\alpha )}{a^{p1n}}<\frac{\beta }{a^{p1n}}<\frac{\beta }{a}\frac{\beta }{a_{p,0}}<\alpha .$$
Therefore, $`x_{n+1}I_L`$. Thus,
$$x_2<x_3<\mathrm{}<x_{p1}<\alpha .$$
Q.E.D.
## E Proof of $`a_ca_{3,0}<a_{4,0}<\mathrm{}<a_{p,0}<a_{p+1,0}<\mathrm{}<a_M`$
In this appendix, we prove the following relations,
$$a_ca_{3,0}<a_{4,0}<\mathrm{}<a_{p,0}<a_{p+1,0}<\mathrm{}<a_M.$$
We only have to prove it for the asymmetric case because $`\gamma =1`$ reduces to the symmetric case.
Let us estimate $`x_{p+1}(a_{p,0})`$ for $`p3`$,
$$x_{p+1}(a_{p,0})=a_{p,0}x_p(a_{p,0})=a_{p,0}\beta >\beta .$$
Since $`x_{p+1}(a)`$ is decreasing for $`a_{p+1,max}a`$ and $`a_{p+1,max}<a_{p+1,0}`$, $`a_{p,0}<a_{p+1,0}`$ follows.
Let us prove the first relation $`a_ca_{3,0}`$. For $`a_{3,0}<a<a_{3,1}`$, there is a stable periodic three solution. If $`a_c>a`$, there is a stable periodic solution with period $`2^m`$ with some integer $`m`$, and no other stable solution. Thus, $`a_ca_{3,0}`$. Q.E.D. |
warning/0001/physics0001007.html | ar5iv | text | # The preliminary results of fast neutron flux measurements in the DULB laboratory at Baksan
## Introduction
It is well known that one of the main sources of a background in underground physics experiments (such as the investigation of solar neutrino flux, neutrino oscillations, neutrinoless double beta decay, and the search for annual and daily Cold Dark Matter particle flux modulation) are fast neutrons originating from the surrounding rocks. The sources of the fast neutrons are ($`\alpha ,n`$) reactions on the light elements contained in the rock (C, O, F, Na, Mg, Al, Si). Neutrons from spontaneous fission of <sup>238</sup>U take an additional contribution in a total fast neutron flux of about 15-20%. Several research groups have investigated the neutron background at different underground laboratories . Some of them used $`{}_{}{}^{6}Li`$-dopped liquid scintillator technique , and others used in addition a Pulse Shape Discrimination technique .
The measurements of fast neutron flux in the Deep Underground Low Background Laboratory of Baksan Neutrino Observatory (DULB BNO) have been performed with using of a special, high-sensitive fast neutron spectrometer . This laboratory is located under Mt. Andyrchy (Northen Caucasus Mountains, Russia) in a tunnel that penetrates 4.5 km into the mountain, at a depth of 4900 meters of water equivalent.
The results of such measurements lead to a conclusion that a neutron background places a severe limitation on the sensitivity of current and planned experiments. Owing this fact, the development of new cost-effective, high-strength radiation shielding against neutrons becomes a very important task for modern non-accelerator physics experiments. For such purposes the relative neutron shielding properties of several commonly available natural materials were investigated too. Specially, these materials are planned for use in the construction of large-volume underground facilities which will be covered with suitable shielding materials and are situated in the DULB Laboratory at Baksan.
## Neutron detector
The spectrometer was constructed to measure low background neutron fluxes at the level up to $`10^7cm^2s^1`$ in the presence of intensive gamma-ray background.
The detector consists of 30 l liquid organic scintillator viewed by photomultipliers with 19 neutron counters ($`{}_{}{}^{3}He`$ proportional counters) uniformly distributed through the scintillator volume (see for detail). The spectrometer schematic view and the principle of operation are shown in Fig. 1.
Fast neutrons with $`E_n>`$ 1 MeV entering the liquid scintillator (LS) are moderated down to thermal energy, producing a LS signal. Then they diffuse through the detector volume to be captured in $`{}_{}{}^{3}He`$ counters or on protons in the scintillator. The LS signal starts the recording system. After triggering the system waits a signal from any of the helium counters for a specific time. This time window corresponds to the delay time between correlated events in the scintillator and in the helium counters. This is one of specific features of the detector. The signal from the LS is ’marked’ as a coincident with a neutron capture in the $`{}_{}{}^{3}He`$ counters only in the case if a single counter is triggered during the waiting period. An amplitude of the ’marked’ LS signal corresponds to an initial neutron energy. This method allows us to suppress the natural $`\gamma `$-ray background considerably.
The described event discriminating procedure allows us to measure extremely low neutron fluxes at the level up to $`10^7cm^2s^1`$ reliably even if the LS counting rate is as large as several hundred per second. The dead time of the detector is equal to the delay time (variable value, but generally about $`120\mu s`$) plus about $`400\mu s`$, which is needed to analyze a LS event whether it corresponds to neutron or not. The detection efficiency depends in a complicated manner on the response function of the detector. As a rough estimation, we use the value of the efficiency, which is equal to $`0.04\pm 0.02`$ in the energy range from 1 to 15 MeV. This is based on preliminary measurements performed with a Pu-Be source. Owing this fact, an absolute values of the neutron fluxes can be estimated with an uncertainty of 50% on the basis of available calibration data. The delay time is a specific feature of the detector and depends on the detector design. The acquisition system allows us to measure the delay time for the neutron events directly. Such measurements were carried out using a $`PuBe`$ source with a time window selected to be equal to $`300\mu s`$. A typical delay time distribution is shown in Fig. 2 A fitting procedure leads to a time constant of $`T_{1/2}55\mu s`$. According to this result it is sufficient to select the time window to be equal $`120\mu s`$ for an actual measurement.
## Measurements
### A The geometry
It has been mentioned that we have no yet precise information about the detection efficiency, that is why one can calculate absolute value of neutron fluxes with only $`50\%`$ certainty. However, it is possible to measure the relative neutron absorption abilities of various shields with high precision. This information will be very useful for development of new low background experiments and searching for cost-effective neutron absoption shields. Such measurements were carried out in the DULB BNO with using the described neutron spectrometer. This new laboratory, consisted of 8 separate counting facilities, is located under Mt. Andyrchy in a tunnel, which penetrates 4.3 $`km`$ into the mountain, at a depth of 4900 $`mw.e.`$
Quartzite and serpentine were selected as materials to be tested because of comparatively low concentrations of uranium- and thorium-bearing compounds contained in these rocks. For instant, the measured concentrations of uranium and thorium for rock serpentine are about $`10^8`$g/g in comparison with $`10^6`$g/g for the surrounding rock. As for potassium ($`{}_{}{}^{40}K`$) contained in serpentine, it has been found less than $`10^8`$g/g in comparison with $`10^6`$g/g for the surrounding rock. Measurements of gamma-activity of different rock samples have been performed with using a well-type NaI gamma spectrometer with level of sensitivity of about $`10^9`$g/g, operated in one of the underground low counting facilities at BNO . The measured Th, U, and K concentrations in different rock samples are given in Table I.
Four series of measurements were performed with the neutron spectrometer surrounded by different radiation shields. In the first series the spectrometer was surrounded with a lead shield 4 $`cm`$ thickness (to reduce the natural gamma-ray counting rate), and measurements of the natural neutron background radition field existing in the open experimental site were performed. In the second and third series the spectrometer was surrounded with shields of quartzite and serpentine, respectively. The rock shields consisted of broken pieces of various sizes, ranging from 1 $`cm`$ to 15 $`cm`$, with an effective shield thickness of 35 $`cm`$ in all directions. The mean relaxation length of fast neutrons in these shields is about 15 $`cm`$ (25 $`g/cm^2`$ for quartzite and 21 $`g/cm^2`$ for serpentine). In the fourth series we measured the internal background of the detector using a neutron-absorbing shield consisted of 40-cm thick section of polyethylene containing an admixture of boron and water about 30 $`cm`$ thick. Schematic view of one of the investigated neutron shield and cross-section of the DULB experimental site are shown in Fig.3.
### B Calibration
A $`{}_{}{}^{60}Co\gamma `$-source has been used to calibrate the LS-channel. The energy of the middle of the Compton edge was assumed to be equal to 1 $`MeV`$ in the electron energy scale, which corresponds to $`3MeV`$ in the neutron energy scale (see Fig.4a). A Pu-Be source was used to calibrate the NC-channel of $`{}_{}{}^{3}He`$ counters. The spectrum produced by the Pu-Be source in the $`{}_{}{}^{3}He`$ counters has a specific shape due to a wall effect which distorts the counter event spectrum (see Fig. 4b). In spite of this distortion, the range of energies observed for true neutron events is less narrow compared to the broad background spectrum produced by internal alphas. Using of the only events from the neutron window coincident with LS signals makes it possible to suppress the internal background of the detector.
### C Conditions of measurements
Main conditions for all series of measurements, such as measuring times, LS- and NC-counting rates are given in the Table II.
The typical exposure time for each series was a few weeks. The $`\gamma `$-ray background in the open experimental site is high enough that leads to $`\gamma `$-counting rate in the LS-channel of about $`700s^1`$. Due to this fact, following values of dead time were determined for different series: $`12\%`$ of the total exposure time for measurements with the lead shield, $`4.3\%`$ for quartzite series, $`2.7\%`$ for serpentine series, and $`1.5\%`$ for measurements with the polyethylene/water shield. To calculate the true neutron counting rates a proper dead time correction has been performed.
## Data treatment and results
Contamination of $`{}_{}{}^{222}Rn`$ gas inside the experimental site can make a considerable contribution (up to $`20\%`$) to the background $`\gamma `$-counting rate, which can influence results of the performed measurements because $`{}_{}{}^{222}Rn`$ activity can vary significantly for a period of a measurement.
To suppress the count rate variation effect we used a special procedure for treatment of experimental data. It consists of the following steps.
Two types of data files are stored as a result of a measurement. One of them contains the information about neutron energy losses ( LS-signal amplitudes), $`{}_{}{}^{3}He`$ counters signal amplitudes, and delay time for each ’neutron’ candidate event. Data accumulation was stopped every half-hour and overall numbers of NC-counts, LS-counts, LS-counts above 1 $`MeV`$, and elapsed time were saved in a file. Total background $`\gamma `$ \- spectra for every half-hour run were measured simultaneously and saved in a second file to make it possible to take into account a time variation of the background $`\gamma `$-counting rate.
We consider three contributions into the experimentally measured counting rate $`R_{meas}`$: the random coincidence rate $`R_{rnd}`$, the internal detector background counting rate $`R_{bkg}`$, and the ’neutron’ counting rate $`R_n`$, so that
$$R_n=R_{meas}R_{rnd}R_{bkg}$$
(1)
We have made obvious assumption that the total background $`\gamma `$-spectrum and the random coincidence spectrum have the same shapes. To obtain random coincidence spectrum for further subtraction procedure the total background $`\gamma `$-spectrum has been normalized with a factor corresponding to the calculated random coincidence rate. The maximal evaluation for the random coincidence rate, if the LS- and NC- events are absolutely independent, can be calculated by the following way:
$$R_{rand}=r_\gamma r_n^w\mathrm{\Delta }t,$$
(2)
where $`r_\gamma `$ is the $`\gamma `$-rate, $`r_n^w`$ is the $`{}_{}{}^{3}He`$-counters counting rate in the determined neutron energy window, $`\mathrm{\Delta }t`$ is the time window. In the case of the performed measurements $`(R_{LS}R_{He})`$, this evaluation is very close to the real counting rate of random coincidences. Due to a variation in time of the $`{}_{}{}^{222}Rn`$ activity, the current value of $`r_\gamma ^i`$ depends on time too. Owing to this fact, we applied the described subtraction procedure to each half-hour run with corresponding current value of $`R_{rnd}^i`$, and then summarized resulting neutron spectra in a total serial spectrum. The accumulated LS-spectra of all coincidented events ($`R_{meas}`$) and the recalculated spectra of random coincidence ($`R_{rnd}`$) for the no-shield, quartzite, and serpentine series are presented in Fig. 5.
An internal detector background spectrum $`R_{bkg}`$ has been accumulated inside the neutron-absorbing shield consisting of polyethylene and water. Obtained counting rate of the internal background correlated (neutron-type, but non-neutron) events was measured as 27 counts per hour, which in terms of a neutron flux corresponds to ($`8.1\pm 0.5`$) $`10^7s^1cm^2`$. The residual LS-spectra ($`R_{meas}R_{rnd}`$) in comparison with the internal background LS-spectrum ($`R_{bkg}`$) are presented in Fig. 6.
Performing the total subtraction procedure in according with the equation (1) we obtain values of the neutron counting rate $`R_n`$ for the no-shield, quartzite, and serpentine series. Taking into account the detection efficiency uncertainty ($`\epsilon =0.04\pm 0.02`$) the obtained values of fast neutron fluxes (above 700 $`keV`$ of neutron energy) are presented here in a following way:
$`a(3.5\pm 1.1)`$ $`10^7s^1cm^2`$ for the no-shield measurement,
$`a(2.9\pm 1.1)`$ $`10^7s^1cm^2`$ for quartzite shield,
$`a(0.6\pm 0.7)`$ $`10^7s^1cm^2`$ for serpentine shield,
were a = ($`\epsilon `$ \+ $`\mathrm{\Delta }\epsilon `$)/$`\epsilon `$.
One can see that the resulting neutron flux measured when the serpentine shield was in place were found to be at about the minimum level of sensitivity of the spectrometer. It means that a neutron background inside the serpentine shield is consisted with a neuron flux less than $`0.7`$ $`10^7s^1cm^2`$. It indicates that serpentine is indeed clear from uranium and thorium, and is, therefore, the most likely candidate for use as a cost-effective neutron shield component material for large-scale low background experiments.
A delay time distribution analysis was performed to understand the origin of a high level of the internal detector background.
## Delay time distributions
Decays of Bi and Po radioactive isotopes, such as
$${}_{}{}^{214}Bi(e,\stackrel{~}{\nu })\frac{164\mu s}{}^{214}Po(\alpha )\mathrm{},$$
(3)
which can take place in the helium counter walls, have been considered as main possible sources of the significant internal background. To imitate an actual neutron event beta decay of $`{}_{}{}^{214}Bi`$ can fire the liquid scintillator, followed by a delayed capture $`\alpha `$ \- signal from Po decay in helium counters. The delay time distribution of the neutron-type coincidented events obtained for the series in the water shield is shown in Fig. 7. Fitting procedure leads to the time constant $`T_{1/2}=164\mu s`$.
It means that, as it was supposed, the origin of the internal background of our detector is mostly due to contamination of $`{}_{}{}^{214}Bi`$ in the $`{}_{}{}^{3}He`$-counter walls. The delay time distributions for other series of measurements are shown in Fig. 8.
The following fitting function was used to analize these distributions ($`t`$ is expressed in $`\mu `$s):
$$A+Ne^{tln2/55}+Be^{tln2/164},$$
(4)
where A is a constant, N is an amplitude corresponding to neutrons and B corresponds to internal background. The ratio N/B, which was obtained in this manner, decreases from measurements in the lead shield to the measurements in the serpentine shield.
## Conclusions
The main results of the measurements can be summarized as follows.
(I). The preliminary results obtained from the fast neutron spectrum accumulated in the open experimental site of the DULB Laboratory at Baksan is consisted with a neutron flux (for neutrons with energy above 700 keV) estimated as values from $`5.3\times 10^7cm^2s^1`$ to $`1.8\times 10^7cm^2s^1`$ depending on the present uncertainty in determination of the detection efficiency.
(II). The neutron spectrometer sensitivity in a shielded experimental site is estimated as $`0.5\times 10^7cm^2s^1`$ for a measuring time of about 1000 h.
(III). It is shown that the main source of the detection sensitivity limitation, rather then random coincidences, is the internal background of the spectrometer, which is mostly due to the presence of $`\alpha `$-particle emitters ($`{}_{}{}^{214}Bi^{214}Po`$ decays) in the $`{}_{}{}^{3}He`$-counters walls.
(IV). The achieved neutron background inside the serpentine shield is consisted with a neutron flux less than $`0.7`$ $`10^7s^1cm^2`$. It indicates that serpentine is one of the more likely candidate for use as a cost-effective neutron shield component material for large-scale low background experiments.
We have obtained the presented results using the simple event discrimination procedure and did not use pulse shape discrimination yet. Nevertheless, it takes us a possibility to measure extremely low neutron fluxes up to $`10^7cm^2s^1`$ even when external $`\gamma `$-counting rate is more than $`200s^1`$.
###### Acknowledgements.
We are grateful to I.I.Pyanzin for the management in proving of reserves and quarrying of the domestic ultra basic rock samples. We thank P.S.Wildenhain for careful reading of this article and his critical remarks. We acknowledge the support of the Russian Foundation of Basic Research. This research was made possible in part by the grants of RFBR No. 98–02 16962 and No. 98-02-17973. |
warning/0001/cond-mat0001429.html | ar5iv | text | # Tractable non-local correlation density functionals for flat surfaces and slabs
## I Introduction
The density-functional theory (DFT), with its local-density (LDA) and semilocal generalized-gradient approximations (GGA), is not only successful in numerous applications to individual molecules and dense solids. It is also under intense development, for instance, in order to include non-local effects, such as van der Waals (vdW) forces. The latter are needed in order to allow DFT to describe sparse matter. A unified treatment of vdW forces at large and asymptotic molecular separations is available, and a description at short distances and at overlap is striven for. An accurate calculation for the interaction of two He atoms has been given, and recently the first microscopic (RPA) calculation of the vdW interaction between two self-consistent jellium slabs has been reported and given a density functional (DF) account. The ultimate challenge is to construct an approximate vdW DF that is generally applicable, efficient, and accurate.
We here propose an explicit form for the vdW DF that applies to flat surfaces, test it successfully against these slab results, and apply it to two parallel flat semi-infinite metal surfaces. This is a case with relevance for many physical situations, including wetting and atomic-force microscopy (AFM). Compared to the vdW-DF approximation proposed by Dobson and Wang, the virtues of our functional are the computational simplifications gained from choosing a particular sub-class of response functions, utilizing a differential formulation and sparse matrices, and recognizing the insensitivity to the details of the density profiles, simplifications which might transfer even to three dimensions.
The ubiquitous van der Waals force plays an important role for numerous physical, chemical, and biological systems, such as physisorption, vdW complexes, vdW bonds in crystals, liquids, adhesion, and soft condensed matter (e.g., biomacromolecules, biosurfaces, polymers, and membranes).
The DFT expresses the ground-state energy of an interacting system in an external potential $`v(𝐫)`$ as a functional $`E[n]`$ of the particle density $`n(𝐫)`$, which has its minimum at the true ground-state density. The Kohn-Sham form of the functional makes the scheme a tractable one, as it leads to equations of one-electron type, and accounts for the intricate interactions among the electrons with an exchange-correlation (XC) functional $`E_{\mathrm{xc}}[n]`$. This XC-energy functional can be expressed exactly as an integral over a coupling constant ($`\lambda `$), the so-called adiabatic connection formula (ACF),
$$E_{\mathrm{xc}}=_0^1𝑑\lambda _0^{\mathrm{}}\frac{du}{2\pi }\mathrm{Tr}\left[\chi (\lambda ,iu)V\right]E_{\mathrm{self}},$$
(1)
where $`V(𝐫,𝐫^{})=1/\left|𝐫𝐫^{}\right|`$ and where the density-density correlation function is denoted by $`\chi (𝐫,𝐫^{},iu;\lambda )`$. $`E_{\mathrm{self}}`$ is the Coulomb self-energy of all electrons, which is exactly cancelled by a corresponding term in $`\chi (\lambda ,iu)V`$. Equation (1) shows a truly non-local XC interaction and is a starting point for approximate treatments, local (LDA), semilocal (GGA) and non-local ones.
The LDA and GGA are completely unable to express the vdW interaction in a physically sound way. The exact XC energy functional, on the other hand, of course encompasses such interactions. The basic problem of making DFT a working application tool also for sparse matter is to express the truly non-local vdW interactions between the electrons in the form of a simple, physical, and tractable DF. Equation (1) is then the starting point. Along these lines, we have proposed extensions of the DFT to include van der Waals interactions, with very promising results for the interaction between two atoms or molecules, an atom and a surface, and two parallel surfaces, respectively. This now unified approach applies for separated systems, i.e. when the electrons of the interacting fragments have negligible overlaps.
The corresponding asymptotic expressions have singular behaviors at short separations $`d`$. Yet one knows that the vdW forces are finite. They should go smoothly over to the XC forces that apply in the interior of each electron system. This phenomenon is often called damping, or *saturation*. Approximate saturation functions have been proposed, in particular for the cases of vdW molecules and physisorbed particles.
The key difficulty in extracting the vdW DF from Eq. (1) is the computational complexity. A direct solution gives simply too many operations on the computer. The guideline for our reduction of the number of such operations is to exploit analytical advantages of RPA-like approximations, to focus on the key quantity, to recast the integral formulation into a differential one, leading to a sparse-matrix computation, and to make maximal use of symmetry.
Our exploratory study here concerns cases with vdW forces between two flat parallel model systems. We first test our approximate functional on the model system of two self-consistent jellium slabs, utilizing the recent RPA results, which gives the size of the correlation-interaction energy per unit area, showing saturation. We then test our DF against accurate calculations of the surface correlation energy, showing an excellent agreement. After these successful tests we make predictions on two parallel semi-infinite jellia.
## II General Theory
The $`\lambda `$ integration of Eq. (1) can be performed analytically in some cases, such as in the random-phase approximation (RPA). In 1957 Gell-Mann and Bruckner (GMB) presented the RPA correlation energy as a selected summation of ring diagrams, which gives a logarithmic form. Their study concerns the *homogeneous* electron gas, where equations simplify thanks to the three-dimensional translational invariance and plane waves. Here we treat systems with less symmetry.
By virtue of the fluctuation-dissipation theorem, the density-density correlation function $`\chi `$ is equal to the density change $`\delta n`$ induced by an external potential $`\mathrm{\Phi }_{\mathrm{ext}}`$, i.e. $`\delta n=\chi \mathrm{\Phi }_{\mathrm{ext}}`$. It satisfies
$$\chi (\lambda ,iu)=\stackrel{~}{\chi }(iu)+\lambda \stackrel{~}{\chi }(iu)V\chi (\lambda ,iu),$$
(2)
where $`\stackrel{~}{\chi }`$ is the density response to a fully screened potential $`\mathrm{\Phi }`$, i.e. $`\delta n=\stackrel{~}{\chi }\mathrm{\Phi }`$. We assume here that the coupling dependence of $`\stackrel{~}{\chi }`$ can be neglected, when performing the $`\lambda `$ integral in Eq. (1). This is true in the random-phase approximation, where $`\stackrel{~}{\chi }`$ is the density response function for $`\lambda =0`$, and is also true for the approximate dielectric functions, which we use here. Equation (1) then becomes
$$E_{\mathrm{xc}}=_0^{\mathrm{}}\frac{du}{2\pi }\mathrm{Re}\left[\mathrm{Tr}\left[\mathrm{log}(1\stackrel{~}{\chi }(iu)V)\right]\right]E_{\mathrm{self}},$$
(3)
where the real part means the principal branch.
To simplify this expression and to get a functional in terms of the electron density $`n(𝐫)`$, we have to focus on the key target, the non-local part, introduce key quantities, and rewrite the expressions, in order to make physically sound and computationally efficient approximations. It is more convenient to introduce the polarizability or dielectric function instead of $`\stackrel{~}{\chi }`$. The polarizability $`\alpha `$ (a matrix in the spatial positions) is defined by the relation $`𝐏=\alpha 𝐄`$, where $`𝐏`$ is the polarization. We have
$$\delta n=𝐏=\alpha 𝐄=\alpha \mathrm{\Phi },$$
(4)
so that from the definition of $`\stackrel{~}{\chi }`$, one has $`\stackrel{~}{\chi }=\alpha `$. In turn the dielectric function is given by $`ϵ\mathrm{𝟏}+4\pi \alpha `$. In terms of $`ϵ`$ Eq. (3) then transforms to
$$E_{\mathrm{xc}}=_0^{\mathrm{}}\frac{du}{2\pi }\mathrm{Re}\left[\mathrm{Tr}\left[\mathrm{log}(ϵG)\right]\right]E_{\mathrm{self}},$$
(5)
where we have introduced the Coulomb Green’s function $`G=V/4\pi `$ and used $`^2G=1`$. The $`\mathrm{Tr}\left[\mathrm{log}\right]`$ expression gives great advantage for the further analytical and numerical treatment. The only approximation made so far is the neglect of the coupling constant dependence of $`\stackrel{~}{\chi }`$ when doing the coupling constant integration. This is not an additional approximation either in the RPA or for the approximate $`ϵ`$’s we use here.
In order to develop long-range functionals, one may substitute approximations for the dielectric function based on the free electron gas into Eq. (5). To obtain tractable expressions it will normally be necessary to make still further approximations. In this case it is desirable to use the additional approximations only for the non-local part of $`E_{\mathrm{xc}}`$ so as to avoid destroying the accuracy of the LDA in the high-density regions. Ideally one would subtract from Eq. (5) the LDA version of the same approximation, and would envisage adding back a better version of the LDA. Here we make a similar, but more tractable subtraction, which takes the form
$$E_{\mathrm{xc}}^0=_0^{\mathrm{}}\frac{du}{2\pi }\mathrm{Re}\left[\mathrm{Tr}\left[\mathrm{log}(ϵ)\right]\right]E_{\mathrm{self}}.$$
(6)
$`E_{\mathrm{xc}}^0`$ has the property that it is a good approximation for a slowly varying system, becoming exact for a uniform system. For density variations slow on the scale of the range or width of $`ϵ(𝐫,𝐫^{})`$, it agrees with the LDA, the trace in Eq. (6) replacing the integral over density.
Subtracting Eq. (6) from Eq. (5), one obtains
$$E_{\mathrm{xc}}^{\mathrm{nl}}=_0^{\mathrm{}}\frac{du}{2\pi }\mathrm{Re}\left[\mathrm{Tr}\left[\mathrm{log}(ϵ^1ϵG)\right]\right].$$
(7)
We will call this the non-local exchange-correlation energy, although for models more general than those used in this paper, an additional short-range correction must be applied to make $`E_{\mathrm{xc}}^0`$ correspond precisely to the LDA, and hence to make $`E_{\mathrm{xc}}^{\mathrm{nl}}`$ the deviation from the LDA.
The approximations considered in this paper contain no non-local exchange component, in effect making Eq. (7) our approximation for the non-local correlation energy $`E_c^{nl}`$. Using this fact, together with $`\mathrm{Tr}\left[\mathrm{log}x\right]=\mathrm{log}(detx)`$ and $`^2G=1`$, we obtain
$$E_c^{\mathrm{nl}}=_0^{\mathrm{}}\frac{du}{2\pi }\mathrm{log}\left|det(1+ϵ^1[,ϵ]G)\right|$$
(8)
where the notation $`[A,B]`$ means the commutator. We will later use the fact that Eq. (8) involves only the determinant to good advantage.
## III Method for Planar Geometries
Now we are in a situation to discuss what $`ϵ`$ to use. We will in this paper concentrate on the simple case of jellium systems. Our aim is to find an *efficient* way of exploring the planar translational invariance of the jellium system – not only to decrease the number of spatial integrations. The major difficulty in evaluating Eq. (8) is the determinant, which is $`𝒪(N^3)`$ in the general case, $`N`$ being the number of grid points in a discrete representation. This holds true even in one dimension. So, instead of allowing a completely general $`ϵ`$, we aim at approximations resulting in *differential* operators only, for which the determinant is known to be $`𝒪(N)`$, hence significantly simpler to calculate.
In the particular case of planar translational invariance, i.e. for planar surfaces or slabs, we use an approximate form that is made local in the coordinate perpendicular to the surface,
$$ϵ(𝐫,𝐫^{})=\delta (zz^{})\frac{d^2k}{(2\pi )^2}ϵ_k(z)e^{i𝐤(𝐫𝐫^{})},$$
(9)
where $`𝐤`$ is a wave vector parallel to the surface. Keeping the fully non-local form along the symmetry plane allows, e.g., the effect of the cutoff, which was introduced artificially in previous approximations of this type, to occur in a natural way laterally.
The first thing we note about Eq. (9) is that we easily form the inverse
$$ϵ^1(𝐫,𝐫^{})=\delta (zz^{})\frac{d^2k}{(2\pi )^2}ϵ_k(z)^1e^{i𝐤(𝐫𝐫^{})}.$$
(10)
Evaluating the commutator then yields
$$ϵ^1[,ϵ]=\widehat{z}\delta (zz^{})\frac{d^2k}{(2\pi )^2}\frac{ϵ_k^{}(z)}{ϵ_k(z)}e^{i𝐤(𝐫𝐫^{})},$$
(11)
where the prime indicates differentiation with respect to $`z`$. In what follows we shall substitute $`l(z)=\mathrm{log}(ϵ(z))`$, yielding $`l^{}(z)=ϵ^{}(z)/ϵ(z)`$.
In the same basis we express the Green’s function,
$$G(𝐫𝐫^{})=\frac{d^2k}{(2\pi )^2}G_k(zz^{})e^{i𝐤(𝐫𝐫^{})},$$
(12)
where
$$G_k(zz^{})=\frac{1}{2k}e^{k\left|zz^{}\right|}.$$
(13)
Since the logarithm of Eq. (7) can be expanded in powers, the integration over $`k`$ may be singled out, and we may express the non-local correlation energy per surface area ($`A`$) as
$$E_c^{\mathrm{nl}}/A=_0^{\mathrm{}}\frac{du}{2\pi }\frac{d^2k}{(2\pi )^2}\mathrm{log}\left|det(1+ł_k^{}_zG_k)\right|.$$
(14)
In Eq. (14), the determinant is given in terms of integral operators. To take advantage of the locality of the Laplacian, we use $`(_z^2k^2)G_k=1`$ to express it in terms of differential operators
$$det\left(1+l_k^{}_zG_k\right)=\frac{\varphi }{\varphi _0},$$
(15)
where
$$\varphi =det(_z^2k^2+l_k^{}_z)$$
(16)
and where $`\varphi _0`$ is the empty space ($`ϵ=1`$) value of Eq. (16). The step from Eq. (14) to Eq. (15) requires that the differential operators are defined throughout the whole space.
Our ultra-fast method is made possible by the observation that the determinants in Eq. (15) can be written down, not only for the full system, but also for a subdivision of it. Related determinants for the subsystem satisfy a simple second-order differential equation as a function of subsystem size. Thus by a simple renormalization, one may evaluate $`E_c^{\mathrm{nl}}`$ with the same effort as finding the charge induced by an applied electric field. A similar relation holds also in several dimensions, which will be explored in another paper.
To make this more concrete, let us suppose that $`ϵ_k(z)`$ varies only in the interval $`0<z<L`$ (which will eventually be extended to infinity) and takes the same value at either end point. This is the case for parallel surfaces or slabs of identical materials. Then for each value of $`z`$ we can define a determinant $`\varphi (z)`$ for the subsystem extending from $`0`$ to $`z`$. It is clear then from Eqs. (14), (15), and (16) that $`E_c^{\mathrm{nl}}`$ is given by
$$E_c^{\mathrm{nl}}/A=\underset{L\mathrm{}}{lim}_0^{\mathrm{}}\frac{du}{2\pi }\frac{d^2k}{(2\pi )^2}\mathrm{log}\frac{\varphi (L)}{\varphi _0(L)}.$$
(17)
As discussed in Appendix A, the determinants $`\varphi (z)`$ and $`\varphi _0(z)`$ individually have oscillating signs that do not occur in their quotient. However, the envelope determinants $`\stackrel{~}{\varphi }(z)`$ and $`\stackrel{~}{\varphi }_0(z)`$ can be scaled so that they satisfy the simple differential equation
$$(ϵ_k\stackrel{~}{\varphi }^{})^{}=k^2ϵ_k\stackrel{~}{\varphi },$$
(18)
together with the boundary conditions that $`\stackrel{~}{\varphi }(0)=0`$ and $`\stackrel{~}{\varphi }(L)=1`$. In terms of $`\stackrel{~}{\varphi }`$, we obtain
$$E_c^{\mathrm{nl}}/A=\underset{L\mathrm{}}{lim}_0^{\mathrm{}}\frac{du}{2\pi }\frac{d^2k}{(2\pi )^2}\mathrm{log}\frac{\stackrel{~}{\varphi }^{}(0)}{\stackrel{~}{\varphi }_0^{}(0)},$$
(19)
where the prime indicates differentiation with respect to $`z`$, which in this case is the subsystem size. However, note that $`\stackrel{~}{\varphi }`$ is also just the electrostatic potential as a function of distance $`z`$, within a system having a potential difference across it along with a prescribed variation in $`ϵ`$. Thus the calculation of the determinant becomes a simple electrostatic problem which is easily solved.
To illustrate this, consider the case of two identical parallel surfaces a distance $`d`$ apart, when $`d`$ is much larger than the thicknesses of the surface-healing layers. Solving Eq. (18) for the described boundary conditions then becomes a trivial matching problem (see Appendix B), which after insertion into Eq. (19) immediately leads to the Lifshitz formula
$$E_c^{\mathrm{nl}}/A=_0^{\mathrm{}}\frac{du}{2\pi }\frac{d^2k}{(2\pi )^2}\mathrm{log}\left|1\rho ^2e^{2kd}\right|+2\gamma _{\mathrm{nl}}.$$
(20)
Here $`\rho =(ϵ_b1)/(ϵ_b+1)`$, $`ϵ_b`$ being the bulk dielectric function, and $`\gamma _{\mathrm{nl}}`$ is defined by
$$\gamma _{\mathrm{nl}}=(E_c^{\mathrm{nl}}(d\mathrm{})E_c^{\mathrm{nl}}(0))/2A.$$
(21)
Since, by construction, $`E_c^{\mathrm{nl}}=0`$ for a uniform ($`d=0`$) system, $`\gamma _{\mathrm{nl}}`$ may equivalently be defined as the non-local correlation contribution to the surface tension of a single surface.
The original ACF Eq. (1) is now reduced to a set of simple electrostatic calculations, each one being an $`𝒪(N)`$ operation instead of $`𝒪(N^2)`$, a major simplification. Of course the success of the method depends on how well we can reproduce the true dielectric function using our approximate form Eq. (9). The only approximation made so far is the assumption of a local dielectric function perpendicular to the surface.
## IV Approximate Dielectric Function
Equation (17) or (19) provides the basis for a functional that describes the van der Waals interaction between planar objects. To turn these equations into density functionals, we have to introduce quantities that depend on the density $`n(𝐫)`$. Our suggestion is based on an approximate dielectric function $`ϵ_k`$ that depends on the local density $`n(𝐫)`$. It utilizes experiences from the homogeneous electron gas and from experimental studies of the dynamical structure factor $`S(𝐪,\omega )\mathrm{Im}\left[\frac{1}{ϵ(𝐪,\omega )}1\right]`$, where $`\mathrm{}𝐪`$ and $`\mathrm{}\omega `$ are the momentum and energy losses, respectively, of a photon or a charged particle being scattered while passing a bulk sample. There is a peak in $`S(𝐪,\omega )`$, the plasmon peak, sharp in the ideal electron gas and of varying width in real materials. This peak carries most of the spectral strength and has $`\omega `$ equal to the plasma frequency $`\omega _p`$ as $`q0`$, and then a dispersion with a limiting behavior $`\omega _qq^2/2m`$, the kinetic energy of one electron, in the impulse-approximation, valid in the Compton-scattering limit, $`q\mathrm{}`$. In the electron-plasmon coupling one focuses on the *magnitude* and *position* of the sharp plasmon peak, and neglects the broadening, i.e. $`ϵ`$ is described in a plasmon-pole approximation. A dispersion law like
$$\omega _q^2=\omega _p^2+(v_Fq)^2/3+(q^2/2m)^2$$
(22)
has been shown to efficiently account for the average behavior of plasmon-like excitations and for correlation properties of the homogeneous electron gas. Introducing the electron density via $`\omega _p^2(z)=4\pi n(z)`$ and $`v_F(z)=(3\pi ^2n(z))^{1/3}`$ in Hartree units, $`n(z)`$ being the electron density profile in this planar case, the dielectric function can be written as
$$ϵ(z,k,iu)=1+\frac{\omega _p^2(z)}{u^2+(v_F(z)q(k))^2/3+q^4(k)/4},$$
(23)
where the imaginary frequency $`\omega =iu`$ is introduced.
Alternatively, Eq. (23) may be viewed as an interpolation between the exact small- and large-$`q`$ behavior of the Lindhard expression for the frequency-dependent dielectric function. However, we see little point in using such an elaborate expression, since our concern here is to investigate how well local approximations to the dielectric function work in highly non-uniform systems.
In directions parallel to the surface, our approximation Eq. (23) allows fully for the non-diagonality of $`ϵ`$ with respect to the corresponding spatial coordinates, as implied by a Fourier transform with respect to the parallel wave vector $`k`$. However, in directions perpendicular to the surface, our approximation takes $`ϵ`$ to be diagonal in the coordinates $`z`$ and $`z^{}`$. It is thus taken to be local not only in this sense, but in the additional sense that it is a function only of the local density. To compensate, we retain the corresponding component $`q_{}`$ of the wave vector $`q`$ in the right side of Eq. (23) as a parameter, so that everywhere $`q^2=k^2+q_{}^2`$. We thus take $`1/q_{}`$ to be a constant measure of length over which $`ϵ`$ is effectively nonlocal. The dispersion perpendicular to the surface in Eq. (22) is in this way replaced by a parameter that we will fix to some length scale appropriate for the surface.
For physical reasons such a length scale should be associated with intrinsic electron-gas parameters like the screening length or the extent of the correlation hole. There are of course several choices available. It must be kept in mind that we are after long range surface properties in a variety of environments. These properties are determined by various response functions introduced by Feibelman, of which the simplest,
$$d(iu)=\frac{𝑑zzn_{\mathrm{ind}}(iu,z)}{𝑑zn_{\mathrm{ind}}(iu,z)},$$
(24)
is the centroid of induced charge when a uniform electric field is applied perpendicularly to the surface. However, for the van der Waals properties of a planar surface, a related function $`D(iu)`$ as defined by Hult et al.,
$$D(iu)=\frac{ϵ_b(iu)1}{ϵ_b(iu)+1}\frac{ϵ_b(iu)}{ϵ_b(iu)+1}d(iu),$$
(25)
is more important. In particular the $`D`$-function arises in connection with the calculation of van der Waals planes, which are determined not only by the $`D`$ of the surface in question, but also by a response function of the other body. For example, for a surface in the vicinity of an isotropic atom, the van der Waals plane $`Z`$ is given by
$$Z_{\mathrm{vdW}}=\frac{1}{4\pi C_3}_{\mathrm{}}^{\mathrm{}}𝑑u\alpha (iu)D(iu),$$
(26)
where $`\alpha (iu)`$ is the polarizability of the atom and $`C_3`$ is the coefficient of the leading $`1/z^3`$ term in the asymptotic form of the interaction energy. In order for the surface calculation to scale correctly for a wide variety of atoms with different $`\alpha (iu)`$’s, it is obviously important for our approximation to well reproduce $`D(iu)`$. Similarly, for two parallel surfaces labeled A and B, the van der Waals plane for surface A is given by
$$Z_{\mathrm{vdW}}^A=\frac{1}{(4\pi )^2C_2}_{\mathrm{}}^{\mathrm{}}𝑑u\rho ^B(iu)D^A(iu)+\mathrm{\Delta }Z^A,$$
(27)
where $`\rho ^B(iu)`$ is the long wavelength surface response function of surface B, $`\rho ^B(iu)=(ϵ_b^B(iu)1)/(ϵ_b^B(iu)+1)`$, and $`C_2`$ is the coefficient of the leading $`1/z^2`$ term in the asymptotic form of the interaction energy. The main difference in cases such as this one, where two infinite bodies are involved, is that terms involving multiple reflections no longer vanish in the asymptotic limit, and should be explicitly included as indicated by the $`\mathrm{\Delta }Z^A`$ term in Eq. (27). However, it is common experience that in the asymptotic limit multiple reflection terms are typically small enough to treat in perturbation theory, a conclusion that we agree with. Therefore, the first term in Eq. (27) is dominant, and the conclusion that $`D(iu)`$ is the key quantity to obtain accurately remains. Note, however, that in the numerical calculations presented later, the full form of Eq. (27), including all multiple reflections (see Eqs. (8) and (9) of Ref.), is used.
Thus we opt to choose $`q_{}`$ based on the premise that the $`D`$-function should be reproduced accurately. We implement this by choosing $`q_{}`$ so that $`D(0)`$ agrees exactly with full LDA calculations of this quantity, a procedure precisely analogous to that of Ref.. In this way, the parameter $`q_{}`$ may be indirectly determined as a function of the electron gas parameter $`r_s`$. This determination gives $`q_{}1/r_sk_\mathrm{F}`$ as expected from previous arguments.
The constant $`q_{}`$ smoothly limits the response at small $`k`$ and small $`u`$ values. It thus replaces the sharp cutoff used in the earlier scheme. In doing this, we consciously violate the Lifshitz limit, obtaining somewhat smaller vdW coefficients than Ref., since the choice $`q_{}>0`$ affects the $`k0`$ limit of Eq. (23). This also means according to Eqs. (26) and (27) that the van der Waals planes will be predicted to be somewhat too far from the surfaces in this approximation. The $`d`$-function, i.e. Eq. (24) will also be too large, significantly so at small $`u`$.
Far more important, though, is the fact that our simple dielectric function reproduces the dynamic properties of the $`D`$-function very well, as shown in Figures 1 and 2. Thus the overall scaling properties of our theory in a variety of non-uniform van der Waals environments should continue to be correct.
In order to investigate the approximation with regard to key properties, we use as input to our functional a set of self-consistent single-surface jellium densities in the metallic range ($`r_s=25`$). The exact form of the density profiles is however not that important; by construction, a non-local energy is expected to depend less sensitively on the exact form of the density, and results in this paper show that this is indeed the case.
The optimal $`q_{}`$ values found by fitting to the values of $`D(0)`$ of Ref. are shown in Fig. 3 as a function of the electron-gas parameter $`r_s`$, well accounted for by the simple interpolation formula
$$q_{}=0.416e^{0.217r_s}+0.168.$$
(28)
The variation is roughly 50 percent over the whole metallic range, indicating a rather small overall sensitivity.
Table I shows how the quantities $`D(0)`$ and $`Z_{\mathrm{vdW}}`$ varies with $`r_s`$, using the parameterization Eq. (28). In the fourth and fifth columns, the values of the present method for the van der Waals plane $`Z_{\mathrm{vdW}}`$ are compared to those of Ref.. Although the slightly larger coefficients are expected from considerations earlier in this section, the fact that the numbers come out almost *independent* of $`r_s`$ is somewhat of a surprise.
## V Results
To test our approximate DF, we solve Eq. (18), using Eq. (23), and insert the result into Eq. (19) for a known system, consisting of two parallel jellium slabs, separated by a distance $`d`$. Figure 4 shows results in ergs/cm<sup>2</sup> ($`\mathrm{ergs}/\mathrm{cm}^2=0.6423`$ $`\mu `$Ha$`/a_0^2`$) for the non-local correlation-interaction energy per surface area, $`(E_\mathrm{c}^{\mathrm{nl}}(d)E_\mathrm{c}^{\mathrm{nl}}(\mathrm{}))/A`$. The results using the $`r_s=2.07`$ value of Eq. (28) are compared to those of a recent RPA calculation of a slab system of $`r_s=2.07`$ by Dobson and Wang, as well as to those of their approximate DF (IGADEL). The saturation effects are found to be substantial in this small-separation region, and we judge all the proposed density functionals in the figure to give good accounts of the non-local correlation energy. The agreement with IGADEL reflects the inherent similarities between IGADEL and our approximation. The tractability of the latter is however not reflected in the table, but has to be stated here (about a thousand times faster than IGADEL for this particular system, due to the overall lower computational complexity of our DF), together with the claim that this gives great prospects for a tractable future general DF.
The calculation presented in Fig. 4 is performed using a simple superposition of the densities of two separate slabs (obtained from $`d=12`$ data) as input. The difference between the non-local correlation energy for the self-consistent density of the slab system and that obtained by superposition turns out to be very small, as indicated by the stars in Fig. 4.
The minor deviations at $`d`$ close to zero between our results (solid curve) and those of the full RPA calculation (circles) are due to our somewhat crude treatment of the width of the exchange-correlation hole perpendicular to the surface, a small price one has to pay for a tractable DF.
Another important property is the surface correlation energy. This quantity has recently been calculated for a jellium surface within the RPA, an approximation thought to give the long range correlation effects accurately. The long-range part ($`\gamma _{\mathrm{nl}}`$) may be extracted using the data of Kurth and Perdew, by subtracting the LDA contribution to the same quantity. In Table II, the result (column 3) is compared with our approximation (column 2), calculated as the non-local correlation energy given by Eqs. (19) and (23), for single-surface jellium at various values of $`r_s`$, using Eq. (28). We note that the two approximations differ on average by only $`13`$ percent. This is perhaps the strongest indication that indeed the physical approximations made in this paper are robust.
A remarkable fact already indicated in Fig. 4 is that the non-local correlation energy is quite insensitive to the exact form of the density profile. To further test this assumption, we use a linear superposition of two identical self-consistent LDA single-surface densities (SLDA) to calculate the non-local correlation surface energy according to Eq. (21). The result closely follows the column 2 result of Table II, with a mean error of only 3 percent. This observation adds to the accumulated findings supporting the use of superpositions of single-fragment densities in a future general DF.
After these successful tests of the predictive power of the DF, defined by Eqs. (19), (23), and (28), applications to other systems where no other results are available might be done with confidence. Here we present results of an application to two semi-infinite jellia of identical $`r_s`$, having their parallel surfaces a distance $`d`$ apart (Figs. 5 and 6). The input densities are obtained by linear superposition of two self-consistent single-surface LDA densities (SLDA). We stress that our DF is very fast; obtaining a single value for a given density profile and a given separation only takes a few seconds on a typical workstation of today.
Figure 5 shows the variation of the calculated non-local correlation-interaction energy for two semi-infinite jellia of $`r_s=2.07`$ as a function of the separation $`d`$. It illustrates two facts: the significant deviation from the corresponding quantity for two thin jellium slabs (same as in Fig. 4; $`5`$ a.u. wide) and the substantial saturation effects, the latter by comparing with the results of the asymptotic DF formula, applying the dielectric function Eq. (23).
In Fig. 6, we present the normalized non-local correlation-interaction energy $`(E(d)E(\mathrm{}))/2A\gamma _{\mathrm{nl}}`$ in the metallic range, showing how the interaction varies with $`r_s`$. The curves depend only weakly on $`r_s`$ when scaled in this manner.
## VI Concluding remarks
In summary, we have studied the basis for a DF accounting for vdW interactions, by starting in the manner of our previous work but with essential generalizations to small separations between interacting objects. A systematic approach for the construction of such a DF is described, together with a very efficient method to calculate the resulting expressions. In the case of flat surfaces, results for the interaction of two parallel jellium slabs are shown to agree with those of a recent RPA calculation, and we show that input densities can be successfully approximated by a superposition of the electron densities of the interacting fragments. Results for the surface energy of jellium are compared favorably with other studies. As a prediction of the theory, the interaction energy between two parallel jellia is calculated for all separations $`d`$ and in the whole metallic range. The well known asymptotic behavior ($`E1/z^2`$) is obtained for large $`d`$, and as $`d`$ becomes smaller, substantial saturation effects are predicted.
The major significance of these results is the demonstration that such numbers can be calculated accurately at a reduced computational complexity and hence greatly improved speed. We have shown that for a subclass of dielectric functions, the resulting expressions for the non-local correlation energy may be calculated very efficiently, and that even a simple approximation to the dielectric function yields valuable insight and reproduces several physical properties of flat surface and slab models. Furthermore, we have indicated that generalizations to three-dimensional systems are possible, and that the results here suggest such an attempt to be a fruitful one. In this way there should be a basis for applications to numerous physical, chemical, and biological systems, such as vdW bonds in crystals, liquids, adhesion, soft condensed matter (e.g., biomacromolecules, biosurfaces, polymers, and membranes), and scanning-force microscopy.
###### Acknowledgements.
We thank J. Dobson for providing us with several numerical density profiles obtained in Ref. for parallel slabs, and J. Perdew for providing computer code to generate density profiles for isolated jellium surfaces. Work at Rutgers supported in part by NSF Grant DMR 97-08499. Financial support from the Swedish Natural Science Research Council and the Swedish Foundation for Strategical Research through Materials Consortium no. 9 is also acknowledged.
## A Details of the evaluation of the determinants $`\varphi `$ and $`\varphi _0`$
Here we show how ratios of determinants like Eq. (14) can be efficiently evaluated. To eliminate the oscillating sign problem mentioned in the main text, we need to first go to a discrete representation for the operator $`(_z^2k^2+l_k^{}_z)`$ occurring in Eq. (16). Specifically we take $`N`$ points between $`0`$ and $`L`$ (for the full determinant), $`n`$ points between $`0`$ and $`z`$ (for subsystem determinants), so that $`z=(n/N)L`$, with the spacing between points $`h=L/N`$. We use a similar representation for the empty space operator $`(_z^2k^2)`$. Thus, replacing the subsystem determinant $`\varphi (z)`$ by the discretisized $`\varphi _n`$, we have
$$\varphi _n=det\left(\begin{array}{ccccc}a_1& b_1& 0& \mathrm{}& 0\\ c_1& a_2& \mathrm{}& \mathrm{}& \mathrm{}\\ 0& \mathrm{}& \mathrm{}& b_{n2}& 0\\ \mathrm{}& \mathrm{}& c_{n2}& a_{n1}& b_{n1}\\ 0& \mathrm{}& 0& c_{n1}& a_n\end{array}\right),$$
(A1)
where
$`a_n`$ $`=`$ $`2h^2k^2`$ (A2)
$`b_n`$ $`=`$ $`1+{\displaystyle \frac{h}{2}}l_{k,n}^{}`$ (A3)
$`c_n`$ $`=`$ $`1{\displaystyle \frac{h}{2}}l_{k,n+1}^{},`$ (A4)
along with a similar relation for the empty space determinant $`\varphi _{0,n}`$. Being tridiagonal, the determinant can be evaluated in $`𝒪(N)`$ time, as contrasted to the $`𝒪(N^3)`$ time applicable for a general determinant.
The determinant is readily expanded by minors, giving the recursion formula
$$\varphi _{n+1}=a_{n+1}\varphi _nb_nc_n\varphi _{n1},$$
(A5)
with $`\varphi _0=1`$ and $`\varphi _1=a_1`$. It is clear from (A5) and (A4) that $`\varphi _n`$ oscillates in sign from order to order, an oscillation that is cancelled in the ratio Eq. (15) by a similar oscillation in $`\varphi _0`$.
However, the envelope determinant $`(1)^n\varphi _n`$ satisfies a simple differential equation in the continuum limit. To make the exact form of this differential equation identical to Eq. (18), we further scale the envelope determinant $`\stackrel{~}{\varphi }_n`$ through
$$\varphi _n(1)^nS_n\stackrel{~}{\varphi }_n,$$
(A6)
where the scaling function $`S_n\mathrm{\Pi }_{i=1}^nb_i`$ satisfies
$$S_n=b_nS_{n1}$$
(A7)
with $`S_0=1`$. Using (A4), we see that in the continuum $`h0`$ limit, (A7) becomes
$$\frac{dS(z)}{dz}=\frac{1}{2}l_k^{}(z)S(z),$$
(A8)
with the boundary condition $`S(0)=1`$. Equation (A8) has the solution
$$S(z)=\left(\frac{ϵ_k(z)}{ϵ_k(0)}\right)^{\frac{1}{2}}.$$
(A9)
This means that if $`ϵ_k(L)=ϵ_k(0)`$, then the scaling due to $`S`$ has no effect on the final result. Otherwise, the argument of the logarithm in Eq. (19) should be multiplied by $`S(L)`$ obtained from (A9).
The difference equation for $`\stackrel{~}{\varphi }`$ is obtained by use of (A6) and (A7) in (A5), yielding
$$b_{n+1}\stackrel{~}{\varphi }_{n+1}+a_{n+1}\stackrel{~}{\varphi }_n+c_n\stackrel{~}{\varphi }_{n1}=0,$$
(A10)
with
$$\stackrel{~}{\varphi }_0=1$$
(A11)
and
$$\stackrel{~}{\varphi }_1\stackrel{~}{\varphi }_0=(1+a_1/b_1).$$
(A12)
Substitution of (A4) into (A10) yields
$`(\stackrel{~}{\varphi }_{n+1}2\stackrel{~}{\varphi }_n+\stackrel{~}{\varphi }_{n1})h^2k^2\stackrel{~}{\varphi }_n`$ (A13)
$`+{\displaystyle \frac{h}{2}}l_{k,n+1}^{}(\stackrel{~}{\varphi }_{n+1}\stackrel{~}{\varphi }_{n1})=0.`$ (A14)
In the continuum ($`h0`$) limit, this becomes
$$\stackrel{~}{\varphi }^{\prime \prime }(z)k^2\stackrel{~}{\varphi }(z)+l_k^{}(z)\stackrel{~}{\varphi }^{}(z)=0.$$
(A15)
This equation is the same as Eq. (18) in the main text, although there we used the notation $`\stackrel{~}{\varphi }`$ for a particular solution, while in this appendix it represents the general solution.
This general solution to Eq. (18) can be written in the form
$$\stackrel{~}{\varphi }(z)=\alpha \stackrel{~}{\varphi }_\alpha (z)+\beta \stackrel{~}{\varphi }_\beta (z)$$
(A16)
where $`\stackrel{~}{\varphi }_\alpha (0)=1`$ and $`\stackrel{~}{\varphi }_\alpha (L)=0`$, while $`\stackrel{~}{\varphi }_\beta (0)=0`$ and $`\stackrel{~}{\varphi }_\beta (L)=1`$. The boundary condition (A11) implies that $`\alpha =1`$, while (A12) combined with (A4) implies that $`\stackrel{~}{\varphi }^{}(0)=1/h`$ as $`h`$ approaches zero. This means that $`\stackrel{~}{\varphi }_\alpha `$ can be neglected, since $`\beta `$ will be of order $`1/h`$. Specifically we have $`\stackrel{~}{\varphi }^{}(0)\beta \varphi _\beta ^{}(0)=1/h`$, so that
$$\stackrel{~}{\varphi }(L)\beta =\frac{1}{h\stackrel{~}{\varphi }_\beta ^{}(0)}.$$
(A17)
The large coefficient ($`1/h`$) cancels out in the ratio of (A17) and the analogous free-space expression, so we may write
$$\frac{\varphi (L)}{\varphi _0(L)}=\frac{\stackrel{~}{\varphi }(L)}{\stackrel{~}{\varphi }_0(L)}S(L)=\frac{\stackrel{~}{\varphi }_{0,\beta }^{}(0)}{\stackrel{~}{\varphi }_\beta ^{}(0)}S(L),$$
(A18)
where the continuum version of (A6) was used to obtain the first identity, after noting that the oscillating signs of the numerator and denominator cancel in the ratio, and that the free space value of $`S`$ is unity. Finally substituting this into Eq. (17), we obtain, using (A9),
$$E_c^{\mathrm{nl}}=_0^{\mathrm{}}\frac{du}{2\pi }\frac{d^2k}{(2\pi )^2}\mathrm{log}\frac{\stackrel{~}{\varphi }_\beta ^{}(0)\sqrt{ϵ_k(L)}}{\stackrel{~}{\varphi }_{0,\beta }^{}(0)\sqrt{ϵ_k(0)}}.$$
(A19)
This is our most general result, which reduces to Eq. (19) when $`ϵ_k(0)=ϵ_k(L)`$. Note that in the main text, we used the notation $`\stackrel{~}{\varphi }`$ and $`\stackrel{~}{\varphi }_0`$ for the particular solutions that are called $`\stackrel{~}{\varphi }_\beta `$ and $`\stackrel{~}{\varphi }_{0,\beta }`$ here.
## B Details on the Lifshitz limit
Let $`ϵ_k(z)=1`$ in between two surfaces a distance $`d`$ apart, and let $`ϵ_k(z)=ϵ_b`$ in the bulk of each surface. Furthermore, let $`d`$ be large enough so we can assume sharp boundaries between the three regions. Then, the interaction energy $`(E_c^{\mathrm{nl}}(d)E_c^{\mathrm{nl}}(\mathrm{}))/A`$ may be calculated exactly. Note that although the interaction energy may be calculated this way, the constant component contributing to the surface energy may not, but needs a much more involved calculation.
The matching problem becomes solving Eq. (18) for $`\varphi `$ and $`\varphi _0`$ in the regions left bulk (lb), middle region (m), and right bulk (rb). Let $`\varphi _0=e^{kx}`$, with the origin lying on the boundary between the left and middle regions. Now we want to find the $`\varphi `$ that approaches zero far into the left bulk, and $`e^{kx}`$ far into the right bulk. In the left bulk, we must have
$$\varphi _{\mathrm{lb}}=ae^{kx},$$
(B1)
$`a`$ being an arbitrary constant. Note that the wanted quantity $`\varphi _0^{}(0)/\varphi ^{}(0)`$ equals $`1/a`$. On the boundary between the left bulk and the middle region, $`\varphi `$ must be continuous, and $`\varphi ^{}`$ must have a discontinuous step with the size of $`ϵ_b`$, reflecting the fact that the displacement field $`ϵ_b\varphi ^{}`$ must also be continuous. The solution in the middle region now becomes
$$\varphi _\mathrm{m}=a(\mathrm{cosh}(kz)+ϵ_b\mathrm{sinh}(kz)).$$
(B2)
The same matching conditions apply between the middle region and the right bulk, although here, we only need to consider the solution growing to the right, since we want to match $`\varphi `$ to the value of $`\varphi _0`$ infinitely far into the right bulk. Obeying the matching conditions means the solution in the right bulk becomes
$$\varphi _{\mathrm{rb}}^{\mathrm{growing}}=\frac{\varphi _\mathrm{m}(d)k+\varphi _\mathrm{m}^{}(d)/ϵ_b}{2k}e^{k(xd)}.$$
(B3)
Equating $`\varphi _0`$ with (B3) determines the coefficient a, and the solution becomes
$$\varphi _0^{}(0)/\varphi ^{}(0)=\frac{\varphi _\mathrm{m}(d)k+\varphi _\mathrm{m}^{}(d)/ϵ_b}{2kae^{kd}}=\frac{1\rho ^2e^{2kd}}{1\rho ^2},$$
(B4)
with $`\rho =(ϵ_b1)/(ϵ_b+1)`$. It is clear from (B4) that the constant contribution to $`\mathrm{log}(\varphi _0^{}(0)/\varphi ^{}(0))`$ and hence to the surface energy is given by $`\mathrm{log}(1\rho ^2)`$. As discussed earlier in this appendix, that is not the correct constant and should be excluded, yielding Eq. (20). |
warning/0001/hep-th0001107.html | ar5iv | text | # Randall-Sundrum Choice in the Brane World
## I Introduction
Recently there has been much interest in the Randall-Sundrum brane-world. The key idea of this model is that our universe may be a brane embedded in higher dimensional space. A concrete model is a single 3-brane embedded in five-dimensional anti-de Sitter space ($`\mathrm{AdS}_5`$). Randall and Sundrum have shown that the longitudinal part ($`h_{\mu \nu }`$) of the metric fluctuations satisfy the Schrödinger-like equation with an attractive delta-function. As a result, the massless zero mode which describes the localized gravity on the brane was found. Furthermore, the massive modes lead to the correction to the Newtonian potential as of $`V(r)=G_N\frac{m_1m_2}{r}(1+\frac{1}{r^2k^2})`$.
However, we point out that this has been done with the RS choice (a four-dimensional transverse-tracefree gauge). It seems that this choice is so restrictive that the RS model can describe the tensor fluctuation only. Furthermore, in order to have the well-defined theory on the brane, one has to consider the transverse parts of $`h_{5\mu },h_{55}`$. More recently, Ivanov and Volovich found that the equation for $`h_{55}`$ takes the Schrödinger-like equation with a repulsive delta-function. But their linearized equation is not correct.
In the massless and massive cases, $`h_{55}`$ is a four-dimensional scalar and $`h_{5\mu }`$ is a four-dimensional vector. Hence it is not natural to set these fields to be zero, as is shown in the RS choice. At the first sight, the RS choice does not seem to be consistent with the massive states. This is because $`h_{55}`$ and $`h_{5\mu }`$ belong to the physical fields and these cannot be gauged away because the general covariance is broken. In both cases we choose the other gauge such as the de Donder gauge (a five-dimensional transverse-tracefree gauge) instead of the RS choice in the beginning.
In this paper, we find the correct linearized equation including $`h_{5\mu },h_{55}`$. We point out the validity of the RS choice in describing the massless states as well as massive ones in the RS brane model. Also we discuss its connection to the stability of the RS solution.
## II Perturbation Analysis
We start with the Einstein equation with the bulk cosmological constant $`\mathrm{\Lambda }`$ and the brane tension $`\stackrel{~}{\sigma }`$
$$R_{MN}\frac{1}{2}g_{MN}R=\mathrm{\Lambda }g_{MN}+\sigma \sqrt{g_{55}}g_{\mu \nu }\delta _M^\mu \delta _N^\nu ,$$
(1)
which is derived from the action
$$I=\frac{1}{2}d^5x\sqrt{g_5}(R+2\mathrm{\Lambda })+\stackrel{~}{\sigma }d^4x\sqrt{g_4}.$$
(2)
The RS solution is given by
$$\overline{g}_{MN}=H^2\eta _{MN}$$
(3)
with $`H=k|z|+1`$ and $`\eta _{MN}=\mathrm{diag}(+)`$. Further $`\mathrm{\Lambda }=6k^2`$ and $`\sigma =\stackrel{~}{\sigma }\delta (z)`$ with $`\stackrel{~}{\sigma }=6k`$. Here the capital indices $`M,N,\mathrm{}`$ are split into $`\mu ,\nu ,\mathrm{}`$ (four-dimensions: $`x^\mu `$) and $`5(x^5=z)`$.
After the conformal transformation of $`g_{MN}=\mathrm{\Omega }^2\stackrel{~}{g}_{MN}`$ with $`\mathrm{\Omega }=H^1`$, let us introduce the perturbation
$$\stackrel{~}{g}_{MN}=\eta _{MN}+h_{MN}.$$
(4)
Its linearized equation for Eq.(1) takes the form
$`\mathrm{}h_{MN}+3{\displaystyle \frac{_KH}{H}}\eta ^{KL}\left(_Nh_{KM}+_Mh_{KN}_Kh_{MN}\right)`$ (5)
$`\left({\displaystyle \frac{2\mathrm{\Lambda }+2\sigma }{H^2}}\right)h_{55}\eta _{MN}{\displaystyle \frac{2\sigma }{H^2}}\left\{h_{MN}\left(h_{\mu \nu }+{\displaystyle \frac{h_{55}}{2}}\eta _{\mu \nu }\right)\delta _M^\mu \delta _N^\nu \right\}=0.`$ (6)
Ivanov and Volovich in the version 2 of ref. have missed the second line of Eq.(6). They in the version 3 have missed all of $`\sigma `$-dependent terms but have included the $`\mathrm{\Lambda }`$-dependent term. This appears because the terms without $``$ arising from the LHS of Eq.(1) cannot be cancelled against those$`\left(H^2\left[\mathrm{\Lambda }h_{MN}+\sigma \left(h_{\mu \nu }+\frac{h_{55}}{2}\eta _{\mu \nu }\right)\delta _M^\mu \delta _N^\nu \right]\right)`$ from the RHS of Eq.(1). This line vanishes if $`h_{MN}`$ reduces to $`h_{\mu \nu }`$ with $`h_{55}=h_{5\mu }=0`$.
Here we use the de Donder gauge
$$^Mh_{MN}=0,h_P^P=0.$$
(7)
This means that
$$h_\mu ^\mu =h_{55},^\mu h_{\mu 5}=_5h_{55},^\mu h_{\mu \nu }=_5h_{5\nu }.$$
(8)
From Eq.(6) we obtain three equations,
$`\left(\mathrm{}{\displaystyle \frac{12k^2}{H^2}}3f_5\right)h_{55}=0,`$ (9)
$`\left(\mathrm{}{\displaystyle \frac{12k}{H^2}}\delta (z)\right)h_{5\mu }3f_\mu h_{55}=0,`$ (10)
$`\left(\mathrm{}+3f_5\right)h_{\mu \nu }3f\left(_\mu h_{5\nu }+_\nu h_{5\mu }\right)+{\displaystyle \frac{12}{H^2}}\left(k^2{\displaystyle \frac{k}{2}}\delta (z)\right)h_{55}\eta _{\mu \nu }=0`$ (11)
with $`f=_5H/H`$. Taking the trace of (11) and comparing it with Eq.(9), one finds that $`h_{55}`$ should vanish. In deriving $`h_{55}=0`$, we use the de Donder gauge in (8). With $`h_{55}=0`$, Eq.(10) becomes a decoupled equation. In analyzing the perturbations, if one finds a decouple one, then one should solve it first. And then one has to check its consistency with the remaining equation (11). In order to solve Eq.(10) first, we introduce the separation of variables as
$$h_{5\mu }(x,z)=C_\mu (z)\psi _5(x).$$
(12)
Then Eq.(10) takes the form:
$`\left(\mathrm{}_4+m_5^2\right)\psi _5(x)=0,`$ (13)
$`C_\mu ^{\prime \prime }(z)+\left\{{\displaystyle \frac{12k}{H^2}}\delta (z)+m_5^2\right\}C_\mu (z)=0`$ (14)
with the gauge condition of $`C_\mu (z)^\mu \psi (x)=0`$. Here the prime() means the differentiation with respect to its argument. Now let us solve Eq.(14) first. This is exactly the case of ref.. The solution must satisfy the equation $`C_\mu ^{\prime \prime }(z)+m_5^2C_\mu (z)=0`$ at everywhere, except $`z=0`$. And then we assume its plane wave solution as
$$C_\mu (z)=A_\mu e^{im_5|z|};C_\mu ^{z>0}(z)=A_\mu e^{im_5z},C_\mu ^{z<0}(z)=\stackrel{~}{A}_\mu e^{im_5z}.$$
(15)
We note that this solution keeps the reflection symmetry of the RS solution as $`C_\mu (z^{})=C_\mu (z)`$, under $`z^{}z`$. The coefficients in front are the same $`A_\mu =\stackrel{~}{A}_\mu `$ because of the continuity of the wave function. The derivative of $`C_\mu (z)`$ is no longer continuous because of the presence of the delta-function. That is, one has
$$\frac{C_\mu }{z}|_{z=0^+}\frac{C_\mu }{z}|_{z=0^{}}=12kA_\mu ,$$
(16)
which leads to
$$im_5=6k$$
(17)
This admits the tachyonic mass of $`C_\mu (z)`$ as $`m_5^2=36k^2<0`$. In other words, the normalizable bound-state solution to Eq.(14) is allowed if its energy($`m_5^2`$) is negative. As a check, $`C_\mu ^t(z)=A_\mu e^{6k|z|}`$ satisfies
$$C_{\mu }^{t}{}_{}{}^{\prime \prime }(z)+12k\delta (z)C_\mu ^t(z)+m_5^2(=36k^2)C_\mu ^t(z)=0.$$
(18)
But it remains to check whether this solution is or not consistent with Eq.(11). Acting $`^\mu `$ on Eq.(11) and using Eqs.(8) and (10), one gets the condition
$$\left(\delta (z)C_\mu \right)^{}3\mathrm{s}\mathrm{g}\mathrm{n}(z)\delta (z)C_\mu =0.$$
(19)
We note that $`\mathrm{sgn}(z)\delta (z)`$ is not well defined at $`z=0`$ and thus one requires
$$C_\mu (0)=0.$$
(20)
An alternative solution which satisfies Eqs.(14) and (20) is the plane wave as Eq.(15) but $`C_\mu ^p(0)=0`$,
$$C_\mu ^p(z)=A_\mu \mathrm{sin}m_5|z|.$$
(21)
At this stage, we remind the reader that our background is AdS<sub>5</sub> with $`\delta (z)`$-source. This means that the solution to the linearized equations should carry at least the parameter “ $`k`$” because the size of AdS<sub>5</sub> box is $`1/k`$ approximately and the brane tension is $`\stackrel{~}{\sigma }=6k`$. However this plane wave solution misses “$`k`$”. This seems to be a solution for 5D Minkowski but not for AdS<sub>5</sub> background. This is so because, due to the condition (20) this does not account for the presence of the brane at $`z=0`$($`12k\delta (z)C_\mu `$-term in Eq.(14)) appropriately. On the other hand, if $`C_\mu (0)0`$, $`\delta (z)C_\mu (z)`$ can be taken into account(as in our tachyonic solution $`C_\mu ^t`$). That is, there is no solution which satisfies both Eqs.(10) and (11). Hence we are in a dilemma if $`h_{5\mu }`$ is truely a massive vector in the RS brane world.
Consequently, the tachyonic solution $`C_\mu ^t(z)`$ is not a physical one because it is incompatible with the tensor equation (11). As it stands, the presence of this solution says that $`h_{5\mu }`$ should be rejected to have a well-defined theory. Fortunately the consistency with Eq.(11) leads to $`h_{5\mu }=0`$ on the whole space $`z`$ as in the RS choice. Furthermore, the analysis for the massless case ($`m_5=0`$) in Eq.(14) leads to $`A_\mu =0`$. This implies that there is no massless vector state on the brane. Hence it is obvious that $`h_{5\mu }`$ should not be a propagating vector in the RS background. From now on we set $`h_{5\mu }=0`$.
## III Massless States
These states correspond to $`_5h_{MN}=0`$. Before we proceed, we are willing to count the number of independent components of the graviton $`h_{MN}`$. For $`D=5`$ dimensions, a symmetric tensor field $`h_{MN}`$ has $`5(5+1)/2=15`$ independent components, some of which can be eliminated by the gauge conditions (7). This is $`(5+1)`$. Further, after choosing the guage (7), there exists a residual gauge degrees of freedom as
$$h_{MN}^{}=h_{MN}_M\xi _N_N\xi _M.$$
(22)
Notice that $`h_{MN}^{}`$ satisfy the de Donder gauge (7) provided that
$$_M\xi ^M=0,\mathrm{}\xi ^M=0.$$
(23)
Thus $`(51)`$ are eliminated by our freedom. Hence the number of massless degrees of freedom in $`D=5`$ is
$$\frac{56}{2}(5+1)(51)=5.$$
(24)
In order to see how $`5`$ is composed, let us consider the conventional Kaluza-Klein (KK) model. This corresponds to $`\mathrm{}h_{MN}=0`$. Its massless bound state ($`_5h_{MN}=0`$) in $`D=5`$ dimensions can be described by
$$\mathrm{}h_{MN}=J_{MN}$$
(25)
with the external source $`J_{MN}`$. The general covariance of massless case can be represented as a source conservation law of $`^NJ_{MN}=0`$ with $`J_M^M=0`$. In this case we choose a Lorentz frame as
$$_1=_4,_2=_3=_5=0.$$
(26)
In this frame, the effective interaction reduces to the positive-definite form
$$_{\mathrm{massless}}^{\mathrm{KK}}=\frac{1}{4}h_{MN}J^{MN}=\frac{1}{4}\underset{\lambda =2}{\overset{2}{}}J_\lambda \frac{1}{_1^2_4^2}J_\lambda ,$$
(27)
where $`\lambda `$ refers to $`\mathrm{O}(2)`$ helicity and
$`J_{\pm 2}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(J_{22}J_{33})\pm iJ_{23},`$ (28)
$`J_{\pm 1}`$ $`=`$ $`J_{52}\pm iJ_{53},`$ (29)
$`J_0`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}J_{55}.`$ (30)
The terms in (27) with (28), (29), and (30) describe the exchanges of spin-2 graviton, spin-1 photon, and spin-0 scalar. Here we have 2 components (spin-2), 2 (spin-1), and 1 (spin-0) and summing up these leads to 5 in (24). According to the stability analysis, it is stable if each pole in (27) is positive-definite. Hence the KK model is classically stable.
Now let us consider the same issue using the RS choice as follows:
$$h_{55}=0,h_\mu ^\mu =0,h_{5\mu }=0,^\mu h_{\mu \nu }=0.$$
(31)
These eliminate $`10`$ in $`15`$. Furthere we point out that there exists a $`D=4`$ residual gauge as
$$h_{\mu \nu }^{}=h_{\mu \nu }_\mu \xi _\nu _\nu \xi _\mu $$
(32)
with
$$_\mu \xi ^\mu =0,\mathrm{}_4\xi ^\mu =0.$$
(33)
This eliminates $`3`$ degrees of freedom (DOFs). Hence we have left 2 DOFs ($`=15103`$) in the RS choice. This is appropriate for describing the graviton $`h_{\mu \nu }`$ only. Under this gauge, one finds from Eq.(11)
$$\mathrm{}_4h_{\mu \nu }=J_{\mu \nu }$$
(34)
with the source relations
$$J_{55}=0,J_\mu ^\mu =0,J_{5\mu }=0,^\mu J_{\mu \nu }=0.$$
(35)
Using these, Eq.(27) reduces
$$_{\mathrm{massless}}^{\mathrm{RS}}=\frac{1}{4}\left[J_2\frac{1}{_1^2_4^2}J_2+J_2\frac{1}{_1^2_4^2}J_2\right]$$
(36)
with $`J_{\pm 2}=J_{22}\pm iJ_{23}`$. As a result, the RS choice can describe the massless spin-2 modes of $`h_{\pm 2}=h_{22}\pm ih_{23}`$ only as it can do best. One cannot find the vector and scalar fields. Three modes of $`h_{\pm 1}=h_{52}\pm ih_{53}`$ and $`h_0=\sqrt{3/2}h_{55}`$ are missed, in comparison with the conventional KK model. Furthermore, it is shown that the RS background is stable because $`_{\mathrm{massless}}^{\mathrm{RS}}`$ in (36) is positive-definite.
## IV Massive states
In this case we start with the de Donder gauge in (7). But it turns out that the set of perturbation equations (9)-(11) become
$$(\mathrm{}+3f_5)h_{\mu \nu }=0,h_{55}=h_{5\mu }=0,$$
(37)
which corresponds to the RS massive case. At this stage, it is convenient to introduce the new variables $`h_{\mu \nu }=H^{3/2}(z)\widehat{h}_{\mu \nu }(x,z)=H^{3/2}(z)\psi _h(z)\widehat{\widehat{h}}_{\mu \nu }(x)`$. $`\widehat{\widehat{h}}_{\mu \nu }`$ corresponds to the canonical form of $`h_{\mu \nu }`$. Then one finds
$$\left[\frac{\mathrm{}_4}{2}+[\frac{_5^2}{2}+V(z)]\right]\widehat{h}_{\mu \nu }=0,$$
(38)
with
$$V(z)=\frac{15k^2}{8H^2}\frac{3k}{2H}\delta (z).$$
(39)
Considering the equation of $`(\mathrm{}_4+m_h^2)\widehat{\widehat{h}}_{\mu \nu }=J_{\mu \nu }`$ with $`J_{55}=J_{5\mu }=0`$, we find the source conservation law as follows:
$$^\mu J_{\mu \nu }=0J_\mu ^\mu =0.$$
(40)
Further the mass $`m_h^2`$ is determined by the equation
$$\left[\frac{1}{2}_5^2+V(z)\right]\psi _h(z)=\frac{1}{2}m_h^2\psi _h(z).$$
(41)
It was shown that $`V(r)`$ guarentees $`m_h^20`$. This implies that these are no normalizable negative energy graviton modes. In this case we choose a massive Lorentz frame in which
$$_1=_5\mathrm{and}_i=0\mathrm{for}i=2,3,4.$$
(42)
It follows that, in the neighborhood of the pole, the effective interaction reduces to
$$_{\mathrm{massive}}^{\mathrm{RS}}=\frac{1}{4}J_{ij}\frac{1}{_1^2+m_h^2}J_{ij},$$
(43)
where $`J_{ij}`$ is a symmetric traceless tensor in three dimensions.
One finds a massive tensor with 5 DOFs because $`J_{ij}`$ has 5 ($`=34/21`$) components. It is interesting to ask how we can interpret this DOFs. This is clear from the fact that in the massive case the global symmetry of spacetime is spontaneously broken. The gauge parameters $`\xi _\mu (x,z)`$ and $`\xi _5(x,z)`$ in (22) are associated to spontaneously broken generators. The $`h_{\mu \nu }`$ with 2 DOFs acquire mass $`m_h^2`$ by eating 2 DOFs of $`h_{5\mu }`$-vector and 1 DOF of $`h_{55}`$ scalar. Thus one finds a pure spin-2 massive particle with 5 DOFs. Explicitly, these are $`h_{23},h_{24},h_{34}`$, and other two satisfying $`h_{22}+h_{33}+h_{44}=0`$. All these have positive-definite norm states. Hence all of the massive states in the RS model are classically stable.
## V Discussion
We study the validity of the RS choice in the RS model. For this purpose we start with the de Donder gauge. Using the RS choice for the massive case, one finds 5 DOFs in $`h_{\mu \nu }`$. These all turn out to be the physical massive modes. Hence there remains no residual gauge symmetry. In the massless case, we have three gauge degrees of freedom upon choosing the RS one. This corresponds to a residual gauge degrees of freedom. Using these, we can always find the massless spin-2 with 2 DOFs; for example, see Ref.. Hence we always have a localized gravity in a 3-brane.
For the stability of the RS solution of $`ds_{\mathrm{RS}}^2=H^2\eta _{MN}dx^Mdx^N`$, we find that $`_{\mathrm{massless}}^{\mathrm{RS}}`$ has positive norm states for the RS choice, and also $`_{\mathrm{massive}}^{\mathrm{RS}}`$ is positive-definite with the RS choice. This means that the RS choice is regarded as a unique one to establish the stability of the RS solution.
Finally, we comment on the Schrödinger-like equation (10) with $`h_{55}=0`$. It may imply that the RS model is not classically stable because it has a tachyonic mass. However, considering its consistency with Eq.(11), this puzzle can be resolved. If the RS model makes sense, $`h_{5\mu }=0`$.
## Acknowledgments
We would like to thank Volovich for sending his paper to us. This work was supported by the Brain Korea 21 Program, Ministry of Education, Project No. D-0025. |
warning/0001/cond-mat0001186.html | ar5iv | text | # Nonsingular vortices in (s+d)-wave superconductors
## Abstract
The structure of a single flux line in (s+d)-wave superconductors has been analyzed within the Ginzburg-Landau (GL) model generalized for two order parameter components. The fourfold symmetric singular vortex solution is shown to be unstable in a certain range of the GL parameters with respect to the mutual shift of s- and d- wave unit vortices. The resulting nonsingular vortex structure is studied both analytically and numerically.
Recently the distinctive characteristics of vortices in unconventional superconductors which can be described by the phenomenological Ginzburg-Landau (GL) theory with a multicomponent order parameter (OP) are of great interest in connection with the investigations of the mixed state structure in high-$`T_c`$ and heavy fermion compounds which are strong candidates for unconventional superconductivity. A single flux line in such systems is known to contain a set of unit vortices of different OP components . One can indentify two possible types of flux lines: (i) singular vortices (which have at least one point where the superconducting gap is zero) and (ii) nonsingular vortices (where the gap is nonzero everywhere in the vortex core). In this paper we focus on the case of $`s+d_{x^2y^2}`$-wave pairing (which can be relevant to the case of high-$`T_c`$ superconductors) and consider the range of GL functional parameters where the singular vortices (studied in ) become unstable according to the scenario analogous to the one proposed in for heavy fermion compounds. The goal of this paper is to analyse the detailed structure of resulting nonsingular vortices using both numerical and analytical methods.
We start with the GL free energy functional generalized for two components of the OP $`\mathrm{\Psi }_d`$ and $`\mathrm{\Psi }_s`$ corresponding to the $`d_{x^2y^2}`$-wave and s-wave pairing, respectively :
$$F=\{a_d|\mathrm{\Psi }_d|^2+a_s|\mathrm{\Psi }_s|^2+\frac{b_d}{2}|\mathrm{\Psi }_d|^4+\frac{b_s}{2}|\mathrm{\Psi }_s|^4$$
$$+\alpha |\mathrm{\Psi }_d|^2|\mathrm{\Psi }_s|^2+\frac{\beta }{2}(\mathrm{\Psi }_d^2\mathrm{\Psi }_s^2+\mathrm{\Psi }_d^2\mathrm{\Psi }_s^2)$$
$$+K_s|𝚷\mathrm{\Psi }_s|^2+K_d|𝚷\mathrm{\Psi }_d|^2+\gamma [(\mathrm{\Pi }_x^{}\mathrm{\Psi }_s^{}\mathrm{\Pi }_x\mathrm{\Psi }_d$$
$$\mathrm{\Pi }_y^{}\mathrm{\Psi }_s^{}\mathrm{\Pi }_y\mathrm{\Psi }_d)+c.c.]+\frac{𝐇^2}{8\pi }\}d𝐫,$$
(1)
where $`𝚷=i\frac{2\pi }{\mathrm{\Phi }_0}𝐀`$, $`𝐇=curl𝐀`$, $`𝐫=(x,y)`$, $`a_s=\alpha _s(TT_{cs})`$, $`a_d=\alpha _d(TT_{cd})`$ (we assume $`T_{cs}<T_{cd}`$ and the magnetic field is applied parallel to the $`c`$-axis). If the GL parameters are chosen so that there is an one-component homogeneous state ($`𝐇=0`$), then for $`𝐇0`$ the subdominant OP can appear in the vortex core regions with the inhomogeneous dominant OP component either due to the mixing gradient terms ($`\gamma 0`$) or due to the instability of the inhomogeneous state against the formation of the subdominant OP nucleus . Let us consider the latter mechanism (which is obviously responsible for the formation of nonsingular vortices) for the simplest case $`\gamma =0`$, $`\beta =0`$ and continue with the analysis of the linearised GL equation of $`\mathrm{\Psi }_s`$:
$$K_s^2\mathrm{\Psi }_s+\alpha |\mathrm{\Psi }_d|^2\mathrm{\Psi }_s=a_s\mathrm{\Psi }_s,$$
(2)
where $`\mathrm{\Psi }_d`$ describes the vortex solution within the conventional single-component GL theory, also we neglect $`𝐀`$. The ”lowest energy eigenstate” of the Schrödinger equation (2) defines the temperature $`T^{}`$ of the phase transition into nonsingular vortex state. Using the following approximation $`|\mathrm{\Psi }_d|=|\mathrm{\Psi }_d|_{\mathrm{}}r/\sqrt{r^2+2\xi _{d}^{}{}_{}{}^{2}}`$ (here $`|\mathrm{\Psi }_d|_{\mathrm{}}=\sqrt{|a_d|/b_d}`$, $`\xi _d=\sqrt{K_d/|a_d|}`$) for the cases $`K_sK_da_s/a_d`$ and $`\alpha b_dK_s/K_d`$, we obtain $`a(T^{})=\sqrt{2K_s\alpha /(K_db_d)}`$ and $`a(T^{})\alpha /b_d`$, respectively (here $`a(T)=a_s(T)/|a_d(T)|`$). We propose the following extrapolation for the phase transition curve on $`\alpha T`$ plane (Fig.1), which separates the regions with singular and nonsingular vortices: $`\alpha =b_da(T^{})(1a(T^{})K_d/(2K_s))`$.
This expression is valid for the above two limiting cases and for the particular case ($`\alpha _s=1.25\alpha _d`$, $`T=0`$, $`T_{cs}=0.8T_{cd}`$, $`b_s=b_d`$, $`\beta =0`$, $`K_s=K_d`$, $`\gamma =0`$) studied in . The analysis of the structure of nonsingular vortices has been carried out using numerical calculations based on the time-dependent GL theory. It was found out that for $`\gamma 0`$ the nonsingular vortices with broken fourfold symmetry (see Fig.2) are stable below
the phase transition curve (solid line in Fig.1), whereas above the curve it is the singular fourfold symmetric vortices that are energetically favourable.
In conclusion, we have considered the structure of nonsingular vortices in $`s+d_{x^2y^2}`$ superconductor and found the range of parameters where these vortices exist. The nonsingular vortices in such materials have neither the fourfold symmetry, nor the normal core region. This can lead to some interesting effects (e.g. nontrivial dynamics, moderate pinning effects etc.). |
warning/0001/physics0001070.html | ar5iv | text | # Interaction of a vortex ring with the free surface of ideal fluid
## 1 Introduction
The study of interaction between vortex structures in a fluid and the free surface is important both from practical and theoretical points of view. In general, a detailed investigation of this problem is very hard. Even the theories of potential surface waves and the dynamics of vortices in an infinite space taken separately still have a lot of unsolved fundamental problems on their own. Only the consideration of significantly simplified models can help us to understand the processes which take place in the combined system.
In many cases it is possible to neglect the compressibility of the fluid as well as the energy dissipation. Therefore the model of ideal homogeneous incompressible fluid is very useful for hydrodynamics. Because of the conservative nature of this model the application of the well developed apparatus of Hamiltonian dynamics becomes possible . An example of effective use of the Hamiltonian formalism in hydrodynamics is the introduction of canonical variables for investigations of potential flows of perfect fluids with a free boundary. V.E.Zakharov showed at the end of the sixties that the surface shape $`z=\eta (x,y,t)`$ and the value of the velocity potential $`\psi (x,y,t)`$ on the surface can be considered as generalized coordinate and momentum, respectively.
It is important to note that a variational formulation of Hamiltonian dynamics in many cases allows to obtain good finite-dimensional approximations which reflect the main features of the behavior of the original system. There are several possibilities for a parameterization of non-potential flows of perfect fluid by some variables with dynamics determined by a variational principle. All of them are based on the conservation of the topological characteristics of vortex lines in ideal fluid flows which follows from the freezing-in of the vorticity field $`𝛀(𝐫,t)=\text{curl}𝐯(𝐫,t)`$. In particular, this is the representation of the vorticity by Clebsch canonical variables $`\lambda `$ and $`\mu `$
$$𝛀(𝐫,t)=[\lambda \times \mu ]$$
However, the Clebsch representation can only describe flows with a trivial topology (see, e.g., ). It cannot describe flows with linked vortex lines. Besides, the variables $`\lambda `$ and $`\mu `$ are not suitable for the study of localized vortex structures like vortex filaments. In such cases it is more convenient to use the parameterization of vorticity in terms of vortex lines and consider the motion of these lines ,, even if the global definition of canonically conjugated variables is impossible due to topological reasons.
This approach is used in the present article to describe the interaction of deep (or small) vortex rings of almost ideal shape in the perfect fluid with the free surface. In the case under consideration the main interaction of the vortex rings with the surface can be described as the dipole-dipole interaction between ”point” vortex rings and their ”images”. Moving rings interact with the surface waves, leading to radiation due to the Cherenkov effect. Deep rings disturb the surface weakly, so the influence of the surface can be taken into account as some small corrections in the equations of motion for the parameters of the rings.
In Sec.2 we discuss briefly general properties of vortex line dynamics, which follow from the freezing-in of the vorticity field. In Sec.3 possible simplifications of the model are made and the point ring approximation is introduced. In Sec.4 the interaction of the ring with its image is considered. In Sec.5 we calculate the Fourier-components of Cherenkov surface waves radiated by a moving vortex ring and determine the non-conservative corrections caused by the interaction with the surface for the vortex ring equations of motion.
## 2 Vortex lines motion in perfect fluid
It is a well known fact that the freezing-in of the vorticity lines follows from the Euler equation for ideal fluid motion
$$𝛀_t=\text{curl}[𝐯\times 𝛀],𝐯=\text{curl}^1𝛀$$
Vortex lines are transported by the flow ,,. They do not appear or disappear, neither they intersect one another in the process of motion. This property of perfect fluid flows is general for all Hamiltonian systems of the hydrodynamic type. For simplicity, let us consider temporally the incompressible fluid without free surface in infinite space. The dynamics of the system is specified by a basic Lagrangian $`L[𝐯]`$, which is a functional of the solenoidal velocity field. The relations between the velocity $`𝐯`$, the generalized vorticity $`𝛀`$, the basic Lagrangian $`L[𝐯]`$ and the Hamiltonian $`[𝛀]`$ are the following <sup>2</sup><sup>2</sup>2 For the ordinary ideal hydrodynamics in infinite space the basic Lagrangian is
$$L_{Euler}[𝐯]=\frac{𝐯^2}{2}𝑑𝐫𝛀=\text{curl}𝐯$$
The Hamiltonian in this case coincides with the kinetic energy of the fluid and in terms of the vorticity field it reads
$$_{Euler}[𝛀]=1/2𝛀\mathrm{\Delta }^1𝛀𝑑𝐫=\frac{1}{8\pi }\frac{𝛀(𝐫_1)𝛀(𝐫_2)}{|𝐫_1𝐫_2|}𝑑𝐫_1𝑑𝐫_2$$
where $`\mathrm{\Delta }^1`$ is the inverse Laplace operator. Another example is the basic Lagrangian of Electron Magneto-hydrodynamics which takes into account the magnetic field created by the current of electron fluid through the motionless ion fluid.
$$L_{EMHD}[𝐯]=\frac{1}{2}𝐯(1\mathrm{\Delta }^1)𝐯𝑑𝐫𝛀=\text{curl}(1\mathrm{\Delta }^1)𝐯$$
$$_{EMHD}[𝛀]=\frac{1}{2}𝛀(1\mathrm{\Delta })^1𝛀𝑑𝐫=\frac{1}{8\pi }\frac{e^{|𝐫_1𝐫_2|}}{|𝐫_1𝐫_2|}𝛀(𝐫_1)𝛀(𝐫_2)𝑑𝐫_1𝑑𝐫_2$$
The second example shows that the relation between the velocity and the vorticity can be more complex than in usual hydrodynamics.
$$𝛀=\text{curl}\left(\frac{\delta L}{\delta 𝐯}\right)𝐯=𝐯[𝛀]$$
(1)
$$[𝛀]=\left(𝐯\left(\frac{\delta L}{\delta 𝐯}\right)d^3𝐫L[𝐯]\right)|_{𝐯=𝐯[𝛀]}$$
(2)
$$𝐯=\text{curl}\left(\frac{\delta }{\delta 𝛀}\right)$$
(3)
and the equation of motion for the generalized vorticity is
$$𝛀_t=\text{curl}[\text{curl}(\delta /\delta 𝛀)\times 𝛀]$$
(4)
This equation corresponds to the transport of frozen-in vortex lines by the velocity field. In this process all topological invariants of the vorticity field are conserved. The conservation of the topology can be expressed by the following relation
$$𝛀(𝐫,t)=\delta (𝐫𝐑(𝐚,t))(𝛀_0(𝐚)_𝐚)𝐑(𝐚,t)𝑑𝐚=\frac{(𝛀_0(𝐚)_𝐚)𝐑(𝐚,t)}{\text{det}𝐑/𝐚}|_{𝐚=𝐚(𝐫,t)}$$
(5)
where the mapping $`𝐑(𝐚,t)`$ describes the deformation of lines of some initial solenoidal field $`𝛀_0(𝐫)`$. Here $`𝐚(𝐫,t)`$ is the inverse mapping with respect to $`𝐑(𝐚,t)`$. The direction of the vector $`𝐛`$
$$𝐛(𝐚,t)=(𝛀_0(𝐚)_𝐚)𝐑(𝐚,t)$$
(6)
coincides with the direction of the vorticity field at the point $`𝐑(𝐚,t)`$. The equation of motion for the mapping $`𝐑(𝐚,t)`$ can be obtained with the help of the relation
$$𝛀_t(𝐫,t)=\text{curl}_𝐫\delta (𝐫𝐑(𝐚,t))[𝐑_t(𝐚,t)\times 𝐛(𝐚,t)]𝑑𝐚,$$
(7)
which immediately follows from Eq.(5). The substitution of Eq.(7) into the equation of motion (4) gives
$$\text{curl}_𝐫\left(\frac{𝐛(𝐚,t)\times [𝐑_t(𝐚,t)𝐯(𝐑,t)]}{\text{det}𝐑/𝐚}\right)=0$$
One can solve this equation by eliminating the $`\text{curl}_𝐫`$ operator. Using the general relationship between variational derivatives of some functional $`F[𝛀]`$
$$\left[𝐛\times \text{curl}\left(\frac{\delta F}{\delta 𝛀(𝐑)}\right)\right]=\frac{\delta F}{\delta 𝐑(𝐚)}|_{𝛀_0}$$
(8)
it is possible to represent the equation of motion for $`𝐑(𝐚,t)`$ as follows
$$\left[(𝛀_0(𝐚)_𝐚)𝐑(𝐚)\times 𝐑_t(𝐚)\right]=\frac{\delta [𝛀[𝐑]]}{\delta 𝐑(𝐚)}|_{𝛀_0}.$$
(9)
It is not difficult to check now that the dynamics of the vorticity field with topological properties defined by $`𝛀_0`$ in the infinite space is equivalent to the requirement of an extremum of the action ($`\delta S=\delta _{𝛀_0}𝑑t=0`$) where the Lagrangian is
$$_{𝛀_0}=\frac{1}{3}\left(\left[𝐑_t(𝐚)\times 𝐑(𝐚)\right](𝛀_0(𝐚)_𝐚)𝐑(𝐚)\right)𝑑𝐚[𝛀[𝐑]].$$
(10)
In the simplest case, when all vortex lines are closed it is possible to choose new curvilinear coordinates $`\nu _1,\nu _2,\xi `$ in $`𝐚`$-space such that Eq.(5) can be written in a simple form
$$𝛀(𝐫,t)=_𝒩d^2\nu \delta (𝐫𝐑(\nu ,\xi ,t))𝐑_\xi 𝑑\xi .$$
(11)
Here $`\nu `$ is the label of a line lying on a fixed two-dimensional manifold $`𝒩`$, and $`\xi `$ is some parameter along the line. It is clear that there is a gauge freedom in the definition of $`\nu `$ and $`\xi `$. This freedom is connected with the possibility of changing the longitudinal parameter $`\xi =\xi (\stackrel{~}{\xi },\nu ,t)`$ and also with the relabeling of $`\nu `$
$$\nu =\nu (\stackrel{~}{\nu },t),\frac{(\nu _1,\nu _2)}{(\stackrel{~}{\nu }_1,\stackrel{~}{\nu }_2)}=1.$$
(12)
Now we again consider the ordinary perfect fluid with a free surface. To describe the flow entirely it is sufficient to specify the vorticity field $`𝛀(𝐫,t)`$ and the motion of the free surface. Thus, we can use the shape $`𝐑(\nu ,\xi ,t)`$ of the vortex lines as a new dynamic object instead of $`𝛀(𝐫,t)`$. It is important to note that in the presence of the free surface the equations of motion for $`𝐑(\nu ,\xi ,t)`$ follow from a variational principle as in the case of infinite space. It has been shown that the Lagrangian for a perfect fluid, with vortices in its bulk and with a free surface, can be written in the form
$$=\frac{1}{3}_𝒩d^2\nu ([𝐑_t\times 𝐑]𝐑_\xi )𝑑\xi +\mathrm{\Psi }\eta _t𝑑𝐫_{}[𝐑,\mathrm{\Psi },\eta ].$$
(13)
The functions $`\mathrm{\Psi }(𝐫_{},t)`$ and $`\eta (𝐫_{},t)`$ are the surface degrees of freedom for the system. $`\mathrm{\Psi }`$ is the boundary value of total velocity potential, which includes the part from vortices inside the fluid, and $`\eta `$ is the deviation of the surface from the horizontal plane. This formulation supposes that vortex lines do not intersect the surface anywhere. In the present paper only this case is considered.
The Hamiltonian $``$ in Eq.(13) is nothing else than the total energy of the system expressed in terms of $`[𝐑,\mathrm{\Psi },\eta ]`$.
Variation with respect to $`𝐑(\nu ,\xi ,t)`$ of the action defined by the Lagrangian (13) gives the equation of motion for vortex lines in the form
$$[𝐑_\xi \times 𝐑_t]=\frac{\delta [𝛀[𝐑],\mathrm{\Psi },\eta ]}{\delta 𝐑}.$$
(14)
This equation determines only the transversal component of $`𝐑_t`$ which coincides with the transversal component of the actual solenoidal velocity field. The possibility of solving Eq.(14) with respect to the time derivative $`𝐑_t`$ is closely connected with the special gauge invariant nature of the $`[𝐑]`$ dependence which results in
$$\frac{\delta }{\delta 𝐑}𝐑_\xi 0.$$
The tangential component of $`𝐑_t`$ with respect to vorticity direction can be taken arbitrary. This property is in accordance with the longitudinal gauge freedom. The vorticity dynamics does not depend on the choice of the tangential component.
Generally speaking, only the local introduction of canonical variables for curve dynamics is possible. For instance, a piece of the curve can be parameterized by one of the three of Cartesian coordinates
$$𝐑=(X(z,t),Y(z,t),z)$$
In this case the functions $`X(z,t)`$ and $`Y(z,t)`$ are canonically conjugated variables. Another example is the parameterization in cylindrical coordinates, where variables $`Z(\theta ,t)`$ and $`(1/2)R^2(\theta ,t)`$ are canonically conjugated.
Curves with complicated topological properties need a general gauge free description by means of a parameter $`\xi `$.
It should be mentioned for clarity that the conservation of all vortex tube volumes, reflecting the incompressibility of the fluid, is not the constraint in this formalism. It is a consequence of the symmetry of the Lagrangian (13) with respect to the relabeling (12) $`\nu \stackrel{~}{\nu }`$ . Volume conservation follows from that symmetry in accordance with Noether’s theorem. To prove this statement, we should consider such subset of relabelings which forms a one-parameter group of transformations of the dynamical variables. For small values of the group parameter, $`\tau `$, the transformations are determined by a function of two variables $`T(\nu _1,\nu _2)`$ (with zero value on the boundary $`𝒩`$) so that
$$𝐑(\nu _1,\nu _2,\xi )𝐑_T^\tau (\nu _1,\nu _2,\xi )=𝐑(\nu _1\tau \frac{T}{\nu _2}+O(\tau ^2),\nu _2+\tau \frac{T}{\nu _1}+O(\tau ^2),\xi )$$
(15)
Due to Noether’s theorem, the following quantity is an integral of motion
$$I_T=_𝒩d^2\nu \frac{\delta }{\delta 𝐑_t}\frac{𝐑_T^\tau }{\tau }|_{\tau =0}d\xi =\frac{1}{3}_𝒩d^2\nu [𝐑\times 𝐑_\xi ](𝐑_2T_1𝐑_1T_2)𝑑\xi $$
After simple integrations in parts the last expression takes the form
$$I_T=_𝒩d^2\nu T(\nu _1,\nu _2)([𝐑_1\times 𝐑_2]𝐑_\xi )𝑑\xi =_𝒩T(\nu _1,\nu _2)𝒱(\nu _1,\nu _2,t)d^2\nu $$
(16)
where $`𝒱(\nu _1,\nu _2,t)d^2\nu `$ is the volume of an infinitely thin vortex tube with cross-section $`d^2\nu `$. It is obvious that actually the function $`𝒱`$ doesn’t depend on time $`t`$ because the function $`T(\nu _1,\nu _2)`$ is arbitrary <sup>3</sup><sup>3</sup>3 If vortex lines are not closed but form a family of enclosed tori then the relabeling freedom is less rich. In that case one can obtain by the similar way the conservation laws for volumes inside closed vortex surfaces. Noether’s theorem gives integrals of motion which depend on an arbitrary function of one variable $`S(\zeta )`$, where $`\zeta `$ is the label of the tori. .
## 3 Point ring approximation
In general case an analysis of the dynamics defined by the Lagrangian (13) is too much complicated. We do not even have the exact expression for the Hamiltonian $`[𝐑,\mathrm{\Psi },\eta ]`$ because it needs the explicit knowledge of the solution of the Laplace equation with a boundary value assigned on a non-flat surface. Another reason is the very high nonlinearity of the problem.
In this paper we consider some limits where it is possible to simplify the system significantly. Namely, we will suppose that the vorticity is concentrated in several very thin vortex rings of almost ideal shape. For a solitary ring the perfect shape is stable for a wide range of vorticity distributions through the cross-section. This shape provides an extremum of the energy for given values of the volumes of vortex tubes and for a fixed momentum of the ring. As already mentioned, volume conservation follows from Noether’s theorem. Therefore some of these quantities (those of which are produced by the subset of commuting transformations) can be considered as canonical momenta. Corresponding cyclical coordinates describe the relabeling (12) of the line markers, which doesn’t change the vorticity field. Actually these degrees of freedom take into account a rotation around the central line of the tube. This line represents the mean shape of the ring and we are interested in how it behaves in time. For our analysis we don’t need the explicit values of cyclical coordinates, but only the conserved volumes as parameters in the Lagrangian.
A possible situation is when a typical time of the interaction with the surface and with other rings is much larger then the largest period of oscillations corresponding to deviations of the ring shape from perfect one. Under this condition, excitations of all (non-cyclical) internal degrees of freedom are small during all the time, and a variational anzats completely disregarding them reflects the behavior of the system adequately. The circulations
$$\mathrm{\Gamma }_n=_{𝒩_n}d^2\nu $$
of the velocity for each ring don’t depend on time. A perfect ring is described by the coordinate $`𝐑_n`$ of the center and by the vector $`𝐏_n=\mathrm{\Gamma }_n𝐒_n`$, where $`𝐒_n`$ is an oriented area of the ring. We use in this work the Cartesian system of coordinates $`(x,y,z)`$, so that the vertical coordinate is $`z`$, and the unperturbed surface is at $`z=0`$. The corresponding components of the vectors $`𝐑_n`$ and $`𝐏_n`$ are
$$𝐑_n=(X_n,Y_n,Z_n),𝐏_n=(P_{xn},P_{yn},P_{zn})$$
It is easy to verify that the vectors $`𝐏_n`$ are canonically conjugated momenta for the coordinates $`𝐑_n`$. To verify that we can parameterize the shape of each vortex line in the following manner
$$𝐑(\xi ,t)=\underset{m=M}{\overset{M}{}}𝐫_m(t)e^{im\xi },𝐫_m=\overline{𝐫}_m$$
(17)
Here $`𝐫_m(t)`$ are complex vectors. Substituting this into the first term of the Lagrangian (13) gives
$$\frac{1}{3}([𝐑_t\times 𝐑]𝐑_\xi )𝑑\xi =2\pi i\dot{𝐫}_0([𝐫_1\times 𝐫_1]+2[𝐫_2\times 𝐫_2]+\mathrm{})+$$
$$+\frac{d\{\mathrm{}\}}{dt}+2\pi i\dot{𝐫}_1[𝐫_1\times 𝐫_2]2\pi i\dot{𝐫}_1[𝐫_1\times 𝐫_2]+\mathrm{}$$
(18)
If we neglect the internal degrees of freedom which describe deviations of the ring from the ideal shape
$$(𝐫_1)^2=(𝐫_1)^2=0,𝐫_2=𝐫_2=0,\mathrm{}$$
then the previous statement about canonically conjugated variables becomes obvious:
$$𝐑_n=𝐫_{0n},𝐏_n=2\pi \mathrm{\Gamma }_ni[𝐫_{1n}\times 𝐫_{1n}]$$
(19)
Such an approximation is valid only in the limit when sizes of rings are small in comparison with the distances to the surface and the distances between different rings
$$\sqrt{\frac{P_n}{\mathrm{\Gamma }_n}}|Z_n|,|𝐑_n𝐑_l|,ln.$$
(20)
These conditions are necessary for ensuring that the excitations of all internal degrees of freedom are small. Obviously, this is not true when a ring approaches the surface. In that case one should take into account also the internal degrees of freedom for the vortex lines.
The inequalities (20) also imply that vortex rings in the limit under consideration are similar to point magnet dipoles. This analogy is useful for calculation of the Hamiltonian for interacting rings. In the main approximation we may restrict the analysis by taking into account the dipole-dipole interaction only.
It should be mentioned that in some papers (see e.g. and references in that book) the discrete variables identical to $`𝐑_n`$ and $`𝐏_n`$ are derived in a different way and referred as the vortex magnetization variables.
In the expression for the Hamiltonian, several simplifications can be made. Let us recall that for each moment of time it is possible to decompose the velocity field into two components
$$𝐯=𝐕_0+\varphi .$$
(21)
Here the field $`𝐕_0`$ satisfies the following conditions
$$(𝐕_0)=0,\text{curl}𝐕_0=𝛀,(𝐧𝐕_0)|_{z=\eta }=0.$$
The boundary value of the surface wave potential $`\varphi (𝐫)`$ is $`\psi (𝐫_{})`$. In accordance with these conditions the kinetic energy is decomposed into two parts and the Hamiltonian of the fluid takes the form
$$=\frac{1}{2}_{z<\eta }𝐕_0^2d^3𝐫+\frac{1}{2}\psi (\varphi d𝐒)+\frac{g}{2}\eta ^2𝑑𝐫_{}$$
(22)
The last term in this expression is the potential energy of the fluid in the gravitational field. If all vortex rings are far away from the surface then its deviation from the horizontal plane is small
$$|\eta |1,|\eta ||Z_n|$$
(23)
Therefore in the main approximation the energy of dipoles interaction with the surface can be described with the help of so called ”images”. The images are vortex rings with parameters
$$\mathrm{\Gamma }_n,𝐑_n^{}=(X_n,Y_n,Z_n),𝐏_n^{}=(P_{xn},P_{yn},P_{zn})$$
(24)
The kinetic energy for the system of point rings and their images is the sum of the self-energies of rings and the dipole-dipole interaction between them. The expression for the kinetic energy of small amplitude surface waves employs the operator $`\widehat{k}`$ which multiplies Fourier-components of a function by the absolute value $`k`$ of a two-dimensional wave vector $`𝐤`$. So the real Hamiltonian $``$ is approximately equal to the simplified Hamiltonian $`\stackrel{~}{}`$
$$\stackrel{~}{}=\frac{1}{2}(\psi \widehat{k}\psi +g\eta ^2)𝑑𝐫_{}+\underset{n}{}_n(P_n)+$$
$$+\frac{1}{8\pi }\underset{ln}{}\frac{3((𝐑_n𝐑_l)𝐏_n)((𝐑_n𝐑_l)𝐏_l)|𝐑_n𝐑_l|^2(𝐏_n𝐏_l)}{|𝐑_n𝐑_l|^5}+$$
$$+\frac{1}{8\pi }\underset{ln}{}\frac{3((𝐑_n𝐑_l^{})𝐏_n)((𝐑_n𝐑_l^{})𝐏_l^{})|𝐑_n𝐑_l^{}|^2(𝐏_n𝐏_l^{})}{|𝐑_n𝐑_l^{}|^5}$$
(25)
With the logarithmic accuracy the self-energy of a thin vortex ring is given by the expression
$$_n(P_n)\frac{\mathrm{\Gamma }_n^2}{2}\sqrt{\frac{P_n}{\pi \mathrm{\Gamma }_n}}\mathrm{ln}\left(\frac{(P_n/\mathrm{\Gamma }_n)^{3/4}}{A_n^{1/2}}\right)$$
(26)
where the small constant $`A_n`$ is proportional to the conserved volume of the vortex tube forming the ring. This expression can easily be derived if we take into account that the main contribution to the energy is from the vicinity of the tube where the velocity field is approximately the same as near a straight vortex tube. The logarithmic integral should then be taken between the limits from the thickness of the tube to the radius of the ring.
In the relation $`\mathrm{\Psi }=\mathrm{\Phi }_0+\psi `$ the potential $`\mathrm{\Phi }_0`$ is approximately equal to the potential created on the flat surface by the dipoles and their images
$$\mathrm{\Phi }_0(𝐫_{})\mathrm{\Phi }(𝐫_{})=\frac{1}{2\pi }\underset{n}{}\frac{(𝐏_n(𝐫_{}𝐑_n))}{|𝐫_{}𝐑_n|^3}$$
(27)
In this way we arrive at the following simplified system describing the interaction of point vortex rings with the free surface
$$\stackrel{~}{}=\underset{n}{}\dot{𝐑}_n𝐏_n+\dot{\eta }(\psi +\mathrm{\Phi })d^2𝐫_{}\stackrel{~}{}[\{𝐑_n,𝐏_n\},\eta ,\psi ]$$
(28)
It should be noted that due to the condition (20) the maximum value of the velocity $`V_0`$ on the surface is much less then the typical velocities of the vortex rings
$$\frac{P_n}{Z_n^3}\frac{\mathrm{\Gamma }_n^{3/2}}{P_n^{1/2}}$$
Therefore the term $`V_0^2/2`$ in the Bernoulli equation
$$\mathrm{\Psi }_t+V_0^2/2+g\eta +\text{small corrections}=0$$
is small in comparison with the term $`\mathrm{\Psi }_t`$. The Lagrangian (28) is in accordance with this fact because it does not take into account terms like $`(1/2)V_0^2\eta d^2𝐫_{}`$ in the Hamiltonian expansion.
## 4 Interaction of the vortex ring with its image
Now let us for simplicity consider the case of a single ring. It is shown in the next section, that for a sufficiently deep ring the interaction with its image is much stronger than the interaction with the surface waves. So it is interesting to examine the motion of the ring neglecting the surface deviation. In this case we have the integrable Hamiltonian for the system with two degrees of freedom
$$H=\frac{1}{64\pi }\left(\alpha (P_x^2+P_z^2)^{1/4}\frac{2P_z^2+P_x^2}{|Z|^3}\right),Z<0$$
(29)
where $`\alpha \text{const}`$. The system has integrals of motion
$$P_x=p=const,H=E=const$$
so it is useful to consider the level lines of the energy function in the left $`(Z,P_z)`$-half-plane taking $`P_x`$ as the parameter (see the Figure).
One can distinguish three regions of qualitatively different behavior of the ring in that part of this half-plane where our approximation is valid (see Eq.(20)). In the upper region the phase trajectories come from infinitely large negative $`Z`$ where they have a finite positive value of $`P_z`$. In the process of motion $`P_z`$ increases. This behavior corresponds to the case when the ring approaches the surface. Due to the symmetry of the Hamiltonian (29) there is a symmetric lower region, where the vortex ring moves away from the surface. And there is the middle region, where $`P_z`$ changes the sign from negative to positive at a finite value of $`Z`$. This is the region of the finite motion.
In all three cases the track of the vortex ring bends toward the surface, i.e. the ring is ”attracted” by the surface.
## 5 Cherenkov interaction of a vortex ring with surface waves
When the ring is not very far from the surface and not very slow, the interaction with the surface waves becomes significant. Let us consider the effect of Cherenkov radiation of surface waves by a vortex ring which moves from the infinity to the surface. This case is the most definite from the viewpoint of initial conditions choice. We suppose that the deviation of the free surface from the horizontal plane $`z=0`$ is zero at $`t\mathrm{}`$, and we are interested in the asymptotic behavior of fields $`\eta `$ and $`\psi `$ at large negative $`t`$. In this situation we can neglect the interaction of the ring with its image in comparison with the self-energy and concentrate our attention on interaction with surface waves only.
The ring moves in the $`(x,z)`$-plane with an almost constant velocity. In the main approximation the position $`𝐑`$ of the vortex ring is given by the relations
$$𝐑𝐂t,𝐂=𝐂(𝐏)=\frac{}{𝐏}=(C_x,0,C_z)\frac{𝐏}{P^{3/2}},$$
(30)
$$C_x>0,C_z>0,t<0.$$
The equations of motion for the Fourier-components of $`\eta `$ and $`\psi `$ follow from the Lagrangian (28)
$$\dot{\eta }_𝐤=k\psi _𝐤,\dot{\psi }_𝐤+g\eta _𝐤=\dot{\mathrm{\Phi }}_𝐤$$
(31)
Eliminating $`\eta _𝐤`$ we obtain an equation for $`\psi _𝐤`$
$$\ddot{\psi }_𝐤+gk\psi _𝐤=\ddot{\mathrm{\Phi }}_𝐤$$
(32)
where $`\mathrm{\Phi }_𝐤`$ is the Fourier-transform of the function $`\mathrm{\Phi }(𝐫_{})`$. Simple calculations give
$$\mathrm{\Phi }_𝐤=\frac{e^{ik_xX}}{2\pi }\frac{P_zZP_xx}{\sqrt{(x^2+y^2+Z^2)^3}}e^{i(k_xx+k_yy)}𝑑x𝑑y=$$
$$=\frac{e^{ik_xX}}{2\pi }\left(P_zD(k|Z|)+i\frac{P_x}{|Z|}\frac{}{k_x}D(k|Z|)\right)$$
(33)
where
$$D(q)=\frac{e^{iq\alpha }d\alpha d\beta }{\sqrt{(\alpha ^2+\beta ^2+1)^3}}=2\pi e^{|q|}$$
(34)
Finally, we have for $`\mathrm{\Phi }_𝐤`$
$$\mathrm{\Phi }_𝐤=\left(\frac{iP_xk_x}{k}P_z\right)e^{k|Z|ik_xX}=\left(\frac{iP_xk_x}{k}P_z\right)e^{t(kC_zik_xC_x)}$$
(35)
Due to the exponential time behavior of $`\mathrm{\Phi }_𝐤(t)`$ it is easy to obtain the expressions for $`\psi _𝐤(t)`$ and $`\eta _𝐤(t)`$. Introducing the definition
$$\lambda _𝐤=kC_zik_xC_x$$
(36)
we can represent the answer in the following form
$$\psi _𝐤(t)=\frac{\left(\frac{iP_xk_x}{k}P_z\right)\lambda _𝐤^2}{gk+\lambda _𝐤^2}e^{\lambda _𝐤t}=\left(\frac{P}{Ck}\right)\frac{\lambda _𝐤^3}{gk+\lambda _𝐤^2}e^{\lambda _𝐤t}$$
(37)
$$\eta _𝐤(t)=\left(\frac{P}{C}\right)\frac{\lambda _𝐤^2}{gk+\lambda _𝐤^2}e^{\lambda _𝐤t}$$
(38)
The radiated surface waves influence the motion of the vortex ring. The terms produced by the field $`\eta _𝐤(t)`$ in the equations of motion for the ring come from the part $`\dot{\eta }\mathrm{\Phi }d^2𝐫_{}`$ in the Lagrangian (28). Using Eq.(35) for the Fourier-transform of $`\mathrm{\Phi }`$ we can represent these terms as follows
$$\delta \dot{X}=\frac{d^2𝐤}{(2\pi )^2}\dot{\eta }_𝐤\frac{ik_x}{k}e^{kZ+ik_xX}$$
(39)
$$\delta \dot{Z}=\frac{d^2𝐤}{(2\pi )^2}\dot{\eta }_𝐤e^{kZ+ik_xX}$$
(40)
$$\delta \dot{P}_x=\frac{d^2𝐤}{(2\pi )^2}\dot{\eta }_𝐤(ik_x)\left(P_z+\frac{iP_xk_x}{k}\right)e^{kZ+ik_xX}$$
(41)
$$\delta \dot{P}_z=\frac{d^2𝐤}{(2\pi )^2}\dot{\eta }_𝐤k\left(P_z+\frac{iP_xk_x}{k}\right)e^{kZ+ik_xX}$$
(42)
We can use Eq.(38) to obtain the nonconservative corrections for time derivatives of the ring parameters from these expressions. It is convenient to write down these corrections in the autonomic form
$$\delta \dot{X}=\left(\frac{P}{C}\right)\frac{d^2𝐤}{(2\pi )^2}\left(\frac{ik_x}{k}\right)\frac{(kC_zik_xC_x)^3}{gk+(kC_zik_xC_x)^2}e^{2k|Z|}$$
(43)
$$\delta \dot{Z}=\left(\frac{P}{C}\right)\frac{d^2𝐤}{(2\pi )^2}\frac{(kC_zik_xC_x)^3}{gk+(kC_zik_xC_x)^2}e^{2k|Z|}$$
(44)
$$\delta \dot{P}_x=\left(\frac{P}{C}\right)^2\frac{d^2𝐤}{(2\pi )^2}\left(\frac{ik_x}{k}\right)\frac{(kC_zik_xC_x)^2(C_z^2k^2+C_x^2k_x^2)}{gk+(kC_zik_xC_x)^2}e^{2k|Z|}$$
(45)
$$\delta \dot{P}_z=\left(\frac{P}{C}\right)^2\frac{d^2𝐤}{(2\pi )^2}\frac{(kC_zik_xC_x)^2(C_z^2k^2+C_x^2k_x^2)}{gk+(kC_zik_xC_x)^2}e^{2k|Z|}$$
(46)
where $`C_x`$ and $`C_z`$ can be understood as explicit functions of $`𝐏`$ defined by the dependence $`𝐂(𝐏)=/𝐏`$. More exact definition of $`C_x`$ and $`C_z`$ as $`\dot{X}`$ and $`\dot{Z}`$ is not necessary.
To analyze the above integrals let us first perform there the integration over the angle $`\phi `$ in $`𝐤`$-space. It is convenient to use the theory of contour integrals in the complex plane of variable $`w=\mathrm{cos}\phi `$. The contour $`\gamma `$ of integration in our case goes clockwise just around the cut which is from $`1`$ to $`+1`$. We define the sign of the square root $`R(w)=\sqrt{1w^2}`$ so that its values are positive on the top side of the cut and negative on the bottom side. After introducing the quantities
$$a=\frac{C_z}{C_x},\omega _𝐤^2=gk,b_𝐤=\frac{\omega _𝐤}{C_xk}=\frac{1}{C_x}\sqrt{\frac{g}{k}}$$
(47)
we have to use the following relations
$$I_1(a,b)\underset{\gamma }{}\frac{dw}{\sqrt{1w^2}}\frac{w(w+ia)^3}{b^2(w+ia)^2}=$$
$$=\pi (1+2b^2)+\pi i\left(\frac{(bia)b^2}{\sqrt{1(bia)^2}}\frac{(b+ia)b^2}{\sqrt{1(bia)^2}}\right)$$
(48)
$$I_2(a,b)i\underset{\gamma }{}\frac{dw}{\sqrt{1w^2}}\frac{(w+ia)^3}{b^2(w+ia)^2}=$$
$$=\pi \left(2a+\frac{b^2}{\sqrt{1(bia)^2}}+\frac{b^2}{\sqrt{1(bia)^2}}\right)$$
(49)
$$J_1(a,b)i\underset{\gamma }{}\frac{dw}{\sqrt{1w^2}}\frac{w(w+ia)^2(w^2+a^2)}{b^2(w+ia)^2}=$$
$$=4\pi ab^2+\pi \left(\frac{b(bia)(a^2+(bia)^2)}{\sqrt{1(bia)^2}}+\frac{b(b+ia)(a^2+(b+ia)^2)}{\sqrt{1(bia)^2}}\right)$$
(50)
$$J_2(a,b)\underset{\gamma }{}\frac{dw}{\sqrt{1w^2}}\frac{(w+ia)^2(w^2+a^2)}{b^2(w+ia)^2}=$$
$$=2\pi (a^2+b^2+1/2)\pi i\left(\frac{b(a^2+(bia)^2)}{\sqrt{1(bia)^2}}\frac{b(a^2+(b+ia)^2)}{\sqrt{1(bia)^2}}\right)$$
(51)
where the sign of the complex square root should be taken in accordance with the previous choice. It can easily be seen that the integrals $`I_2`$ and $`J_1`$ have resonance structure at $`a1`$ and $`|b|<1`$. This is the Cherenkov effect itself. Now the expressions (43-46) take the form
$$\delta \dot{X}=\frac{P_x}{(2\pi )^2}\underset{0}{\overset{+\mathrm{}}{}}I_1(a,b_𝐤)k^2e^{2k|Z|}𝑑k=\frac{P_x}{(2\pi )^2}\left(\frac{g}{C_x^2}\right)^3F_1(a,\frac{2g|Z|}{C_x^2})$$
(52)
$$\delta \dot{Z}=\frac{P_x}{(2\pi )^2}\underset{0}{\overset{+\mathrm{}}{}}I_2(a,b_𝐤)k^2e^{2k|Z|}𝑑k=\frac{P_x}{(2\pi )^2}\left(\frac{g}{C_x^2}\right)^3F_2(a,\frac{2g|Z|}{C_x^2})$$
(53)
$$\delta \dot{P}_x=\frac{P_x^2}{(2\pi )^2}\underset{0}{\overset{+\mathrm{}}{}}J_1(a,b_𝐤)k^3e^{2k|Z|}𝑑k=\frac{P_x^2}{(2\pi )^2}\left(\frac{g}{C_x^2}\right)^4G_1(a,\frac{2g|Z|}{C_x^2})$$
(54)
$$\delta \dot{P}_z=\frac{P_x^2}{(2\pi )^2}\underset{0}{\overset{+\mathrm{}}{}}J_2(a,b_𝐤)k^3e^{2k|Z|}𝑑k=\frac{P_x^2}{(2\pi )^2}\left(\frac{g}{C_x^2}\right)^4G_2(a,\frac{2g|Z|}{C_x^2})$$
(55)
Here the functions $`F_1(a,Q)..G_2(a,Q)`$ are defined by the integrals
$$F_1(a,Q)=\underset{0}{\overset{+\mathrm{}}{}}I_1(a,\frac{1}{\sqrt{\xi }})\mathrm{exp}\left(Q\xi \right)\xi ^2𝑑\xi $$
(56)
$$F_2(a,Q)=\underset{0}{\overset{+\mathrm{}}{}}I_2(a,\frac{1}{\sqrt{\xi }})\mathrm{exp}\left(Q\xi \right)\xi ^2𝑑\xi $$
(57)
$$G_1(a,Q)=\underset{0}{\overset{+\mathrm{}}{}}J_1(a,\frac{1}{\sqrt{\xi }})\mathrm{exp}\left(Q\xi \right)\xi ^3𝑑\xi $$
(58)
$$G_2(a,Q)=\underset{0}{\overset{+\mathrm{}}{}}J_2(a,\frac{1}{\sqrt{\xi }})\mathrm{exp}\left(Q\xi \right)\xi ^3𝑑\xi $$
(59)
and $`Q=2g|Z|/C_x^2`$ is a dimensionless quantity <sup>4</sup><sup>4</sup>4 If we consider a fluid with surface tension $`\sigma `$, then two parameters appear: $`Q`$ and $`T=g\sigma /C_x^4`$. In that case one should substitute $`b_𝐤\sqrt{1/\xi +T\xi }`$ as the second argument of the functions $`I_1,I_2,J_1,J_2`$ in the integrals (56-59) . The Cherenkov effect is most clear when the motion of the ring is almost horizontal. In this case $`a+0`$, and it is convenient to rewrite these integrals without use of complex functions
$$F_1(+0,Q)=\pi \underset{0}{\overset{+\mathrm{}}{}}\left(\xi ^2+2\xi \right)\mathrm{exp}\left(Q\xi \right)𝑑\xi 2\pi \underset{0}{\overset{1}{}}\frac{\xi d\xi }{\sqrt{1\xi }}\mathrm{exp}\left(Q\xi \right)$$
(60)
$$F_2(+0,Q)=G_1(+0,Q)=2\pi \underset{1}{\overset{+\mathrm{}}{}}\frac{\xi ^{3/2}d\xi }{\sqrt{\xi 1}}\mathrm{exp}\left(Q\xi \right)$$
(61)
$$G_2(+0,Q)=\pi \underset{0}{\overset{+\mathrm{}}{}}\left(\xi ^3+2\xi ^2\right)\mathrm{exp}\left(Q\xi \right)𝑑\xi +2\pi \underset{0}{\overset{1}{}}\frac{\xi ^2d\xi }{\sqrt{1\xi }}\mathrm{exp}\left(Q\xi \right)$$
(62)
Here the square root is the usual positive defined real function. We see that only resonant wave-numbers contribute to the functions $`F_2`$ and $`G_1`$, while $`F_1`$ and $`G_2`$ are determined also by small values of $`\xi `$ which correspond to the large scale surface deviation co-moving with the ring. So the effect of the Cherenkov radiation on the vortex ring motion is the most distinct in the equations for $`\dot{Z}`$ and $`\dot{P}_x`$. Especially it is important for $`P_x`$ because the radiation of surface waves is the only reason for change of this quantity in the frame of our approximation.
The typical values of $`Q`$ are large in practical situations. In this limit asymptotic values of the integrals above are
$$F_1(+0,Q)\frac{9\pi }{2Q^4},G_2(+0,Q)\frac{18\pi }{Q^5}$$
$$F_2(+0,Q)=G_1(+0,Q)2\pi \sqrt{\pi }\frac{\mathrm{exp}(Q)}{\sqrt{Q}}$$
and
$$\delta \dot{X}\frac{9}{64\pi }\frac{P}{|Z|^3}\frac{1}{Q},\delta \dot{Z}\frac{1}{16\sqrt{\pi }}\frac{P}{|Z|^3}Q^{2+1/2}\mathrm{exp}(Q),$$
$$\delta \dot{P}_x\frac{1}{32\sqrt{\pi }}\frac{P^2}{|Z|^4}Q^{3+1/2}\mathrm{exp}(Q),\delta \dot{P}_z+\frac{9}{32\pi }\frac{P^2}{|Z|^4}\frac{1}{Q}.$$
It follows from these expressions that the interaction with the surface waves is small in comparison with the interaction between ring and its image, if $`Q1`$. The corresponding small factors are $`1/Q`$ for $`X`$ and $`P_z`$, and $`Q^{2+1/2}\mathrm{exp}(Q)`$ for $`Z`$. As against the flat boundary, now $`P_x`$ is not conserved. It decreases exponentially slowly and this is the main effect of Cherenkov radiation.
We see also that the interaction with waves turns the vector $`𝐏`$ towards the surface which results in a more fast boundary approach by the ring track.
## 6 Conclusions and acknowledgments
In this paper we have derived the simplified Lagrangian for the description of the motion of deep vortex rings under free surface of perfect fluid. We have analyzed the integrable dynamics corresponding to the pure interaction of the single point vortex ring with its image. It was found that there are three types of qualitatively different behaviour of the ring. The interaction of the ring with the surface has an attractive character in all three regimes. The Fourier-components of radiated Cherenkov waves were calculated for the case when the vortex ring comes from infinity and has both horizontal and vertical components of the velocity. The non-conservative corrections to the equations of motion of the ring, due to Cherenkov radiation, were derived. Due to these corrections the track of the ring bends towards the surface faster then in the case of flat surface. For simplicity, all calculations in Sec.5 were performed for a single ring. The generalization for the case of many rings is straightforward.
The author thanks professor J.J. Rasmussen for his attention to this work and for helpful suggestions. This work was supported by the INTAS (grant No. 96-0413), the Russian Foundation for Basic Research (grant No. 97-01-00093), and the Landau Postdoc Scholarship (KFA, Forschungszentrum, Juelich, Germany). |
warning/0001/astro-ph0001422.html | ar5iv | text | # Gas Stripping of Dwarf Galaxies in Clusters of Galaxies
## 1. Introduction
Explaining the morphology-density relation (Dressler 1980), that is the higher early-type galaxy fraction in clusters of galaxies in contrast to the higher late-type galaxy fraction in the fields, remains one of the most important problems in cosmology. This relation holds also for dwarf galaxies. Gas-poor dwarf galaxies are strongly clustered, while gas-rich dwarf galaxies appear to be the most weakly clustered objects (Binggeli et al. 1987; Ferguson & Sandage 1988; Binggeli et al. 1990, Thuan et al. 1991).
Since many clusters of galaxies contain an appreciable amount of hot gas, the intracluster medium (ICM), gas will be stripped from galaxies that move through the ICM, if the ram pressure exceeds the internal gravitational force. There are actually many observational evidences for the ram pressure stripping of giant galaxies in clusters of galaxies (Irwin et al. 1987; White et al. 1991; Böhringer et al. 1995). In addition, the physics of ram pressure stripping of giant galaxies has been studied in detail using hydrodynamic simulations (e.g. Balsara, Livio, & O’Dea 1994 and references therein). Since dwarf galaxies have smaller escape velocities, ram pressure stripping is expected to be more efficient than for giant galaxies. However, no attempts have ever been made to examine the amount of gas that could be stripped from low-mass galaxies with shallow gravitational potential wells during their passage through an ICM.
This situation motivated us to consider the ram pressure stripping of the diffuse gas phase of dwarf galaxies by the ICM which may be of prime importance to understand the morphology-density relation for dwarf galaxies. We estimate the critical condition for gas removal by the interaction between an ICM and gaseous dwarf galaxies confined by a surrounding cold dark matter (CDM) halo. Moreover, we verified that condition using an axial symmetric two-dimensional hydrodynamic code which is based on the piecewise parabolic method (PPM) described by Colella & Woodward (1984).
In §2, prior to the main subject, we describe the model of dwarf galaxies which is treated as a one-parameter family defined by the core mass of the CDM halos, using the observed scaling relation (Burkert 1995). In §3, we give the analytical estimates of the critical conditions for the gas ablation from dwarf galaxies due to ram pressure stripping by the ICM. In §4, we demonstrate the complex interaction between dwarf galaxies and the ICM using hydrodynamical simulations and compare the simulation results with the analytical predictions. In §5, we summarise the results of this paper, and discuss observational and theoretical implications.
## 2. Model of dwarf galaxies
The hierarchical clustering model of galaxies has been successful in explaining the clustering pattern of galaxies revealed by redshift surveys. On the other hand, the hierarchical model predicts a large number of low-mass galaxies formed at high redshifts beyond that estimated from the observed luminosity function of galaxies (White & Frenk 1991; Cole et al. 1994). Therefore, several mechanisms for suppressing and/or delaying the formation of dwarf galaxies have been proposed in order to remove this serious discrepancy.
Using a three-dimensional $`N`$-body/SPH simulation code combined with stellar population synthesis, Mori et al. (1997) showed that an energy feedback via supernovae and stellar winds from massive stars keeps a large fraction of the gas hot and suppresses the formation of dwarf galaxies. Moreover, they pointed out that this energy feedback is a necessary mechanism to reproduce the internal structure and the photometric quantities of the nearby dwarf galaxies (see also Dekel & Silk 1986; Yoshii & Arimoto 1987; Burkert & Ruiz-Lapuernnte 1997; Mori, Yoshii & Nomoto 1999).
Babul & Rees (1992) and Efstathiou (1992) argued that the formation epoch of dwarf galaxies is delayed until $`z<1`$ due to the photoionization of the gas by the ultraviolet background radiation at high redshifts. The ionizing background at $`z>1`$ is high enough to keep the gas in dwarf galaxy halos confined and stable, neither able to escape, nor able to collapse (see also Rees 1986; Ikeuchi 1986; Katz, Weinberg & Hernquist 1996; Navarro & Steinmetz 1997).
Following this scenario, we assume that a (proto-) dwarf galaxy has an already virialized CDM halo and a large amount of extended hot gas heated either due to a first population of supernovae or due to the photoionization by ultraviolet background radiation at the formation epoch of the galaxy clusters.
### 2.1. Structure of the CDM halos
The formation of the CDM halos through hierarchical clustering predicts that the equilibrium density profile of a CDM halo has a central cusp (Dubinski & Carlberg 1991; Navarro, Frenk & White 1996; Fukushige & Makino 1997; Moore et al. 1998). However, recent observed rotation curves of nearby dwarf galaxies rule out the singular density profile of dark matter halos (Flores & Primack 1994; Moore 1994). Moreover, Burkert (1995) points out that the density profile
$`\rho _\mathrm{d}(r)={\displaystyle \frac{\rho _{\mathrm{d0}}r_0^3}{(r+r_0)(r^2+r_0^2)}},`$ (1)
nicely reproduces the rotation curves of nearby dwarf galaxies and the central density $`\rho _{\mathrm{d0}}`$ is correlated with the core radius $`r_0`$ through a simple scaling relation.
In this paper we assume that the density distribution of the CDM halos is represented by Burkert’s (1995) profile for a wide range of core masses from $`10^6M_{}`$ to $`10^{10}M_{}`$. For the density distribution of the CDM halo in equation (1),the potential is found by integrating Poisson’s equation as
$`\mathrm{\Phi }_\mathrm{d}(r)`$ $`=`$ $`\pi G\rho _{\mathrm{d0}}r_0^2[\pi 2(1+{\displaystyle \frac{r_0}{r}})\mathrm{arctan}{\displaystyle \frac{r_0}{r}}`$ (2)
$`+`$ $`2\left(1+{\displaystyle \frac{r_0}{r}}\right)\mathrm{ln}\left(1+{\displaystyle \frac{r}{r_0}}\right)`$
$``$ $`(1{\displaystyle \frac{r_0}{r}})\mathrm{ln}\{1+\left({\displaystyle \frac{r}{r_0}}\right)^2\}].`$
The central potential is given by $`\mathrm{\Phi }_\mathrm{d}(0)=\pi ^2G\rho _{\mathrm{d0}}r_0^2`$, where $`G`$ is the gravitational constant. The mass distribution of the CDM halo is given by
$`M_\mathrm{d}(r)=\pi \rho _{\mathrm{d0}}r_0^3[2\mathrm{arctan}{\displaystyle \frac{r}{r_0}}`$
$`+2\mathrm{ln}(1+{\displaystyle \frac{r}{r_0}})+\mathrm{ln}\{1+\left({\displaystyle \frac{r}{r_0}}\right)^2\}].`$ (3)
Using the observed scaling relation derived by Burkert (1995), $`\rho _{\mathrm{d0}}`$ and $`r_0`$ are calculated as
$`r_0=3.07\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{3}{7}}\mathrm{kpc},`$ (4)
and
$`\rho _{\mathrm{d0}}=1.46\times 10^{24}\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{2}{7}}\mathrm{g}\mathrm{cm}^3,`$ (5)
where the core mass $`M_0`$ is the total mass of the CDM halo inside $`r_0`$.
The circular velocity $`v_\mathrm{c}^2=GM_\mathrm{d}(r)/r`$ of the CDM halo has a maximum value
$`v_{\mathrm{c},\mathrm{max}}=48.7\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{2}{7}}\mathrm{km}\mathrm{s}^1,`$ (6)
at the radius of $`r_{\mathrm{c},\mathrm{max}}=3.24r_0`$ (see Fig. 1). This velocity and the radius are related by
$`v_{\mathrm{c},\mathrm{max}}=10.5\left({\displaystyle \frac{r_{\mathrm{c},\mathrm{max}}}{\mathrm{kpc}}}\right)^{\frac{2}{3}}\mathrm{km}\mathrm{s}^1.`$ (7)
### 2.2. Structure of the gas
We focus on the general question of the interaction between the intracluster medium and a hot gaseous component in dwarf galaxies that is generated e.g. by the supernova heating (Dekel & Silk 1986; Burkert & Ruiz-Lapuente 1997; Mori et al. 1997; Mori et al. 1999), or by the photoionization due to the ultraviolet background radiation (Efstathiou 1992; Babul & Rees 1992). In this case the gas is pressure supported justifying the assumption of an initially spherically symmetric distribution and of axi-symmetry.
The self-gravity of the gas is neglected for simplicity and the gas is assumed to be in hydrostatic equilibrium with the constant temperature
$`T`$ $`=`$ $`{\displaystyle \frac{\mu m_\mathrm{p}}{3k_\mathrm{B}}}{\displaystyle \frac{GM_0}{r_0}},`$ (8)
$`=`$ $`3.45\times 10^4\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{4}{7}}\mathrm{K},`$ (9)
where $`\mu `$ is the mean molecular weight, $`m_\mathrm{p}`$ is the proton mass, and $`k_\mathrm{B}`$ is the Boltzmann’s constant. The density distribution of the gas is given by
$`\rho _\mathrm{g}(r)=\rho _{\mathrm{g0}}\left[{\displaystyle \frac{\mu m_\mathrm{p}}{k_\mathrm{B}T}}\left\{\mathrm{\Phi }_\mathrm{d}(r)\mathrm{\Phi }_\mathrm{d}(0)\right\}\right],`$ (10)
where $`\rho _{\mathrm{g0}}`$ is the central density of the gas. The mass ratio between the gas and the CDM halo within a core radius is defined as
$`F={\displaystyle \frac{M_{\mathrm{g0}}}{M_0}},`$ (11)
where $`M_{\mathrm{g0}}`$ is the total gas mass inside a core radius. The relation between the central density of the gas and of the dark matter is given by the numerical solution of equation (11): $`\rho _{\mathrm{g0}}=1.40F\rho _{\mathrm{d0}}`$. We assumed initially $`F=0.1`$ in this paper. Consequently each dwarf galaxy is described as a one-parameter family by the core mass $`M_0`$ of the CDM halos.
## 3. Analytical model
The stripping process is generally classified into two distinct types such as the instantaneous stripping phase and the continuous stripping phase due to Kelvin-Helmholtz instability. In this section, we estimate analytically the gas ablation affects on the evolution of a dwarf galaxy.
### 3.1. Instantaneous ram pressure stripping
We assume that the CDM halo associated with the dwarf galaxies moves in a homogeneous ICM with number density of $`n_{\mathrm{CG}}=\rho _{\mathrm{CG}}/(\mu m_\mathrm{p})10^4`$ cm<sup>-3</sup> and relative velocity $`v_{\mathrm{gal}}10^3`$ km s<sup>-1</sup> corresponding to the three-dimensional velocity dispersion of the galaxy cluster.
The complete condition for the ram pressure stripping from a dwarf galaxy requires that the ram pressure of the ICM exceeds the thermal pressure $`\rho _{\mathrm{g0}}k_\mathrm{B}T/(\mu m_\mathrm{p})`$ at the center of the gravitational potential well of the CDM halo. Using equation (8), such a condition is described by
$`\rho _{\mathrm{CG}}v_{\mathrm{gal}}^2>{\displaystyle \frac{GM_0\rho _{\mathrm{g0}}}{3r_0}}.`$ (12)
Thus, if the core mass $`M_0`$ is smaller than the critical core mass
$`M_{\mathrm{IS}}=1.27\times 10^9\left({\displaystyle \frac{F}{0.1}}\right)^{\frac{7}{2}}\left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^{\frac{7}{2}}`$
$`\times \left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^7M_{},`$ (13)
the gas in a dwarf galaxy is totally stripped by the ram pressure of the ICM.
The collision between the gas of the dwarf galaxy and the ICM generates a forward shock that propagates through the gas of the dwarf galaxy and a reverse shock that propagates through the ICM. The time-scale of the gas removal is roughly estimated by dividing the diameter of the core by the velocity of the forward shock $`v_{\mathrm{fs}}`$. The analysis of the one-dimensional shock problem gives
$`v_{\mathrm{fs}}={\displaystyle \frac{4}{3}}\sqrt{{\displaystyle \frac{\rho _{\mathrm{CG}}}{\rho _{\mathrm{g0}}}}}v_{\mathrm{gal}},`$ (14)
for the high-speed collision of two homogeneous non-gravitating media with large density ratio ($`\rho _{\mathrm{CG}}\rho _{\mathrm{g0}}`$). The time-scale of the mass removal from dwarf galaxies is estimated by
$`\tau _{\mathrm{IS}}={\displaystyle \frac{2r_0}{v_{\mathrm{fs}}}},`$ (15)
$`=2.02\times 10^8\left({\displaystyle \frac{F}{0.1}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{2}{7}}`$
$`\times \left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^{\frac{1}{2}}\left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^1\mathrm{yr}.`$ (16)
Consequently, the gas is totally stripped from less-massive dwarf galaxies ($`M_010^9M_{}`$) in the typical environment of the ICM because the ram pressure of the ICM exceeds the gravitational force of these galaxies. The time-scale of the gas removal is small compared with the characteristic time-scale (several Gyr) a galaxy spends within potentials, therefore, lose their gas easily within this dwarf galaxies ($`M_010^9M_{}`$) the gas might be significantly removed except around the central region of the gravitational potential well.
### 3.2. Kelvin-Helmholtz instability
Gas in the massive dwarf galaxies surviving the instantaneous ram pressure stripping is subsequently removed by Kelvin-Helmholtz instability occurring at the interface between the gas in the dwarf galaxy and the ICM. Murray et al.(1993) estimated the characteristic growth time of Kelvin-Helmholtz instability for a dense cloud embedded in a low-density background. They showed that the growth time is comparable to the sound crossing time of the gas cloud if gravity is negligible. On the contrary, the presence of a gravitational field by the surrounding CDM halo tends to suppress the instability and to stabilize the gas against removal from the dwarf galaxy.
We suppose that the process of the instanteneous ram pressure stripping removes the gas distributed outside a core radius $`r_0`$. Moreover, a sharp discontinuity of the gas density is established between the ICM and the galactic gas having nearly constant density after the shock passage. The unstable wave number of Kelvin-Helmholtz instability at the interface for the incompressible fluid is given by
$`k>g{\displaystyle \frac{\rho _{\mathrm{CG}}^2\rho _{\mathrm{g},\mathrm{avg}}^2}{\rho _{\mathrm{CG}}\rho _{\mathrm{g},\mathrm{avg}}v_{\mathrm{gal}}^2}},`$ (17)
where $`\rho _{\mathrm{g},\mathrm{avg}}`$ is the mean gas density inside $`r_0`$ and $`g=GM_0/r_0^2`$ is the gravitational acceleration at the fluid interface (cf. Chandrasekhar 1961). Since the dominant wavelength for the gas ablation by Kelvin-Helmholtz instability is the order of $`r_0`$ (Murray et al. 1993), the inequality (17) for $`\rho _{\mathrm{CG}}\rho _{\mathrm{g},\mathrm{avg}}`$ is transformed as
$`\rho _{\mathrm{CG}}v_{\mathrm{gal}}^2>{\displaystyle \frac{GM_0\rho _{\mathrm{g},\mathrm{avg}}}{2\pi r_0}}.`$ (18)
Using this inequality, it is found that the diffuse gas component in these galaxies is removed if their core mass $`M_0`$ does not exceed a critical value given by
$`M_{\mathrm{KH}}=1.60\times 10^{12}\left({\displaystyle \frac{F}{0.1}}\right)^{\frac{7}{2}}\left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^{\frac{7}{2}}`$
$`\times \left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^7M_{}.`$ (19)
This value corresponds to masses characteristic for a giant galaxy which indicates that dwarf galaxies should in general be affected by Kelvin-Helmholtz stripping. The gas ablation due to Kelvin-Helmholtz instability is effective for massive dwarf galaxies if the mass-loss time-scale is smaller than the characteristic time-scale a galaxy spends spends within a cluster environment.
According to Nulsen (1982) the mass-loss rate from the galaxy through Kelvin-Helmholtz instability is estimated as
$`\dot{M}_{\mathrm{KH}}=\pi r_0^2\rho _{\mathrm{CG}}v_{\mathrm{gal}},`$ (20)
the gas removal time-scale of the dwarf galaxies is therefore given by
$`\tau _{\mathrm{KH}}`$ $`=`$ $`{\displaystyle \frac{FM_0}{\dot{M}_{\mathrm{KH}}}},`$ (22)
$`=`$ $`2.19\times 10^9\left({\displaystyle \frac{F}{0.1}}\right)\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{1}{7}}\left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^1`$
$`\times \left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^1\mathrm{yr}.`$
Moreover, it should be noted that the characteristic time-scale, normalized by the dynamical time
$`\tau _{\mathrm{dyn}}`$ $`=`$ $`\sqrt{{\displaystyle \frac{3\pi }{32G\rho _{\mathrm{d},\mathrm{avg}}}}},`$ (23)
$`=`$ $`8.93\times 10^7\left({\displaystyle \frac{M_0}{10^9M_{}}}\right)^{\frac{1}{7}}\mathrm{yr},`$ (24)
where $`\rho _{\mathrm{d},\mathrm{avg}}`$ is the mean density of the CDM halo inside $`r_0`$, is given by
$`{\displaystyle \frac{\tau _{\mathrm{KH}}}{\tau _{\mathrm{dyn}}}}=24.5\left({\displaystyle \frac{F}{0.1}}\right)\left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^1`$
$`\times \left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^1.`$ (25)
This equation does not have an explicit dependence on the mass of the CDM halo. Massive dwarf galaxies lose the extended gas within $`25\tau _{\mathrm{dyn}}`$ for $`n_{\mathrm{CG}}=10^4`$ cm<sup>-3</sup> and $`v_{\mathrm{gal}}=1000`$ km s<sup>-1</sup>. Consequently, we expect that even for the massive dwarf galaxy the extended gas will be stripped through Kelvin-Helmholtz instability within a short time-scale.
On the other hand, Balsara, Livio, & O’Dea (1994) demonstrated that the gas accretes from downstream into the core in their hydrodynamical simulations of the interaction between the ICM and the giant galaxy. Since the dwarf galaxies have shallower potential wells than the giant galaxies, we neglect the effect of the gas accretion form backward into the galaxy. It may play however a role for the massive dwarf galaxies where our estimation of the gas ablation might not be suitable.
## 4. Numerical Model
In addition to the effect of the accretion inflow into the core from downstream, the above analytic arguments neglect the geometrical effects between the gas in dwarf galaxies and the ICM and the presence of the Rayleigh-Taylor instability. Moreover, the compressibility of the gas may become important because the each galaxy moves in the cluster of galaxies with the transonic velocity that is corresponding to the velocity dispersion of the cluster of galaxies. Therefore, a realistic analysis that uses a hydrodynamic simulation is necessary to examine the effect of the gas removal due to the interaction between a dwarf galaxy and an ICM. In this section, we describe the numerical model of ram pressure stripping from dwarf galaxies.
### 4.1. Initial conditions and simulation method
For modeling the interaction between a galaxy and an ICM, Portnoy, Pistinner & Shaviv (1993) showed that a treatment of the protons and the electrons as two fluids resulted in negligible differences for the mass of the gas inside the galaxy. In this paper, therefore, the evolution of the gas is described by the hydrodynamic equations for the single perfect fluid. The continuity equation, the momentum equation, and the thermal energy equation are given by
$`{\displaystyle \frac{\rho }{t}}+(\rho 𝒗)=0,`$ (26)
$`{\displaystyle \frac{\rho 𝒗v}{t}}+(\rho 𝒗𝒗)+P=\mathrm{\Phi }_\mathrm{d},`$ (27)
and
$`{\displaystyle \frac{\rho e}{t}}+(\rho e𝒗)+(P𝒗)=\rho 𝒗\mathrm{\Phi }_\mathrm{d},`$ (28)
where $`\rho `$ is the gas density, $`𝒗`$is the gas velocity, $`P`$ is the gas pressure, $`\gamma (=5/3)`$ is the adiabatic index, and $`e`$ is the total specific energy
$`e={\displaystyle \frac{1}{2}}v^2+{\displaystyle \frac{1}{\gamma 1}}{\displaystyle \frac{P}{\rho }}.`$ (29)
Initial conditions are generated using the descriptions in §2. We neglect the effects of the self-gravity of the gas, radiative cooling, the heating by supernovae and stellar winds from massive stars, and the photoionization by ultraviolet radiation background. We discuss them in §5. The equations are solved by a finite difference code VH-1. VH-1 is based on the piecewise parabolic method (PPM) described by Colella & Woodward (1984) and was written and tested by the numerical astrophysics group at the Virginia Institute for Theoretical Astrophysics. Since the PPM scheme has a great advantage due to the reduction of numerical viscosity, all fluid interfaces are sharply preserved and small-scale features can be resolved. This scheme is, therefore, suited for this class of problems.
The system is assumed to be axial symmetric and described by cylindrical geometry. The center of the fixed gravitational potential given by equation (2) is located on the symmetric axis and the gas distribution of the dwarf galaxy is set up using equation (10). The reflecting boundary condition is adopted at the bottom boundary that corresponds to the symmetric axis and the ICM flows continuously from the left stream boundary parallel to the axis. The outflow boundary conditions, by imposing for each variable a zero gradient ($`d/dr=0`$), are adopted at the top and the right stream boundaries. Unfortunately, this condition does not eliminate reflection waves from the boundaries. Thus, the simulation box is separated by two parts such as the inner rectangle $`8r_0`$ long and $`3r_0`$ wide and an outer surrounding part to keep the boundaries far away. The grids are equally spaced with 100 zones per core radius for the higher resolution runs or 50 zones for the lower resolution runs in the inner part, and are exponentially spaced in the outer part. A useful additional effect of this non-uniform grid is the increased numerical dissipation of disturbances that propagate to a large distance. This helps solving the problem of residual reflection waves from the boundary.
### 4.2. Results of simulations
We have performed a parameter study varying the core mass $`M_0`$ of the CDM halo from $`10^6M_{}`$ to $`10^{10}M_{}`$, varying the relative velocity of the dwarf galaxy $`v_{\mathrm{gal}}`$ from 500 km s<sup>-1</sup> to 1000 km s<sup>-1</sup>, and varying the number density $`n_{\mathrm{CG}}`$ of the ICM from $`10^5`$ cm<sup>-3</sup> to $`10^3`$ cm<sup>-3</sup>. Table 1 displays the model parameters of the ICM in this paper.
Figure 2 shows snapshots of the run for the model ($`c`$) of $`M_0=10^7M_{}`$, $`n_{\mathrm{CG}}=10^4`$ cm<sup>-3</sup>, and $`v_{\mathrm{gal}}=1000`$ km s<sup>-1</sup> as a function of elapsed time at $`5.61\times 10^6,2.86\times 10^7,4.58\times 10^7,7.42\times 10^7`$, and $`9.63\times 10^7`$ yr. The left and the right colour images show the logarithmic density distribution and the logarithmic pressure distribution, respectively. Arrows in the right panels indicate the velocity vector of the gas flow at each point and their lengths are proportional to the absolute values of their velocity. The CDM halo is located at the origin and the ICM flows in from the left side.
The upper panel shows the early phase of the interaction between the ICM and the gas of the dwarf galaxy. It is clearly seen that the forward shock that propagates through the gas of dwarf galaxy, and the reverse shock that propagates through the ICM are formed in the downstream and the upstream of the contact discontinuity respectively. The propagation of the forward shock causes the strong compression of the gas in the dwarf galaxy as seen in the upper and the second panel. Below the second panel, Kelvin-Helmholtz instabilities are observed at the contact discontinuity.
Since the instabilities are suppressed by the gravity, the growing modes are only the disturbances of the large wave numbers that have the shortest growth time. Disturbances with small wave numbers grow while moving away from the galaxy. The third panel shows that the gas associated with the dwarf galaxy is accelerated by the ram pressure of the ICM and is pushed out of the potential well. Since the front of the accelerated gas is highly Rayleigh-Taylor unstable, there are many small irregularities in front of the contact discontinuity. Accurately, the feature seen on the symmetric axis ($`R=0`$) at $`Z0.5r_0`$ is caused by pure Rayleigh-Taylor instability. The combination of Rayleigh-Taylor instability and Kelvin-Helmholtz instability forms other irregularities because the gas flows tangentially to the fluid interface as seen in the third and right panel. However, the perturbations do not grow significantly because the gravity of the CDM halo suppresses Rayleigh-Taylor instability. The gas removal time-scale nicely agrees with the analytical estimate in equation (16) which predicts $`\tau _{\mathrm{IS}}=5.42\times 10^7`$ yr.
Figure 3 shows snapshots of the run for the model ($`c`$) of $`M_0=10^{10}M_{}`$, $`n_{\mathrm{CG}}=10^4`$ cm<sup>-3</sup>, and $`v_{\mathrm{gal}}=1000`$ km s<sup>-1</sup> as a function of elapsed time at $`6.84\times 10^7,6.99\times 10^8,1.15\times 10^9,2.18\times 10^9`$, and $`3.26\times 10^9`$ yr.
The early evolution of the interaction between the gas and the ICM is almost the same as in the case of $`M_0=10^7M_{}`$. However, since the gravitational potential is deep and the central thermal pressure is larger than the ram pressure of the ICM as shown in the analytical model, the mass-loss process is not instantaneous but mild ablation due to Kelvin-Helmholtz instabilities occur. Moreover, the gas interface is Rayleigh-Taylor stable because the gas acceleration is week. The velocity field in the right bottom panel reveals that the flow turns back and an accretion inflow into the core develops close to the symmetry axis on the downstream side. Accretion occurs quasi-periodically as a radial-pumping mode. The condition of the flow is summarized as follows.
1. This flow from the downstream side and the post-shock gas from the upstream side cause increased compression mainly along the symmetry axis.
2. The gas in the galaxy, therefore, acts by expanding sideways as seen in the third panel. Since this expansion increases the cross section of the interaction between the galactic gas and the ICM, the mass-loss rate increases.
3. The expansion decreases the pressure gradient around the galaxy center and the gas contracts to the galaxy center by the gravitational force. Then the system recovers the quasi-static states again.
4. This cycle repeats again.
This process is effective in the simulations of a massive dwarf galaxy or of a small ram pressure of the ICM. For this specific process, the analytic condition given by §3.2 fails to describe the gas ablation from the dwarf galaxy.
Figure 4 shows the evolution of the gas mass inside the core radius around the center of the dark matter halo as a function of time for core masses of the CDM halo from $`10^6M_{}`$ to $`10^{10}M_{}`$. Each curve is categorized by the model parameters (see Table 1). The gas is instantaneously stripped in low-mass dwarfs or for large ram pressures. Due to the accretion inflow into the galactic core from the down stream, the net mass-loss rate becomes small. Moreover, we can observe oscillations of the total mass inside a core radius due to the mass accretion. These oscillation modes have periods of several dynamical timescales. Even for the massive dwarf galaxy ($`M_0=10^{10}M_{}`$), the gas could be instantaneously stripped if the ram pressure is large as e.g. in the cores of galactic clusters.
Figure 5 summarizes the results of our simulations. This diagram shows the mass-loss rate due to ram pressure stripping as a sequence of the core mass. Filled circles indicate the cases where the gas in the core radius is completely stripped within 1 Gyr. In these cases the ram pressure stripping by the ICM is very effective and the whole gas in the potential well is rapidly removed. Open circles indicate the cases where the gas in the core radius is not completely stripped within 2 Gyr. In these cases, the ram pressure stripping by the ICM is not so effective. The thin line indicates the relation of
$`\rho _{\mathrm{CG}}v_{\mathrm{gal}}^2={\displaystyle \frac{GM_0\rho _{\mathrm{g0}}}{3r_0}},`$ (30)
which is the instantaneous stripping condition described in §3.1. This line divides the parameter space into two regions. Stripping is effective for the upper region and is ineffective for the lower region. This relation roughly agrees with the numerical experiments in the range from $`M_0=10^6M_{}`$ to $`M_0=10^{10}M_{}`$. The dashed line indicates the relation of
$`\rho _{\mathrm{CG}}v_{\mathrm{gal}}^2={\displaystyle \frac{GM_0\rho _{\mathrm{g},\mathrm{avg}}}{2\pi r_0}},`$ (31)
which is the Kelvin-Helmholtz stripping condition described in §3.2. Stripping by Kelvin-Helmholtz instability is effective for the upper region and is not effective for the lower region. Since there is mass accretion from the back of the galaxy, this relation is not well reproduced by the numerical experiments.
Using the data of intermediate cases that are shown by filed triangles in Figure 5, we find a critical core mass for effective ram pressure stripping of
$`M_{\mathrm{cr}}=2.52\times 10^9\left({\displaystyle \frac{n_{\mathrm{CG}}}{10^4\mathrm{cm}^3}}\right)^{\frac{5}{2}}\left({\displaystyle \frac{v_{\mathrm{gal}}}{10^3\mathrm{km}\mathrm{s}^1}}\right)^5M_{}.`$ (32)
The thick line shows the loci of this critical core mass.
In Figure 5, $`r_{\mathrm{CG}}`$ is the distance from the center of galaxy cluster assuming the isothermal $`\beta `$-model
$`n_{\mathrm{CG}}=n_{\mathrm{CG0}}\left\{1+\left({\displaystyle \frac{r_{\mathrm{CG}}}{r_{\mathrm{CG0}}}}\right)^2\right\}^{\frac{3}{2}\beta },`$ (33)
with $`\beta =0.6,n_{\mathrm{CG0}}=2\times 10^3`$ cm$`{}_{}{}^{3},r_{\mathrm{CG0}}=0.25`$ Mpc, and $`\sigma _{\mathrm{CG}}=866`$ km s<sup>-1</sup>, where $`n_{\mathrm{CG0}}`$ and $`r_{\mathrm{CG0}}`$ is the central number density and the core radius of an ICM and $`\sigma _{\mathrm{CG}}`$ is the line-of-sight velocity dispersion of the galaxy cluster which corresponds to the relative velocity of the dwarf galaxy. If the galaxy number density follows an approximate King model as
$`n_{\mathrm{gal}}\left\{1+\left({\displaystyle \frac{r_{\mathrm{CG}}}{r_{\mathrm{CG0}}}}\right)^2\right\}^{3/2},`$ (34)
the median radius of the galaxy distribution for $`r<10r_{\mathrm{CG0}}`$ is about $`3.5r_{\mathrm{CG0}}`$. The shaded region corresponds to ICM conditions inside this radius. The diagram indicates that the process of the ram pressure stripping is very effective for dwarf galaxies in the clusters of galaxies. There, diffuse gas in dwarf galaxies should rapidly be removed.
## 5. Summary and discussion
The physics of the interaction between the ICM and dwarf galaxies confined by a surrounding CDM halo has been studied in detail using analytical estimates and hydrodynamic simulations. We have performed a parameter study varying the core mass of the CDM, the relative velocity of the galaxy and the density of the ICM. We find that the gas in dwarf galaxies is rapidly removed in a typical cluster environment by ram-pressure stripping.
Our results can be applied to clusters of galaxies, where X-ray emission has been observed. Their gas number density $`n_{\mathrm{CG}}`$ at the median radius of the galaxy distribution and the line-of-sight velocity dispersion $`\sigma _{\mathrm{CG}}`$ are plotted in Figure 6. Data for the central gas density and $`\beta `$ are from Jones & Forman (1999) and Briel, Henry & Böhringer (1992), and the velocity dispersions are from the same sources including also Hughes (1989). One particular example, the Coma cluster with $`n_0=2.89\times 10^3`$ cm$`{}_{}{}^{3},\beta =0.75`$ and $`\sigma _{\mathrm{CG}}=1010`$ km s<sup>-1</sup>, is indicated by a filled circle. Dotted lines are loci of constant $`M_{\mathrm{cr}}`$ given by equation (32) ranging from $`10^7M_{}`$ to $`10^{11}M_{}`$. The lines with their various core masses and maximum circular velocities are shown. They indicate that galaxies should rapidly loose all of their diffuse gas by ram pressure stripping. This figure shows that the gas in dwarf galaxies is completely stripped in rich clusters of galaxies within the short time-scale.
The effect of radiative cooling has been neglected in our simulations. we estimate the role of radiative cooling in our models here. Since for low-mass dwarf galaxies ($`M_010^9M_{}`$), the ram pressure already exceeds the gravitational force, radiative processes will not be able to affect the gas stripping. However, it might play a role for massive dwarf galaxies ($`M_010^9M_{}`$) which have deeper potential wells. The cooling time-scale is defined as (Efstathiou 1992)
$`\tau _{\mathrm{cool}}={\displaystyle \frac{3}{2}}{\displaystyle \frac{1}{\mu ^2(1Y)^2}}{\displaystyle \frac{k_\mathrm{B}T}{n_\mathrm{g}\mathrm{\Lambda }}},`$ (35)
where $`Y(=0.25)`$ is the helium mass fraction, $`n_\mathrm{g}`$ is the gas number density, and $`\mathrm{\Lambda }`$ is cooling rate. If we evaluate this time-scale for our standard model $`M_0=10^9M_{}`$ using conditions immediately behind the leading shock, we find that $`\tau _{\mathrm{cool}}=3.18\times 10^9`$ yr assuming $`T=10^7`$ K, $`\mathrm{\Lambda }=10^{23}`$ erg s<sup>-1</sup> cm<sup>-3</sup>, and $`n_\mathrm{g}=1.0`$ cm<sup>-3</sup> from the result of our simulations. This time-scale is larger than the instantaneous stripping time-scale ($`\tau _{\mathrm{IS}}=2.03\times 10^8`$ yr) that is given by equation (16). This implies that gas stripping from the massive dwarf galaxies is affected by the radiative cooling.
However, we would then also need to consider the effects of the star formation and subsequent feedback process such as stellar winds and supernovae heating from massive stars. These feedback processes supply thermal energy to the gas and prevent efficient cooling (Dekel & Silk 1986; Mori et al. 1997; Mori et al. 1999). In addition, the photoionization due to the ultraviolet background radiation prevents gas cooling and keeps the gas hot (Efstathiou 1992; Babul & Rees 1992). Furthermore, though we studied only the interaction between the ICM and a hot interstellar medium in this paper, it is very interesting to examine also the fate of the cool and dense interstellar medium in a dwarf galaxy. In this case, the stripping history of the gas may be quite different even for the less-massive dwarf galaxies ($`M<10^9M_{}`$). In a series of forthcoming studies, we will report the results of taking into account the multi-phase states of the gas with cool components in a galaxy, including the effect of the radiative cooling of the gas, energy input from stars, and the ultraviolet background radiation.
We conclude that the ram pressure stripping of the diffuse gas from dwarf galaxies by the ICM is of prime importance to understand the morphology-density relation for dwarf galaxies and the chemical evolution of dwarf galaxies. Most of the heavy elements ejected by massive stars through Type II supernovae are found in the warm and hot diffuse gas. In order for these heavy elements to be incorporated into stars, the hot gas has to cool and form cold clouds again. Our models indicate that this will be impossible for dwarf galaxies in clusters of galaxies, as the gas is stripped before condensing into clouds. Ram pressure stripping therefore plays a significant role for the chemical evolution of dwarf galaxies. Our results imply that most of the chemical evolution of dwarf galaxies must have occured at high redshift, before the virialized cluster had formed.
This work has been supported in part by the Research Fellowship of the Japan Society for the Promotion of Science for Young Scientists (6867) of the Ministry of Education, Science, Sports, and Culture in Japan. M.M. would like to thank the Max-Planck-Institut für Astronomie for its hospitality throughout this research. Numerical calculations were carried out at Max-Planck-Institut für Astronomie in Germany and at the Astronomical Data Analysis Center of the National Astronomical Observatory in Japan. |
warning/0001/gr-qc0001029.html | ar5iv | text | # Scalar-Tensor Gravity in Two 3-brane System
## I introduction
Randall and Sundrum recently proposed a mechanism to solve the hierarchy problem using two 3-branes of opposite tension residing embedded in a five dimensional spacetime. In this scenario we are assumed to live in the negative tension brane. They also found that even the graviton is trapped on the positive tension brane in the same set up (see also ).
It is natural to ask how gravity look like on the brane. In a recent paper, Garriga and Tanaka have analyzed the metric perturbation equation on the 3-brane in the background spacetime of Randall and Sundrum and have shown that gravity on the brane is described by Brans-Dicke gravity (for a single 3-brane, see ). Charmousis et al. recently proposed an Ansatz for the radion which solves five-dimensional linearized equations of motion .
The purpose of this note is to provide alternative considerations on the brane gravity based on the derivation of the low-energy effective action on each 3-brane from the five dimensional action à la Kaluza-Klein. We take into account the degrees of freedom of the modulus field as well as four-dimensional graviton fluctuations following the metric Ansatz by Charmousis et al. We will see that gravity on the brane is actually a more general scalar-tensor gravity theory, where the Brans-Dicke parameter is a function of the Brans-Dicke scalar field, rather than the Brans-Dicke gravity, where the Brans-Dicke parameter is a constant. The modulus field (the radion) plays the role of the Brans-Dicke scalar field. We believe that our approach is technically simpler. The calculations required in our approach are reduced considerably, and it is easy to generalize to include matter fields. We show that the cosmological attractor mechanism toward the Einstein gravity is realized in gravity on a positive tension brane. We also make a brief comment on the cosmological constant problem in the brane-world view of the universe.
## II Reduction in two three-brane system
### A Background
We briefly review the background metric found by Randall and Sundrum to introduce our notations. The set up is that the five-dimensional Einstein gravity with a cosmological constant with two 3-branes located in the orbifold $`S^1/𝐙_2`$ at the fifth dimensional coordinate $`z=0`$(positive tension brane), $`z=r_c`$(negative tension brane). The action is
$$S=2d^4x_0^{r_c}𝑑\varphi \sqrt{g_5}\left[\frac{M^3}{2}R_52\mathrm{\Lambda }\right]\sigma _{(+)}d^4x\sqrt{g_{(+)}}\sigma _{()}d^4x\sqrt{g_{()}},$$
(1)
where $`M`$ is the five-dimensional Planck mass defined in term of the five-dimensional gravitational constant $`G_5`$ as $`M^3=8\pi G_5`$. $`\sigma _{(\pm )}`$ and $`g_{(\pm )}`$ are the brane tension and the induced metric on the brane, respectively. Randall and Sundrum have shown that there exists a solution that respects four-dimensional Poincaré invariance :
$$ds^2=e^{2k|z|}\eta _{\mu \nu }dx^\mu dx^\nu +dz^2,$$
(2)
only if $`\mathrm{\Lambda },\sigma _{(+)}`$ and $`\sigma _{()}`$ are related as
$$\mathrm{\Lambda }=3M^3k^2,\sigma _{(+)}=\sigma _{()}=3M^3k.$$
(3)
### B Naive Ansatz
We shall derive the four-dimensional low-energy effective action on a 3-brane. In order to include massless gravitational degrees of freedom (the zero modes about the background spacetime Eq.(2)), we replace the Minkowski metric $`\eta _{\mu \nu }`$ with a general metric $`\overline{g}_{\mu \nu }(x)`$ and $`z`$-direction length $`r_c`$ with a modulus field $`T(x)`$ :
$$ds^2=e^{2kT(x)|z|}\overline{g}_{\mu \nu }(x)dx^\mu dx^\nu +T(x)^2dz^2.$$
(4)
Although this metric ansatz does not correctly describe the linearized dynamics of massless fields, we present the analysis here because the results nevertheless seem simple and attractive. The induced metric on the positive (or negative) tension brane is thus $`g_{(+)\mu \nu }=\overline{g}_{\mu \nu }`$ (or $`g_{()\mu \nu }=e^{2kr_cT}\overline{g}_{\mu \nu }`$), respectively. Since we impose $`𝐙_2`$ symmetry ($`zz`$), massless vector fluctuations associated with the off-diagonal part of the metric are absent .
#### 1 Effective Four-dimensional Action
We can perform the $`z`$ integral to obtain the four-dimensional action on the positive tension brane described by $`g_{(+)\mu \nu }`$:
$`S_{(+)}`$ $`=`$ $`2{\displaystyle d^4x_0^{r_c}𝑑z\sqrt{g_5}\left[\frac{M^3}{2}R_52\mathrm{\Lambda }\right]}\sigma _{(+)}{\displaystyle d^4x\sqrt{g_{(+)}}}\sigma _{()}{\displaystyle d^4x\sqrt{g_{()}}}`$ (5)
$`=`$ $`2{\displaystyle d^4x_0^{r_c}𝑑z\sqrt{g_{(+)}}e^{4kTz}T\left[\frac{M^3}{2}e^{2kTz}\left(R_{(+)}6e^{kTz}\mathrm{}_{(+)}e^{kTz}\right)2\mathrm{\Lambda }\right]}`$ (7)
$`\sigma _{(+)}{\displaystyle d^4x\sqrt{g_{(+)}}}\sigma _{()}{\displaystyle d^4x\sqrt{g_{(+)}}e^{4kr_cT}}`$
$`=`$ $`{\displaystyle }d^4x\sqrt{g_{(+)}}[{\displaystyle \frac{M^3(1e^{2kr_cT})}{2k}}R_{(+)}3kr_c^2M^3e^{2kr_cT}\left(_{(+)}T\right)^2`$ (9)
$`{\displaystyle \frac{(1e^{4kr_cT})}{k}}\mathrm{\Lambda }\sigma _{(+)}\sigma _{()}e^{4kr_cT}]`$
$`=`$ $`{\displaystyle d^4x\sqrt{g_{(+)}}\left[\frac{M^3(1e^{2kr_cT})}{2k}R_{(+)}3kr_c^2M^3e^{2kr_cT}\left(_{(+)}T\right)^2\right]},`$ (10)
where we have used the relation Eq.(3) in the last equation.
The case of the negative tension brane is obtained by the conformal transformation such that $`g_{()\mu \nu }=e^{2kr_cT}\overline{g}_{\mu \nu }=e^{2kr_cT}g_{(+)\mu \nu }`$, and the result is
$$S_{()}=d^4x\sqrt{g_{()}}\left[\frac{M^3(e^{2kr_cT}1)}{2k}R_{()}+3kr_c^2M^3e^{2kr_cT}\left(_{()}T\right)^2\right].$$
(11)
The above result is also easily obtained by the change of signature such that $`kk`$, the meaning of which may be intuitively clear: exchange of the locations of the branes.
#### 2 Brans-Dicke parameter
The action of the scalar-tensor gravity theory is given by
$$S=d^4x\sqrt{g}\frac{1}{16\pi }\left(\mathrm{\Phi }_{BD}R\frac{\omega (\mathrm{\Phi }_{BD})}{\mathrm{\Phi }_{BD}}(\mathrm{\Phi }_{BD})^2\right).$$
(12)
The correspondence of the gravity theory on the brane to the Brans-Dicke gravity is then immediate. The Brans-Dicke scalar field (or the inverse of the effective gravitational constant) is
$$\frac{1}{G_{\mathrm{eff}}^{(\pm )}}=\mathrm{\Phi }_{BD}^{(\pm )}=\frac{16\pi M^3}{k}e^{kr_cT}\mathrm{sinh}(kr_cT)=\frac{2}{kG_5}e^{kr_cT}\mathrm{sinh}(kr_cT),$$
(13)
while the corresponding Brans-Dicke function is
$$\omega _{(\pm )}(T)=\pm 3e^{\pm kr_cT}\mathrm{sinh}(kr_cT).$$
(14)
Remarkably, despite the defect of the metric Ansatz, both expressions are in perfect agreement with those of Garriga and Tanaka if the length parameter between the branes is replaced with the modulus field $`r_cT(x)`$: the Brans-Dicke gravity with the corresponding Brans-Dicke parameter $`0<\omega <\mathrm{}`$ (for positive tension brane) and $`3/2<\omega <0`$ (for negative tension brane).
### C CGR Ansatz
Although the previous analysis based on the naive Ansatz seems simple and attractive, we seek after another metric Ansatz which does solve the linearized equations of motion. As an example, we adopt the metric Ansatz proposed by Charmousis et al.(CGR) (we shall consider the region $`0<z<r_c`$; the other region by the orbifold symmetry):
$`ds^2=e^{2kh(x,z)}\overline{g}_{\mu \nu }(x)dx^\mu dx^\nu +h_{,z}^2dz^2,`$ (15)
$`h(x,z)=z+f(x)e^{2kz}.`$ (16)
The induced metric on each brane is thus respectively
$`g_{(+)\mu \nu }=e^{2kf}\overline{g}_{\mu \nu },`$ (17)
$`g_{()\mu \nu }=\alpha ^1e^{2\alpha kf}\overline{g}_{\mu \nu }.`$ (18)
Here we have introduced the notation $`\alpha e^{2kr_c}`$ for convenience.
#### 1 Effective Four-dimensional Action
We then perform the $`z`$ integral to obtain the four-dimensional action on the positive tension brane in terms of $`g_{(+)\mu \nu }`$:
$$S_{(+)}=d^4x\sqrt{g_{(+)}}\left[\frac{M^3}{2k}\left(1\frac{1}{\alpha }e^{2(\alpha 1)kf}\right)R_{(+)}3kM^3\left[\left(\alpha +\frac{1}{\alpha }\right)e^{2(\alpha 1)kf}2\right]\left(_{(+)}f\right)^2\right],$$
(19)
where we have again used the relation Eq.(3).
The case of the negative tension brane is obtained by the conformal transformation such that $`g_{()\mu \nu }=\alpha ^1e^{2\alpha kf}\overline{g}_{\mu \nu }=\alpha ^1e^{2(\alpha 1)kf}g_{(+)\mu \nu }`$, and the result is
$$S_{()}=d^4x\sqrt{g_{()}}\left[\frac{M^3}{2k}\left(\alpha e^{2(\alpha 1)kf}1\right)R_{()}+3\alpha ^2kM^3\left[\left(\alpha +\frac{1}{\alpha }\right)e^{2(\alpha 1)kf}2\right]\left(_{()}f\right)^2\right].$$
(20)
#### 2 Brans-Dicke parameter
The Brans-Dicke scalar field (or the inverse of the effective gravitational constant) is
$$\frac{1}{G_{\mathrm{eff}}^{(\pm )}}=\mathrm{\Phi }_{BD}^{(\pm )}=\frac{2}{kG_5}e^{k[(\alpha 1)f+r_c]}\mathrm{sinh}k[(\alpha 1)f+r_c]$$
(21)
while the corresponding Brans-Dicke function is
$$\omega _{(\pm )}(f)=\pm 3e^{\pm k[(\alpha 1)f+r_c]}\mathrm{sinh}k[(\alpha 1)f+r_c]\left(1\frac{e^{\pm (\alpha 1)kf}\mathrm{sinh}[(\alpha 1)kf]}{\mathrm{sinh}^2(kr_c)}\right).$$
(22)
Both expressions are in agreement with those of Garriga and Tanaka in the limit where the radion fluctuations are vanishing, $`f0`$.
Now let us consider what kind of gravity is induced on the brane if the radion fluctuations are present. As long as $`e^{2\alpha kf}\alpha e^{2kr_c}`$ the metric ansatz Eq.(15) and Eq.(16) correctly describes the linearized dynamics of massless fields. For the positive tension brane, gravity on the brane is a class of scalar-tensor gravity theories with $`\omega >0`$ (for $`e^{2\alpha kf}\alpha `$ we have $`\omega =3\alpha e^{2\alpha kf}/2`$) and can satisfy the constraint by the solar system experiment($`|\omega |>3000`$ ) if $`kr_c>4`$. In the limit $`kr_c\mathrm{}`$, the Einstein gravity is recovered with the gravitational constant $`G_4=kG_5`$.
On the other hand, gravity on the negative tension brane is a class of scalar tensor theories with $`\omega =3/2+9\alpha ^1e^{2\alpha kf}/2<0`$ (in the lowest order of $`𝒪(\alpha ^1e^{2\alpha kf})`$) which does not satisfy the constraint by the solar system experiment. This does not immediately mean that the scenario of is not valid because we do not include massive degrees of freedom that may give rise to a stabilizing potential for the modulus(see ).
Moreover, gravity on the negative tension brane exhibits some peculiar behavior in the limit $`kr_c\mathrm{}`$: $`G_{\mathrm{eff}}0`$ and $`\omega 3/2`$. The vanishing of the gravitational constant would imply the absence of gravity force. Furthermore, $`\omega =3/2`$ means that the scalar field degrees of freedom is completely absorbed by the conformal transformation; the action is reduced to that of vacuum Einstein gravity (no scalar field) by the conformal transformation such that $`g_{\mu \nu }^E=[\alpha e^{2(\alpha 1)kf}1]g_{()\mu \nu }`$.We note that in the limit $`kr_c\mathrm{}`$ the Einstein conformal frame metric $`g_{\mu \nu }^E`$ coincides with the metric on the positive tension brane $`g_{(+)\mu \nu }`$. In this sense, the scalar field degrees of freedom is also frozen.
#### 3 Attractor Mechanism
The Einstein gravity is generically a cosmological attractor in scalar-tensor theories (attractor mechanism ). Therefore, as far as gravity is concerned, the dilaton stabilization is not always necessary. It is interesting to examine whether gravity on the positive tension brane contains a similar attractor mechanism although it requires slight abuse the action Eq.(19) beyond the weak field approximation. To do so, we assume matter on the positive tension brane and shall work in the Einstein conformal frame defined by $`g_{\mu \nu }^E=[1\alpha ^1e^{2(\alpha 1)kf}]g_{(+)\mu \nu }e^{2a(\phi )}g_{(+)\mu \nu },2(da/\kappa d\phi )^2=(2\omega (f)+3)^1`$ so that the action Eq.(19) reduces to
$$S_{(+)}=d^4x\sqrt{g_E}\left[\frac{1}{2\kappa ^2}R_E\frac{1}{2}(_E\phi )^2\right],$$
(23)
where $`\kappa ^2=8\pi G_48\pi kG_5`$. An important point of the attractor mechanism by Damour and Nordvedt is that the function $`a(\phi )`$ determines the cosmological dynamics of the dilaton $`\phi `$ and the minimum of it is the Einstein gravity ($`a(\phi )_{,\phi }=0\omega =\mathrm{}`$). Therefore it is sufficient to see the shape of $`a(\phi )`$. For $`e^{2\alpha kf}\alpha `$, we find that $`2a(\phi )=\mathrm{ln}(12\kappa ^2\phi ^2/3)`$ with $`2\kappa ^2\phi ^2=3\alpha ^1e^{2\alpha kf}`$; that is, $`a(\phi )`$ indeed has a minimum at $`\phi =0`$, which corresponds to $`\omega \mathrm{}`$.The attractor mechanism is also realized in the action Eq.(10).
#### 4 Adjusting Mechanism?
Finally let us consider the situation where the bulk cosmological constant $`\mathrm{\Lambda }`$ is slightly perturbed from the assumed value Eq.(3), $`\mathrm{\Lambda }\mathrm{\Lambda }+\delta \mathrm{\Lambda }`$, while the brane tensions are fixed. Such a mismatch would induce an effective potential for the modulus $`f(x)`$ and give rise to the dynamics of $`f(x)`$. However, because the modulus $`f(x)`$ couples non-minimally to the curvature, the effective potential for $`f(x)`$ is given by performing the conformal transformation to the canonical Einstein frame $`g_{\mu \nu }^E=[1\alpha ^1e^{2(\alpha 1)kf}]g_{(+)\mu \nu }`$ (for simplicity we consider the positive tension brane), and it is given by <sup>§</sup><sup>§</sup>§The potential of this type has been discussed in the context of inflation.
$$V_{\mathrm{eff}}(f)=\frac{\delta \mathrm{\Lambda }}{k}\mathrm{coth}k[(\alpha 1)f+r_c].$$
(24)
Since the minimum of the potential is not zero, there is no dynamical mechanism to enforce $`V_{\mathrm{eff}}=0`$. Hence, contrary to the argument in , brane-world picture of the universe seems to be compatible with a nonzero four-dimensional cosmological constant (similar observations are made in ): the four-dimensional background metric can be Minkowski spacetime, de Sitter spacetime, or anti-de Sitter spacetime, depending both on the bulk cosmological constant and on the brane tensions. The cosmological constant problem (in the present context the fine-tuning between the bulk five-dimensional cosmological constant and the brane tension) remains the problem.For stability of the classical Minkowski background, see .
## III summary
We have given simple considerations on the brane gravity based on the dimensional reduction. We have shown that gravity on the brane belongs to a class of scalar-tensor gravity theory where the radion field plays the role of the Brans-Dicke scalar field. We have pointed out the possibility of the cosmological attractor mechanism on the positive tension brane and have also made a brief comment on the cosmological constant problem in the brane world scenario.
###### Acknowledgements.
We would like to thank Yasunori Nomura, Tomohiko Takahashi for useful discussions and Takahiro Tanaka for useful correspondence and comments. This work was supported in part by the JSPS under Grant No. 3596. |
warning/0001/cond-mat0001407.html | ar5iv | text | # 1 Chromatic zeros for the 𝐿_𝑦=3, 𝐿_𝑥=𝑚=10 torus graph.
$`T=0`$ Partition Functions for Potts Antiferromagnets on Square Lattice Strips with (Twisted) Periodic Boundary Conditions
Norman Biggs<sup>*</sup><sup>*</sup>*email: n.l.biggs@lse.ac.uk
Centre for Discrete and Applicable Mathematics
London School of Economics
London WC2A 2AE
UK
Robert Shrock<sup>\**</sup><sup>\**</sup>\**email: robert.shrock@sunysb.edu
Institute for Theoretical Physics
State University of New York
Stony Brook, N. Y. 11794-3840
USA
Abstract
We present exact calculations of the zero-temperature partition function for the $`q`$-state Potts antiferromagnet (equivalently, the chromatic polynomial) for two families of arbitrarily long strip graphs of the square lattice with periodic boundary conditions in the transverse direction and (i) periodic (ii) twisted periodic boundary conditions in the longitudinal direction, so that the strip graphs are embedded on a (i) torus (ii) Klein bottle. In the limit of infinite length, we calculate the exponent of the entropy, $`W(q)`$, show it to be the same for (i) and (ii), and determine its analytic structure.
The chromatic polynomial $`P(G,q)`$ counts the number of ways that one can color a graph $`G`$ with $`q`$ colors such that no two adjacent vertices have the same color (for reviews, see -). The least positive integer $`q`$ for which $`P(G,q)`$ is nonzero is the chromatic number, $`\chi (G)`$. Besides its intrinsic mathematical interest, the chromatic polynomial has an important connection with statistical mechanics since it is the zero-temperature partition function of the $`q`$-state Potts antiferromagnet (AF) on $`G`$: $`P(G,q)=Z(G,q,T=0)_{PAF}`$. The Potts AF exhibits nonzero ground-state entropy $`S_00`$ (without frustration) for sufficiently large $`q`$ on a given lattice graph and is thus an exception to the third law of thermodynamics. This is equivalent to a ground state degeneracy per site $`W>1`$, since $`S_0=k_B\mathrm{ln}W`$. Denoting the number of vertices of $`G`$ as $`n=v(G)`$ and assuming a reasonable definition of $`\{G\}=lim_n\mathrm{}G`$, we have<sup>1</sup><sup>1</sup>1 At certain special points $`q_s`$ (typically $`q_s=0,1,..,\chi (G)`$), one has the noncommutativity of limits $`lim_{qq_s}lim_n\mathrm{}P(G,q)^{1/n}lim_n\mathrm{}lim_{qq_s}P(G,q)^{1/n}`$, and hence it is necessary to specify the order of the limits in the definition of $`W(\{G\},q_s)`$ . We use the first order of limits here; this has the advantage of removing certain isolated discontinuities in $`W`$. $`W(\{G\},q)=lim_n\mathrm{}P(G,q)^{1/n}`$. Since $`P(G,q)`$ is a polynomial, one can generalize $`q`$ from $`_+`$ to $``$. The zeros of $`P(G,q)`$ in the complex $`q`$ plane are called chromatic zeros; a subset of these may form an accumulation set in the $`n\mathrm{}`$ limit, denoted $``$, which is the continuous locus of points where $`W(\{G\},q)`$ is nonanalytic ($``$ may be null, and $`W`$ may also be nonanalytic at certain discrete points). The maximal region in the complex $`q`$ plane to which one can analytically continue the function $`W(\{G\},q)`$ from physical values where there is nonzero ground state entropy is denoted $`R_1`$. The maximal value of $`q`$ where $``$ intersects the (positive) real axis is labelled $`q_c(\{G\})`$.
We consider strips of the square lattice with arbitrary length $`L_x=m`$ vertices and fixed width $`L_y`$ vertices (with the longitudinal and transverse directions taken to be $`\widehat{x}`$ and $`\widehat{y}`$). The chromatic polynomials for the cyclic and Möbius strip graphs of the square lattice were calculated for $`L_y=2`$ in (see also -) and for $`L_y=3`$ in . After studies of the chromatic zeros for $`L_y=2`$ in , $`W`$ and $``$ were determined for this case in and for $`L_y=3`$ in . An important question concerns the effect of boundary conditions (BC’s), and hence graph topology, on $`P`$, $`W`$, and $``$. We use the symbols FBC<sub>y</sub> and PBC<sub>y</sub> for free and periodic transverse boundary conditions and FBC<sub>x</sub>, PBC<sub>x</sub>, and TPBC<sub>x</sub> for free, periodic, and twisted periodic longitudinal boundary conditions. The term “twisted” means that the longitudinal ends of the strip are identified with reversed orientation. These strip graphs can be embedded on surfaces with the following topologies:<sup>2</sup><sup>2</sup>2For the triangular lattice with cylindrical BC’s, $`W`$ and $``$ were calculated in . Other calculations of $`P`$, $`W`$, and $``$ have been performed for strips having BC’s of type (i) -, (ii) , (iii) -, (iv) . : (i) (FBC<sub>y</sub>,FBC<sub>x</sub>): strip; (ii) (PBC<sub>y</sub>,FBC<sub>x</sub>): cylindrical; (iii) (FBC<sub>y</sub>,PBC<sub>x</sub>): cylindrical (denoted cyclic here); (iv) (FBC<sub>y</sub>,TPBC<sub>x</sub>): Möbius; (v) (PBC<sub>y</sub>,PBC<sub>x</sub>): torus; and (vi) (PBC<sub>y</sub>,TPBC<sub>x</sub>): Klein bottle.<sup>3</sup><sup>3</sup>3 These BC’s can all be implemented in a manner that is uniform in the length $`L_x`$; the case (vii) (TPBC<sub>y</sub>,TPBC<sub>x</sub>) with the topology of the projective plane requires different identifications as $`L_x`$ varies and will not be considered here. For connections between topology and graph properties, see e.g. . Here we present and analyze chromatic polynomials for the strip graph of the square lattice with $`L_y=3`$ (i.e. cross sections forming triangles) and boundary conditions of type (v) and (vi): torus and Klein bottle. We recall that unlike graphs of type (i)-(v), the Klein bottle graph (vi) cannot be embedded without self-intersection in $`^3`$. For $`L_x=m2`$ where they are well defined, the $`L_y=3`$ torus and Klein bottle graphs have $`n=L_xL_y`$ vertices, $`e=2n`$ edges, the same girth $`g`$ (length of minimum closed circuit) and number $`k_g`$ of circuits of length $`g`$, and the respective chromatic numbers $`\chi =3`$ and $`\chi =4`$.
We label a particular type of strip graph as $`G_s`$ and the specific graph of length $`L_x=m`$ repeated subgraph units, e.g. columns of squares in the case of the square strip, as $`(G_s)_m`$. If one thinks of the graph as embedded on a rectangular strip of paper, with its upper and lower sides glued together and its longitudinal ends glued with direct or reversed orientation, then $`L_x`$ is the length of this strip of paper in subgraph units. Writing
$$P((G_s)_m,q)=\underset{j=0}{\overset{n1}{}}(1)^jh_{nj}q^{nj}$$
(1)
and using the results that $`h_{nj}=\left(\genfrac{}{}{0pt}{}{e}{j}\right)`$ for $`0j<g1`$ (whence $`h_n=1`$ and $`h_{n1}=e`$) and $`h_{n(g1)}=\left(\genfrac{}{}{0pt}{}{e}{g1}\right)k_g`$, it follows that for $`m`$ greater than the above-mentioned minimal value, these $`h_j`$’s are the same for the torus and Klein bottle of each type $`G_s`$. For a given $`G_s`$, as $`m`$ increases, the $`h_{nj}`$’s for the torus and Klein bottle graphs become equal for larger $`j`$.
A generic form for chromatic polynomials for recursively defined families of graphs, of which strip graphs $`G_s`$ are special cases, is
$$P((G_s)_m,q)=\underset{j=1}{\overset{N_\lambda }{}}c_j(q)(\lambda _j(q))^m$$
(2)
where $`c_j(q)`$ and the $`N_\lambda `$ terms $`\lambda _j(q)`$ depend on the type of strip graph $`G_s`$ but are independent of $`m`$.
For an $`L_y=3`$, $`L_x=m`$ strip with (PBC<sub>y</sub>,FBC<sub>x</sub>) one has $`P(sq(L_y=3)_m,PBC_y,FBC_x,q)=q(q1)(q2)(q^36q^2+14q13)^{m1}`$, whence
$$W(sq(L_y=3),PBC_y,FBC_x,q)=(q^36q^2+14q13)^{1/3}$$
(3)
with $`=\mathrm{}`$.
In order to calculate $`P`$, one may use recursive methods based on iterative use of deletion-contraction theorems or a coloring compatibility matrix method described in . For the $`L_y=3`$ torus ($`t`$) graphs, one finds
$$P(sq(L_y=3)_m,PBC_y,PBC_x,q)=\underset{j=1}{\overset{8}{}}c_{t,j}(\lambda _{t,j})^m$$
(4)
where
$$\lambda _{t,1}=1,c_{t,1}=q^36q^2+8q1,$$
(5)
$$\lambda _{t,2}=q^36q^2+14q13,c_{t,2}=1,$$
(6)
$$\lambda _{t,3}=q1,c_{t,3}=\frac{(q1)(q2)}{2},$$
(7)
$$\lambda _{t,4}=q4,c_{t,4}=(q1)(q2),$$
(8)
$$\lambda _{t,5}=q2,c_{t,5}=q(q3),$$
(9)
$$\lambda _{t,6}=q5,c_{t,6}=\frac{q(q3)}{2},$$
(10)
$$\lambda _{t,7}=(q^27q+13),c_{t,7}=q1,$$
(11)
$$\lambda _{t,8}=(q2)^2,c_{t,8}=2(q1).$$
(12)
For the $`L_y=3`$ Klein $`(K)`$ bottle graphs one finds
$$P(sq(L_y=3)_m,PBC_y,TPBC_x,q)=\underset{j=1}{\overset{5}{}}c_{K,j}(\lambda _{K,j})^m$$
(13)
where
$$\lambda _{K,1}=\lambda _{t,1}=1,c_{K,1}=(q1),$$
(14)
$$\lambda _{K,2}=\lambda _{t,2}=q^36q^2+14q13,c_{K,2}=c_{t,2}=1,$$
(15)
$$\lambda _{K,3}=\lambda _{t,3}=q1,c_{K,3}=c_{t,3}=\frac{(q1)(q2)}{2},$$
(16)
$$\lambda _{K,4}=\lambda _{t,6}=q5,c_{K,4}=c_{t,6}=\frac{q(q3)}{2},$$
(17)
$$\lambda _{K,5}=\lambda _{t,7}=(q^27q+13),c_{K,5}=c_{t,7}=q1.$$
(18)
The terms $`\lambda _{t,j}`$, $`j=4,5,8`$ do not enter in eq. (13). We contrast this with earlier findings. For a given strip, $`N_\lambda `$ was found to be larger for (FBC<sub>y</sub>,PBC<sub>x</sub>) than (FBC<sub>y</sub>,FBC<sub>x</sub>) . For the $`L_y=3`$ square lattice strip case, $`N_\lambda =2`$ for (FBC<sub>y</sub>,FBC<sub>x</sub>) but $`N_\lambda =1`$ for (PBC<sub>y</sub>,FBC<sub>x</sub>), because of the special feature that the cross sections were complete graphs, $`K_p`$ with $`p=3`$ and hence the intersection theorem led to a factorized, monomial form for $`P`$. It was found that for a given type of lattice strip, $`N_\lambda `$ is the same for the (FBC<sub>y</sub>,PBC<sub>x</sub>) = cyclic and (FBC<sub>y</sub>,TPBC<sub>x</sub>) = Möbius topologies, although the $`c_j`$’s were, in general, different. The present results show that reversal of orientation in the identification of opposite ends of a strip can lead to a change in $`N_\lambda `$.<sup>4</sup><sup>4</sup>4A different sort of change in $`P`$, accompanied by a change in $``$, can be obtained if one considers a homogeneous recursive family and the same family with a finite inhomogeneous subgraph inserted, e.g., the “rope ladder” graphs of or two such subgraphs forming ends, viz., the $`J(H)I`$ strip graphs in .
Let $`C=_{j=1}^{N_\lambda }c_j`$. We find $`C=P(K_3,q)=q(q1)(q2)`$ for the $`L_y=3`$ torus graphs and $`C=0`$ for the $`L_y=3`$ Klein bottle graphs. The zero results from the special constraints introduced by the boundary conditions and is analogous to the fact that $`C=0`$ for the $`L_y=2`$ Möbius square strip and its homeomorphic expansions . However, not all Möbius strip graphs have $`C=0`$; for example, for the $`L_y=3`$ Möbius strips of the square and kagomé lattices, $`C=q(q1)`$ and $`C=q`$, respectively .
Chromatic zeros for the $`L_y=3`$, $`m=10`$ torus graph are shown in Fig. 1; with this value of $`m`$, the chromatic zeros for the Klein bottle graph are quite similar. The locus $``$ and the $`W`$ functions are the same for the torus and Klein bottle graph families. We find $`q_c=3`$, which, interestingly, is the same value as for the infinite 2D square lattice. The locus $``$ has support for $`Re(q)0`$ and separates the $`q`$ plane into three regions. The outermost one, region $`R_1`$, extends to infinite $`|q|`$ and includes the intervals $`q3`$ and $`q0`$ on the real $`q`$ axis. Region $`R_2`$ includes the real interval $`2q3`$ and extends upward and downward to the complex conjugate triple points on $``$ at $`q_t`$ and $`q_t^{}`$, where $`q_t2.5+1.4i`$. Region $`R_3`$ is the innermost one and includes the real interval $`0q2`$. The boundary between $`R_2`$ and $`R_3`$ curves to the right as one increases $`|Im(q)|`$, extending from $`q=2`$ upward to $`q_t`$ and downward to $`q_t^{}`$. As is evident in Fig. 1, the density of chromatic zeros on the $`R_1R_3`$ boundary near $`q=0`$ and on the $`R_3R_2`$ boundary is somewhat smaller than on the right-hand part of the $`R_2R_1`$ boundary.
In region $`R_1`$, $`\lambda _{t,2}=\lambda _{K,2}`$ is the dominant $`\lambda _j`$, so
$$W=(q^36q^2+14q13)^{1/3},qR_1.$$
(19)
The fact that this is the same as $`W`$ for the (PBC<sub>y</sub>,FBC<sub>x</sub>) case, eq. (3), is a general result. The importance of the PBC<sub>y</sub> is evident from the fact that for the same width of three squares, the strip with (FBC<sub>y</sub>,FBC<sub>x</sub>) yields a different $`W`$ .
In region $`R_2`$ $`\lambda _{t,6}=\lambda _{K,4}`$ is dominant, so
$$|W|=|q5|^{1/3},qR_2$$
(20)
(in regions other than $`R_1`$, only $`|W|`$ can be determined unambiguously ). In region $`R_3`$, $`\lambda _{t,7}=\lambda _{K,5}`$ is dominant, so
$$|W|=|q^27q+13|^{1/3},qR_3.$$
(21)
The outer boundary separating $`R_1`$ from the inner two regions is oblate, extending out to a maximum of about $`|Im(q)|2.5`$ for $`Re(q)1.5`$ (and passing through $`q=0`$ and 3). All the three points, $`q=0,2,3`$, where $``$ crosses the real $`q`$ axis, it does so vertically. The present results are in accord with the inference that for a recursive graph with regular lattice structure, a necessary and sufficient condition for $``$ to separate the $`q`$ plane into two or more regions is that it contains a global circuit, i.e. a path along a lattice direction whose length goes to infinity as $`n\mathrm{}`$; here this is equivalent to PBC<sub>x</sub>. The fact that $``$ is the same for these torus and Klein families means that none of $`\lambda _{t,j}`$, $`j=4,5,8`$ is a dominant term.
Our calculations of the zero-temperature Potts antiferromagnet partition functions (chromatic polynomials) and exponential of the entropy, $`W`$, for $`L_y=3`$ strips of the square lattice with periodic transverse and periodic and twisted periodic longitudinal boundary conditions (torus and Klein bottle graphs) thus elucidates the role that these boundary conditions and the associated topologies play; the torus and Klein bottle graphs have interestingly different chromatic polynomials, with different $`N_\lambda `$, but the $`W`$ functions and hence the boundaries $``$ are the same.
The research of R. S. was supported in part by the U. S. NSF grant PHY-97-22101. |
warning/0001/cond-mat0001072.html | ar5iv | text | # Composite Pairings in Chirally Stabilized Critical Fluids
The concept of odd-frequency superconducting ordering was conceived by Berezinskii who considered spin-triplet odd-frequency pairing as an alternate of the conventional theory <sup>3</sup>He pairing. The possibility of such unusual pairing was revived nearly one decade ago by Balatsky and Abrahams in the context of singlet superconductors and by Emery and Kivelson in the two-channel Kondo problem with the emphasis laid on the composite nature of the order parameter for odd-frequency pairing. This idea of composite pairing was further developped by Coleman et al. who proposed that heavy-fermion superconductors might involve odd-frequency triplet pairing.
One of the most important difficulty in finding some specific lattice models for realization of odd-frequency pairing stems from the fact that it requires a controlled analysis in the strong coupling regime. Such an approach is possible in some extreme limits as in one dimension where non-perturbative techniques are available. Strong odd-frequency singlet pair correlations were, in particular, identified within the bosonization approach of the one-dimensional single channel Kondo lattice in an anisotropic limit. Another candidate is the one dimensional single channel Kondo-Heisenberg lattice (KHL) which consists of a one dimensional electron gas (1DEG) interacting with an antiferromagnetic Heisenberg spin-1/2 chain by a Kondo coupling. Away from half filling, this model has a spin gap and is thus expected to exhibit enhanced pairing correlations. In particular, it has been shown recently that the dominant instabilities are of a composite nature in a certain regime of the parameters. In this letter, we shall study two different generalization of the one dimensional KHL with two channels which have no spin gap but exhibit dominant unconventional pairing instabilities. A first generalization is to consider a 1DEG coupled by a Kondo coupling to two non-interacting antiferromagnetic Heisenberg spin-1/2 chains. This model belongs to the more general class of Luttinger liquids in active environments introduced in the context of the striped physics. A second generalization consists of conduction electrons with a two-fold orbital degeneracy interacting with a periodic array of localized spins i.e. the two-channel KHL. It has been shown in Ref. that these two models away from half filling exhibit a critical phase governed by a fixed point which belongs to the class of chirally stabilized fluids. In this work, we shall compute correlation functions of physical observables by means of a Toulouse point solution and characterize the dominant instabilities of both models which turn out to be of a composite nature. Very recently, the one dimensional N-channel KHL has been investigated by Andrei and Orignac using a conformal field theory (CFT) approach. These authors found that the leading instability is of a composite pairing type when $`N4`$. In this respect, the exact solution at the Toulouse point provides an independent and physically transparent approach for the particular 2-channel KHL.
The models. The first model (referred in the following as a 1DEG in a special active environment) consists of a 1DEG away from half filling coupled symmetrically by a Kondo coupling ($`J_K>0`$) with two non-interacting antiferromagnetic ($`J_H>0`$) spin-1/2 Heisenberg chains:
$`_1`$ $`=`$ $`t{\displaystyle \underset{i}{}}(c_{i\sigma }^{}c_{i+1\sigma }+H.c.)+J_K{\displaystyle \underset{a=1}{\overset{2}{}}}{\displaystyle \underset{i}{}}𝐒_{ci}𝐒_{ai}`$ (1)
$`+`$ $`J_H{\displaystyle \underset{a=1}{\overset{2}{}}}{\displaystyle \underset{i}{}}𝐒_{ai}𝐒_{ai+1}.`$ (2)
Here, $`c_{i\sigma }`$ is the conduction electron operator at site $`i`$ with spin $`\sigma =,`$, $`\stackrel{}{S}_{ci}=c_{i\alpha }^{}\stackrel{}{\sigma }_{\alpha \beta }c_{i\beta }/2`$ ($`\stackrel{}{\sigma }`$ being the Pauli matrices) stands for the electron spin operator at site $`i`$ whereas the localized spin operator at site $`i`$ on the chain of index $`a=1,2`$ is denoted by $`𝐒_{ai}`$. In the limit $`J_Kt,J_H`$, the low energy behavior of the model can be determined using the continuum limit of the electron operator in terms of right and left moving fermion fields: $`c_{i\sigma }R_\sigma \text{e}^{ik_Fx}+L_\sigma \text{e}^{ik_Fx},x=ia_0`$ ($`k_F`$ being the Fermi momentum and $`a_0`$ the lattice spacing). The continuum description of the conduction spin density operator $`𝐒_c`$ can then be expressed in terms of a spin current $`𝐉_{0R,L}`$ belonging to the SU(2)<sub>1</sub> Kac-Moody (KM) algebra and a bosonic field $`\mathrm{\Phi }_c=\mathrm{\Phi }_{cL}+\mathrm{\Phi }_{cR}`$ that accounts for the charge degrees of freedom: $`𝐒_c𝐉_{0R}+𝐉_{0L}+\mathrm{cos}(2k_Fx+\sqrt{2\pi }\mathrm{\Phi }_c)𝐧_0`$, $`𝐧_0`$ being the staggered magnetization. In the same way, the spin densities of the surface chains are represented as (see Refs. ): $`𝐒_a(x)=𝐉_a(x)+(1)^{x/a_0}𝐧_a(x)`$ where $`𝐉_a=𝐉_{aR}+𝐉_{aL}`$ and $`𝐧_a`$ are respectively the uniform and staggered parts of the magnetization. The chiral spin currents $`𝐉_{aR,L}`$ belong also to the SU(2)<sub>1</sub> KM algebra. With this low energy description, the Hamiltonian density of the lattice model (2) for incommensurate filling reads as follows in the continuum limit:
$`_1`$ $``$ $`{\displaystyle \frac{v_F}{2}}\left(\left(_x\mathrm{\Phi }_c\right)^2+\left(_x\mathrm{\Theta }_c\right)^2\right)+{\displaystyle \frac{2\pi v_F}{3}}\left(𝐉_{0R}^2+𝐉_{0L}^2\right)`$ (3)
$`+`$ $`{\displaystyle \frac{2\pi v_H}{3}}{\displaystyle \underset{a=1}{\overset{2}{}}}\left(𝐉_{aR}^2+𝐉_{aL}^2\right)+g\left(𝐉_{0R}𝐈_L+𝐉_{0L}𝐈_R\right)`$ (4)
where we have neglected all oscillatory, marginal irrelevant contributions as well as current-current interactions of the same chirality. As shown in Ref., at the strong coupling fixed point, these latter interactions only lead to renormalization of velocities and logarithmic corrections. In Eq. (4), $`\mathrm{\Theta }_c=\mathrm{\Phi }_{cL}\mathrm{\Phi }_{cR}`$ is the dual charge field, $`v_F`$ the Fermi velocity, $`v_H`$ the spin velocity of the magnetic environment, and $`gJ_Ka_0>0`$ is the interacting coupling constant; $`𝐈_{R,L}=𝐉_{1R,L}+𝐉_{2R,L}`$ is a SU(2)<sub>2</sub> KM current being the sum of two SU(2)<sub>1</sub> KM currents.
The second model considered in this letter is the one-dimensional two-channel KHL:
$`_2`$ $`=`$ $`t{\displaystyle \underset{i}{}}{\displaystyle \underset{a=1}{\overset{2}{}}}(c_{ia\sigma }^{}c_{i+1a\sigma }+H.c.)+J_H{\displaystyle \underset{i}{}}𝐒_{0i}𝐒_{0i+1}`$ (5)
$`+`$ $`J_K{\displaystyle \underset{a=1}{\overset{2}{}}}{\displaystyle \underset{i}{}}𝐒_{aci}𝐒_{0i}`$ (6)
where now the electron operator $`c_{ia\sigma }`$ has a channel index $`a=1,2`$ and interacts by a Kondo coupling ($`J_K>0`$) with a periodic array of localized spins $`𝐒_{0i}`$. Away from half filling, the continuum limit of the Hamiltonian (6) proceeds in the same way as in the first model and with the same degree of approximations than in Eq. (4), one has:
$`_2`$ $``$ $`{\displaystyle \frac{v_F}{2}}{\displaystyle \underset{a=1}{\overset{2}{}}}\left(\left(_x\mathrm{\Phi }_{ac}\right)^2+\left(_x\mathrm{\Theta }_{ac}\right)^2\right)+{\displaystyle \frac{2\pi v_H}{3}}\left(𝐉_{0R}^2+𝐉_{0L}^2\right)`$ (7)
$`+`$ $`{\displaystyle \frac{2\pi v_F}{3}}{\displaystyle \underset{a=1}{\overset{2}{}}}\left(𝐉_{aR}^2+𝐉_{aL}^2\right)+g\left(𝐉_{0R}𝐈_L+𝐉_{0L}𝐈_R\right)`$ (8)
where $`\mathrm{\Phi }_{ac}`$ is a bosonic field ($`\mathrm{\Theta }_{ac}`$ being the dual field) associated with charge fluctuations in the ath channel. The chiral SU(2)<sub>1</sub> currents $`𝐉_{aR,L}`$ correspond to the chiral uniform part of the electron spin density whereas $`𝐉_{0R,L}`$ are the chiral SU(2)<sub>1</sub> currents of the local moments.
Toulouse point solution. The next step of the approach is to use the representation of two SU(2)<sub>1</sub> currents in terms of four Majorana fermions $`\xi ^0`$ and $`\stackrel{}{\xi }`$: $`𝐈_{R,L}=\text{i}\stackrel{}{\xi }_{R,L}\stackrel{}{\xi }_{R,L}/2,𝐉_{1R,L}𝐉_{2R,L}=\text{i}\stackrel{}{\xi }_{R,L}\xi _{R,L}^0`$. The Hamiltonians (4,8) share then the same structure: $`_{1,2}=_c^0+^0(\xi ^0)+\overline{}`$ where $`_c^0`$ is the free Hamiltonian of the charge bosonic fields and $`^0(\xi ^0)`$ is the free Hamiltonian of the massless Majorana fermion $`\xi ^0`$. All nontrivial physics is incorporated in $`\overline{}`$ which separates into two commuting and chirally asymmetric parts: $`\overline{}=\overline{}_1+\overline{}_2,[\overline{}_1,\overline{}_2]=0`$ with
$$\overline{}_1=\frac{\pi v_1}{2}𝐈_R^2+\frac{2\pi v_0}{3}𝐉_{0L}^2+g𝐉_{0L}𝐈_R$$
(9)
with $`v_1=v_H,v_0=v_F`$ for the 1DEG in a special active environment and the reverse for the two-channel KHL; $`\overline{}_2`$ is obtained from $`\overline{}_1`$ by inverting chiralities of all the currents. For an antiferromagnetic Kondo coupling ($`g>0`$), the interacting part of $`\overline{}`$ is marginally relevant. Naively one will thus expect that the model will enter a massive strong coupling region. However, due to the chiral asymmetry of $`\overline{}`$, the effective interaction flows towards a conformally invariant intermediate fixed point with a smaller central charge than in the ultraviolet. Stated differently, the interaction produces a mass gap in some sectors but there are still some degrees of freedom that remain critical in the infrared (IR). The nature of the fixed point of chirally asymmetric current-current models with the structure (9) have been determined by Andrei et al. by a combination of CFT and Bethe ansatz techniques. This fixed point belongs to the class of chirally stabilized fluids introduced in Ref. describing the IR behavior of one-dimensional interacting chiral fermions. A simpler way to identify the critical degrees of freedom of model (9) in the IR limit can be done by means of a Toulouse point approach as in the two-channel Kondo model. By allowing anisotropic interaction ($`gg_{},g_{}`$), the U(1) version of the Hamiltonian $`\overline{}`$ can be mapped onto free Majorana fermions for a special value of the interaction $`g_{}^{}`$. Although the position of the solvable point is non-universal, the Toulouse limit solution captures the physical and universal properties of the models including the SU(2) case. The details of the approach can be found in Refs. and we shall now briefly review it to fix the notations. The starting point of the Toulouse solution is the Abelian bosonization of all SU(2)<sub>1</sub> currents $`𝐉_a,a=0,1,2`$ in terms of massless bosonic fields $`\phi _a`$: $`J_{aR,L}^z=(1/\sqrt{2\pi })_x\phi _{aR,L}`$, $`J_{aR,L}^+=(1/2\pi a_0)\text{e}^{\text{i}\sqrt{8\pi }\phi _{aR,L}}`$. Introducing the symmetric combination of the fields: $`\phi _+=(\phi _1+\phi _2)/\sqrt{2}`$, the SU(2)<sub>2</sub> current $`𝐈`$ writes: $`I_{R,L}^z=(1/\sqrt{\pi })_x\phi _{+R,L}`$, $`I_{R,L}^+=(\text{i}/\sqrt{\pi a_0})\xi _{R,L}^3\kappa \text{e}^{\text{i}\sqrt{4\pi }\phi _{+R,L}}`$. A fermionic zero-mode operator $`\kappa `$ has been introduced to ensure the proper anticommutation relations. The following canonical transformation is then performed:
$`\phi _0`$ $`=`$ $`\sqrt{2}\overline{\mathrm{\Phi }}_2\overline{\mathrm{\Phi }}_1,\phi _+=\sqrt{2}\overline{\mathrm{\Phi }}_1\overline{\mathrm{\Phi }}_2`$ (10)
$`\vartheta _0`$ $`=`$ $`\sqrt{2}\overline{\mathrm{\Theta }}_2+\overline{\mathrm{\Theta }}_1,\vartheta _+=\sqrt{2}\overline{\mathrm{\Theta }}_1+\overline{\mathrm{\Theta }}_2`$ (11)
where $`\vartheta _0`$ and $`\vartheta _+`$ (respectively $`\overline{\mathrm{\Theta }}_1`$ and $`\overline{\mathrm{\Theta }}_2`$) are the dual fields associated with $`\phi _0`$ and $`\phi _+`$ (respectively $`\overline{\mathrm{\Phi }}_1`$ and $`\overline{\mathrm{\Phi }}_2`$). For a special positive value (Toulouse point) of $`g_{},(g_{}^{}=4\pi (v_0+v_1)/3)`$, the arguments of the interacting terms become those of free fermions so that they can be refermionized further with the introduction of a pair of Majorana fields ($`\eta ,\zeta `$) using the correspondence: $`\eta _{R,L}+\text{i}\zeta _{R,L}=(\kappa /\sqrt{\pi a_0})\text{e}^{\pm \text{i}\sqrt{4\pi }\overline{\mathrm{\Phi }}_{2R,L}}`$. One finally ends with:
$`\overline{}`$ $`=`$ $`{\displaystyle \frac{u_1}{2}}\left[(_x\overline{\mathrm{\Phi }}_1)^2+(_x\overline{\mathrm{\Theta }}_1)^2\right]{\displaystyle \frac{\text{i}u_2}{2}}\left[\zeta _R_x\zeta _R\zeta _L_x\zeta _L\right]`$ (12)
$``$ $`{\displaystyle \frac{\text{i}v_1}{2}}\left[\xi _R^3_x\xi _R^3\xi _L^3_x\xi _L^3\right]{\displaystyle \frac{\text{i}u_2}{2}}\left[\eta _R_x\eta _R\eta _L_x\eta _L\right]`$ (13)
$`+`$ $`\text{i}m\left[\xi _R^3\eta _L\eta _R\xi _L^3\right]`$ (14)
where $`m=g_{}/2\pi a_0`$, $`u_1=(2v_1v_0)/3`$, and $`u_2=(2v_0v_1)/3`$. The Toulouse solution is stable provided that all velocities are positive i.e. $`1/2v_0/v_12`$. The first two terms in Eq. (14) describe decoupled free massless bosonic and Majorana fields, $`\overline{\mathrm{\Phi }}_1`$ and $`\zeta `$, contributing to criticality with the central charge: $`c=3/2.`$ The remaining part of the Hamiltonian (14) has a spectral gap $`m`$ and describes hybridization of the Majorana $`\xi ^3`$ and $`\eta `$ fermions with different chiralities. Adding the contribution of the excitations that do not participate in the interaction, the total central charge in the IR of the 1DEG in a special active environment (respectively the two-channel KHL) is $`c=3`$ (respectively $`c=4`$). Apart from the charge degrees of freedom, the elementary excitations of the models in the IR are of different nature: The magnetic excitations are spinons defined as $`\sqrt{\pi /2}`$ kinks of the bosonic field $`\overline{\mathrm{\Phi }}_1`$ describing an effective S=1/2 Heisenberg chain and nonmagnetic singlet excitations associated with the two massless Majorana fermions $`\xi ^0,\zeta `$ (critical Ising degrees of freedom) referred as pseudocharge degrees of freedom in Refs.. These $`Z_2`$ excitations play a crucial role in the nontrivial IR physical properties of the models as we shall see now.
Physical properties of the 1DEG in a special active environment. The Green’s function for the right-left moving fermions can be computed using the bosonic representation: $`(R,L)_\sigma =(\kappa _\sigma /\sqrt{2\pi a_0})\text{e}^{\text{i}\tau _\sigma \pi /4}\text{e}^{\pm \text{i}\sqrt{4\pi }\phi _{\sigma R,L}}`$ where $`\tau _,=\pm 1`$, $`\phi _{\sigma R,L}=(\mathrm{\Phi }_{cR,L}+\tau _\sigma \phi _{0R,L})/\sqrt{2}`$, and $`\kappa _\sigma `$ are Klein-factors to insure the correct anticommutation relations between fermion fields of different spin index. Using the Toulouse basis (11), one then obtains at the IR fixed point the estimate:
$`R_\sigma (x,\tau )R_\sigma ^{^{}}^{}(0,0)`$ (15)
$`{\displaystyle \frac{\delta _{\sigma ,\sigma ^{^{}}}}{\left(v_F\tau \text{i}x\right)^{1/2}\left(u_1\tau +\text{i}x\right)^{1/2}\left(u_2\tau \text{i}x\right)}}`$ (16)
which means that the system displays non-Fermi liquid properties; the left-moving Green’s function is obtained from Eq. (16) by the substitution: $`\text{i}\text{i}`$. The most interesting physical quantities are the correlation functions of the various possible order parameters. The spin-spin correlation functions have been computed in Ref. and the slowest ones are the staggered parts that decay as $`x^{3/2}`$. Using the previous Abelian bosonization of the chiral fermions and the Toulouse basis (11), we find in the electronic sector the following representation for the charge density wave (CDW), spin density wave (SDW), singlet (SS) and triplet (ST) superconducting order parameters in terms of the different critical fields at the IR fixed point:
$`𝒪_{CDW}=L_\sigma ^{}R_\sigma \text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Phi }_c}\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)\zeta _L\zeta _R`$ (17)
$`\stackrel{}{𝒪}_{SDW}=L_\alpha ^{}\stackrel{}{\sigma }_{\alpha \beta }R_\beta \text{i}\text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Phi }_c}\zeta _R\zeta _L`$ (18)
$`[\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right),\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right),\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)]`$ (19)
$`𝒪_{SS}=\text{i}L_\alpha \left(\sigma ^y\right)_{\alpha \beta }R_\beta \text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Theta }_c}\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)\zeta _R\zeta _L`$ (20)
$`\stackrel{}{𝒪}_{TS}=\text{i}L_\alpha \left(\stackrel{}{\sigma }\sigma ^y\right)_{\alpha \beta }R_\beta \text{i}\text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Theta }_c}\zeta _R\zeta _L`$ (21)
$`[\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right),\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right),\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)].`$ (22)
All these conventional order parameters have thus the same scaling dimension ($`2`$) and the corresponding pair correlation functions decay as $`x^4`$ i.e. much faster than in the one dimensional metals. One should notice that the $`Z_2`$ pseudocharge excitations contribute in (22) through the density energy operator $`\text{i}\zeta _R\zeta _L`$. In a critical Ising model, there are other primary fields (order and disorder parameters) that have a smaller scaling dimension ($`1/8`$) which are highly nonlocal in terms of the original Majorana fermion. This leads us to investigate the possibility that the ground state might be characterized by composite order parameters: staggered odd-frequency singlet pairing (c-SP) and staggered composite CDW (c-CDW). At the IR fixed point, we find the following correspondence for these order parameters:
$`𝒪_{cSP}`$ $`=`$ $`\stackrel{}{𝒪}_{TS}\left(𝐧_1+𝐧_2\right)\text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Theta }_c}\mu _0\mu _4`$ (23)
$`𝒪_{cCDW}`$ $`=`$ $`\stackrel{}{𝒪}_{SDW}\left(𝐧_1+𝐧_2\right)\text{e}^{\text{i}\sqrt{2\pi }\mathrm{\Phi }_c}\mu _0\mu _4`$ (24)
where $`\mu _0,\mu _4`$ (respectively $`\sigma _0,\sigma _4`$) are the Ising disorder (respectively order) parameters associated with the Majorana fermions $`\xi ^0,\zeta `$. ¿From Eq. (24), we deduce the leading asymptotics of the correlation function corresponding to the composite order parameters:
$`𝒪_{composite}(x,\tau )𝒪_{composite}^{}(0,0)`$ (25)
$`{\displaystyle \frac{1}{\left(v_F^2\tau ^2+x^2\right)^{1/2}\left(v_1^2\tau ^2+x^2\right)^{1/8}\left(u_2^2\tau ^2+x^2\right)^{1/8}}}`$ (26)
so that the composite order parameters induce the dominant instabilities and have enhanced long-range coherence. These fluctuations can be made even more enhanced upon switching on a very small ferromagnetic ($`g_{12}<0`$) exchange interaction between the spin chains: $`𝒪_{12}=g_{12}𝐒_1𝐒_2`$. Indeed, at the IR fixed point, this perturbation mostly affects the pseudocharge sector leaving intact the magnetic one: $`𝒪_{12}g_{12}\text{i}\xi _R^0\xi _L^0`$. The Majorana fermion $`\xi ^0`$ acquires a positive mass i.e. its associated Ising model is in the disorder phase: $`\mu _00`$ and the $`Z_2`$ symmetry with respect to the interchange of the two chains is now broken. All other degrees of freedom remain critical so that the pair correlation of the conventional order parameters (22) still decay as $`x^4`$ whereas from Eq. (24) one immediately observes that the composite ones decay now as $`x^{5/4}`$ instead of $`x^{3/2}`$ at $`g_{12}=0`$.
The two-channel KHL case. The determination of the ground-state physical properties of the two-channel KHL proceeds in the same way as in the previous model. We first use the Abelian bosonisation of the right-left moving fermions: $`(R,L)_{a\sigma }=(\kappa _{a\sigma }/\sqrt{2\pi a_0})\text{e}^{\text{i}\tau _\sigma \pi /4}\text{e}^{\pm \text{i}\sqrt{4\pi }\phi _{a\sigma R,L}}`$ where $`\phi _{a\sigma R,L}=(\mathrm{\Phi }_{acR,L}+\tau _\sigma \phi _{aR,L})/\sqrt{2}`$, and $`\kappa _{a\sigma }`$ are some Klein-factors. The Green’s functions of the chiral fermions have a similar structure than in Eq. (16) apart from they fall now with the distance with a power $`5/4`$. The leading asymptotics of the correlation function of the localized spins have been computed in Ref. and in particular the staggered part decays as $`x^3`$. Using the Toulouse basis (11) and introducing the charge and relative charge bosonic fields $`\mathrm{\Phi }_{\pm cR,L}=(\mathrm{\Phi }_{1cR,L}\pm \mathrm{\Phi }_{2cR,L})/\sqrt{2}`$, the order parameters for the electronic degrees of freedom can also be expressed in terms of the critical fields at the IR fixed point. We find the following correspondence for the conventional order parameters:
$`𝒪_{CDW}=L_{a\sigma }^{}R_{a\sigma }\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Phi }_{+c}}\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)_\mathrm{\Phi }`$ (27)
$`\stackrel{}{𝒪}_{SDW}=L_{a\alpha }^{}\stackrel{}{\sigma }_{\alpha \beta }R_{a\beta }\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Phi }_{+c}}`$ (28)
$`[\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right)_\mathrm{\Phi }^{},\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right)_\mathrm{\Phi }^{},\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)_\mathrm{\Phi }]`$ (29)
$`𝒪_{SS}=\text{i}L_{a\alpha }\left(\sigma ^y\right)_{\alpha \beta }R_{a\beta }\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Theta }_{+c}}\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)_\mathrm{\Theta }^{}`$ (30)
$`\stackrel{}{𝒪}_{TS}=\text{i}L_{a\alpha }\left(\stackrel{}{\sigma }\sigma ^y\right)_{\alpha \beta }R_{a\beta }\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Theta }_{+c}}`$ (31)
$`[\mathrm{cos}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right)_\mathrm{\Theta },\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Theta }}_1\right)_\mathrm{\Theta },\mathrm{sin}\left(\sqrt{2\pi }\overline{\mathrm{\Phi }}_1\right)_\mathrm{\Theta }^{}]`$ (32)
where $`_\mathrm{\Phi }=\mu _0\mu _4\mathrm{cos}\left(\sqrt{\pi }\mathrm{\Phi }_c\right)+\text{i}\sigma _0\sigma _4\mathrm{sin}\left(\sqrt{\pi }\mathrm{\Phi }_c\right)`$ and $`_\mathrm{\Theta }`$ is obtained from $`_\mathrm{\Phi }`$ by replacing $`\mathrm{\Phi }_c`$ by its dual field $`\mathrm{\Theta }_c`$. The scaling dimension of all these conventional order parameters is $`5/4`$ and their corresponding pair correlation functions decay thus as $`x^{5/2}`$ i.e. have no enhanced long-range coherence. Notice, however, that this decay is much slower than in the 1DEG in a special active environment case. Similarly, we find the following description at the IR fixed point for the staggered composite order parameters:
$`𝒪_{cSP}=\stackrel{}{𝒪}_{TS}𝐧_0\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Theta }_{+c}}_\mathrm{\Theta }`$ (33)
$`𝒪_{cCDW}=\stackrel{}{𝒪}_{SDW}𝐧_0\text{e}^{\text{i}\sqrt{\pi }\mathrm{\Phi }_{+c}}`$ (34)
$`\left(\mu _0\mu _4\mathrm{cos}\left(\sqrt{\pi }\mathrm{\Phi }_c\right)+3\text{i}\sigma _0\sigma _4\mathrm{sin}\left(\sqrt{\pi }\mathrm{\Phi }_c\right)\right)`$ (35)
so that the scaling dimension of these operators is $`3/4`$ and the pair correlation function of the composite order parameters is still given by Eq. (26). The dominant instabilities of the two-channel KHL are thus of a composite nature as in the 1DEG in a special active environment. Finally, one should note that the scaling dimensions of the physical observables of the two-channel KHL obtained by the Toulouse point solution coincide with those of the recent CFT approach. In this respect, we stress that, in contrast with the latter approach, the Toulouse point solution unables us to take into account the inherent velocities anisotropy in the problem.
In summary, we have shown that the two generalizations of the one dimensional KHL considered in this letter exhibit a nontrivial non-Fermi liquid low-temperature phase belonging to the class of chirally stabilized fluids with strong enhanced staggered composite pairing correlations. |
warning/0001/astro-ph0001429.html | ar5iv | text | # Comptonization of the Cosmic Microwave Background by Relativistic Plasma
## 1 Introduction
The Sunyaev-Zeldovich (SZ) effect has proven to be an important tool for cosmology and the study of clusters of galaxies. It measures the total thermal energy content of the electron population along a line of sight and does not depend on spectral features of the underlying electron distributions, as long as they are non-relativistic. For a review see Sunyaev & Zeldovich (1980).
Deviations in the spectral shape of the SZ effect occur for higher electron temperatures or energies, increasing from the mildly, through the trans- to the ultra-relativistic regime. Approximation for relativistic corrections to the thermal SZ effect are given in the literature. For reviews see Rephaeli (1995a), Fargion & Salis (1998), Molnar & Birkinshaw (1999), Birkinshaw (1999), and Sazonov & Sunyaev (2000).
The thermal SZ effect (thSZ) is therefore a calorimeter for heat releasing processes in the universe, such as the gravitational compression of matter during structure formation, and the energy output of galaxies. A very violent form of energy release of galaxies are the outflows of relativistic plasma from active galactic nuclei. This plasma fills large volumes, the radio lobes, and thereby forms very low density cavities filled with relativistic particles and magnetic fields (for evidence for these cavities see: Böhringer et al 1995, Carilli & Harris 1996, Carilli & Barthel 1996, Clarke et al. 1997; McNamara et al. 2000). The replaced inter-galactic medium (IGM) is heated by this process (e.g. Kaiser & Alexander 1999). It is one goal of this work to estimate the amount of this energy and its imprint on the CMB (see also Enßlin et al. 1998b; Yamada et al. 1999).
Also the radio plasma is able to Comptonize the CMB. A search for the SZ effect in radio lobes was carried out by McKinnon et al. (1991). After being released from the radio galaxy, the relativistic electrons with the highest energies cool rapidly due to synchrotron and inverse Compton (IC) losses, letting the radio lobes fade away after the AGN stops its activity. Patches of such remnant, undetectable radio plasma were named ‘radio ghosts’ (Enßlin 1999). However, the less energetic electrons may remain relativistic in such an environment for cosmological times. Compton scattering off these electrons removes photons from the CMB Planck spectrum and scatters them to much higher energies, leading to an SZ effect with a characteristic spectral signature.
Since radio ghosts are a still speculative ingredient of the IGM, we discuss different strategies to reveal their presence. We further discuss their possible influence on the SZ effect in clusters of galaxies.
The paper is organized as follows. In Sect. 2 we discuss the theory of transrelativistic inverse Compton (IC) scattering. In Sect. 3 we examine the possible role radio ghosts can have, and estimate their contribution to the CMB-Comptonization. In Sect. 4 we discuss the possible influence remnant radio plasma in the ICM can have on the determination of the Hubble parameter via the SZ effect (Sect. 4.1). Further SZ measurements of cluster of galaxies are proposed as a tool to test for non-thermal electrons in the intra-cluster medium (Sect. 4.2). Conclusions can be found in Sect. 5.
We assume an Einstein de Sitter (EdS) cosmology and $`H_\mathrm{o}=50h_{50}\mathrm{km}\mathrm{s}^1\mathrm{Mpc}^1`$. All IC calculations are done in the optically thin limit.
## 2 Comptonization of the CMB
### 2.1 Phenomenology of SZ-Effects
We distinguish here three processes which change the spectral shape of the CMB by inverse Compton scattering:
* the thermal Sunyaev-Zel’dovich effect (thSZ),
* the kinetic Sunyaev-Zel’dovich effect (kSZ),
* up-scattering of photons by relativistic electrons (rSZ).
The latter happens whenever a CMB photon gets scattered off a relativistic electron. An electron with gamma factor $`\gamma `$ increases the frequency $`\nu `$ of a scattered photon ($`\nu ^{}`$) on average by a factor $`\nu ^{}/\nu =\frac{4}{3}\gamma ^2\frac{1}{3}`$. The photon is therefore scattered to much higher energies and is effectively removed from the spectral range of the CMB. We consider a photon to be removed from the CMB, if its energy is increased by about one order of magnitude or greater, requiring $`\gamma >3`$. Relativistic electrons with lower $`\gamma `$-factors still produce a substantial energy gain of the photon, which also depopulates the Rayleigh-Jeans side of the CMB, but these electrons are less effective scatterers compared to electrons of higher energies. An accurate way to treat this case is given in Sect. 2.2.
Using the usual Kompaneets approximation and neglecting for the moment those electrons with $`\gamma <3`$, the spectral distortion produced by the three processes under considerations are (see e.g. Rephaeli 1995b)
$$\delta i(x)=g(x)y_{\mathrm{gas}}h(x)\overline{\beta }_{\mathrm{gas}}\tau _{\mathrm{gas}}i(x)\tau _{\mathrm{rel}},$$
(1)
where the CMB-blackbody spectrum is
$$I_\nu =i_0i(x)=i_0\frac{x^3}{e^x1},$$
(2)
with $`x=h\nu /kT`$ and $`i_0=2(kT_{\mathrm{cmb}})^3/(hc)^2`$. For comparison: the peak of the spectrum is at $`x=2.82`$, and the up-coming Planck satellite will measure the CMB signal at $`x=`$ 1.8, 2.5, 3.8, 6.2, 9.6, and 15 (Bersanelli et al. 1996; Puget 1998; see also Tab. 2). The thSZ-distortions, which have a spectral shape given by
$$g(x)=\frac{x^4e^x}{(e^x1)^2}\left(x\frac{e^x+1}{e^x1}4\right),$$
(3)
depend in their strength on the Comptonization parameter
$$y_{\mathrm{gas}}=\frac{\sigma _\mathrm{T}}{m_\mathrm{e}c^2}𝑑ln_{\mathrm{e},\mathrm{gas}}kT_\mathrm{e}.$$
(4)
The line-of-sight integral extends from the last scattering surface of the cosmic background radiation at $`z=1100`$ to the observer at $`z=0`$. The kSZ-distortions have the spectral shape
$$h(x)=\frac{x^4e^x}{(e^x1)^2}$$
(5)
and depend on $`\overline{\beta }_{\mathrm{gas}}`$, the average line-of-sight streaming velocity of the thermal gas ($`v_{\mathrm{gas}}=\beta _{\mathrm{gas}}c`$, $`\beta _{\mathrm{gas}}>0`$ if gas is approaching the observer), and the Thomson optical depth
$$\overline{\beta }_{\mathrm{gas}}\tau _{\mathrm{gas}}=\sigma _\mathrm{T}𝑑ln_{\mathrm{e},\mathrm{gas}}\beta _{\mathrm{gas}}.$$
(6)
Finally, the optical depth of relativistic electrons
$$\tau _{\mathrm{rel}}=\sigma _\mathrm{T}𝑑ln_{\mathrm{e},\mathrm{rel}}$$
(7)
determines the fraction of the number of photons, which are effectively removed from the CMB.
### 2.2 Transrelativistic Thomson Scattering
We derive here the exact formulae for the photon redistribution function of the transrelativistic SZ effect in the Thomson regime, which can be applied to get the relativistically correct thSZ, but also the rSZ for any arbitrary isotropic electron distribution. Such a formula has already been given in the literature (Rephaeli 1995a), but not in the compact form provided here (see Eq. (21)). We further give the analytically exact photon redistribution function for a power-law electron distribution.
The change in intensity due to IC-scattering of a blackbody photon distribution at an isotropic distribution of electrons in the optical thin limit is given by
$$\delta i(x)=\left(j(x)i(x)\right)\tau ,$$
(8)
where $`\tau `$ is the Thomson optical depth
$$\tau =\sigma _T𝑑ln_\mathrm{e}.$$
(9)
$`i(x)\tau `$ is the flux scattered to other frequencies, and $`j(x)\tau `$ is the flux scattered from other frequencies to $`x=h\nu /(kT)`$. In Eq. (1) we assumed that $`j(x)i(x)`$ in the region of interest ($`x<10`$) for ultra-relativistic electrons. In the following we drop this approximation.
Eq. (8) can be re-written to include a $`y`$-like parameter
$$\delta i(x)=\stackrel{~}{g}(x)\stackrel{~}{y},$$
(10)
where
$`\stackrel{~}{y}`$ $`=`$ $`{\displaystyle \frac{\sigma _T}{m_\mathrm{e}c^2}}{\displaystyle 𝑑ln_\mathrm{e}k\stackrel{~}{T}_\mathrm{e}}`$ (11)
$`k\stackrel{~}{T}_\mathrm{e}`$ $`=`$ $`{\displaystyle \frac{P_\mathrm{e}}{n_\mathrm{e}}}`$ (12)
$`\stackrel{~}{g}(x)`$ $`=`$ $`\left(j(x)i(x)\right){\displaystyle \frac{m_\mathrm{e}c^2}{k\stackrel{~}{T}_\mathrm{e}}}`$ (13)
$`k\stackrel{~}{T}_\mathrm{e}`$ $`=`$ $`{\displaystyle \frac{𝑑ln_\mathrm{e}k\stackrel{~}{T}_\mathrm{e}}{𝑑ln_\mathrm{e}}}`$ (14)
with $`P_\mathrm{e}`$ the electron pressure (see Eq. (69)). The definition of the pseudo-temperature $`k\stackrel{~}{T}_\mathrm{e}`$ is equal to the thermodynamic temperature in the case of a thermal electron distribution (Eq. (24)).
If one introduces the photon redistribution function $`P(t)dt`$, which gives the probability that the frequency of a scattered photon is shifted by a factor $`t`$, one can express the scattered spectrum as
$$j(x)=_0^{\mathrm{}}𝑑tP(t)i(x/t),$$
(15)
where $`P(t)`$ gives the probability that a photon is scattered to a frequency $`t`$ times its original frequency. For a given electron spectrum $`f_\mathrm{e}(p)dp`$, where $`p=\beta _\mathrm{e}\gamma _\mathrm{e}`$ is the normalized electron momentum, and where $`𝑑pf_\mathrm{e}(p)=1`$, the photon redistribution function can be written as
$$P(t)=_0^{\mathrm{}}𝑑pf_\mathrm{e}(p)P(t;p).$$
(16)
$`P(t;p)`$ is the redistribution function for a momoenergetic electron distribution. $`P(t;p)`$ can be derived following Wright’s (1979) kinematical considerations of the IC scattering in the Thomson regime ($`h\nu \gamma _\mathrm{e}m_\mathrm{e}c^2`$). Scattering of a photon from an isotropic photon field has a probability distribution of the angle $`\theta `$ between photon and electron direction in the electron’s rest frame given by
$$f(\mu )d\mu =[2\gamma _\mathrm{e}^4(1\beta _\mathrm{e}\mu )^3]^1d\mu ,$$
(17)
where $`\mu =\mathrm{cos}\theta `$. The frequency shift from $`\nu `$ to $`\nu ^{}`$ of the photon after scattering into an angle $`\theta ^{}`$ is
$$t=\frac{\nu ^{}}{\nu }=\frac{1+\beta _\mathrm{e}\mu ^{}}{1\beta _\mathrm{e}\mu }.$$
(18)
The probability distribution of $`\mu ^{}=\mathrm{cos}\theta ^{}`$ for a given $`\mu `$ is
$$g(\mu ^{};\mu )d\mu ^{}=\frac{3}{8}[1+\mu ^2\mu ^2+\frac{1}{2}(1\mu ^2)(1\mu ^2)]d\mu ^{}$$
(19)
(Chandrasekhar 1950). The probability for a shift $`t`$ given the electron velocity is
$$P(t;p)=𝑑\mu f(\mu )g(\mu ^{};\mu )\left(\frac{\mu ^{}}{t}\right).$$
(20)
$`\mu ^{}(t,\beta _\mathrm{e},\mu )`$ is given by (18) and the range of integration is given by the conditions $`1<\mu <1`$ and $`1<\mu ^{}<1`$. The integrand in (20) is a rational function of $`\mu `$ and therefore analytically integrable:
$`P(t;p)={\displaystyle \frac{3|1t|}{32p^6t}}\left[1+(10+8p^2+4p^4)t+t^2\right]+`$ (21)
$`{\displaystyle \frac{3(1+t)}{8p^5}}\left[{\displaystyle \frac{3+3p^2+p^4}{\sqrt{1+p^2}}}{\displaystyle \frac{3+2p^2}{2p}}(2\mathrm{arcsinh}(p)|\mathrm{ln}(t)|)\right]`$
The maximal frequency shift is given by
$$|\mathrm{ln}(t)|2\mathrm{arcsinh}(p),$$
(22)
and therefore $`P(t;p)=0`$ for $`|\mathrm{ln}(t)|>2\mathrm{arcsinh}(p)`$. Similar expressions for the photon redistribution function using different variables can be found in the literature (Rephaeli 1995a; Fargion et al. 1996; Enßlin & Biermann 1998; Sazonov & Sunyaev 2000). We verified analytically, that $`P(t;p)`$ has the proper normalization ($`𝑑tP(t;p)=1`$) and checked numerically that it also reproduces the known average energy gain of the scattered photons ($`𝑑tP(t;p)(t1)=\frac{4}{3}p^2`$).
If one is interested in the large frequency shift ($`t1`$) due to IC-scattering by ultra-relativistic electrons it is convenient to use the logarithmic frequency shift $`s=\mathrm{ln}(t)`$. The photon redistribution function then reads
$$P_s(s;p)ds=P(e^s;p)e^sds,$$
(23)
and $`P_s(s;p)=0`$ for $`|s|>s_{\mathrm{max}}(p)=2\mathrm{arcsinh}(p)`$. $`P_s(s;p)`$ is plotted in Fig 2.
We are interested in three distinct forms of the momentum distribution of the scattering electrons:
* a thermal electron distribution:
$$f_{\mathrm{e},\mathrm{th}}(p;\beta _{\mathrm{th}})=\frac{\beta _{\mathrm{th}}}{K_2(\beta _{\mathrm{th}})}p^2\mathrm{exp}(\beta _{\mathrm{th}}\sqrt{1+p^2}),$$
(24)
with $`K_\nu `$ denoting the modified Bessel function of the second kind (Abramowitz & Stegun 1965), introduced here for proper normalization, and $`\beta _{\mathrm{th}}=m_\mathrm{e}c^2/kT_\mathrm{e}`$ the normalized thermal beta-parameter.
* a power-law distribution between $`p_1`$ and $`p_2`$, as a model for a cosmic ray (CR) electron population:
$$f_{\mathrm{e},\mathrm{cr}}(p;\alpha ,p_1,p_2)=\frac{(\alpha 1)p^\alpha }{p_1^{1\alpha }p_2^{1\alpha }}$$
(25)
* a thermal distribution with a high-energy power-law tail, as a model for a thermal population modified by in-situ particle acceleration:
$`f_{\mathrm{e},\mathrm{th}\&\mathrm{cr}}(p;\beta _{\mathrm{th}},\alpha ,p_1,p_2)=C_\mathrm{e}(\beta _{\mathrm{th}},\alpha ,p_1,p_2)\times `$ (26)
$`\left\{\begin{array}{cc}f_{\mathrm{e},\mathrm{th}}(p;\beta _{\mathrm{th}})\hfill & ;pp_1\hfill \\ f_{\mathrm{e},\mathrm{th}}(p_1;\beta _{\mathrm{th}})(p/p_1)^\alpha \hfill & ;p_1<p<p_2\hfill \\ 0\hfill & ;p_2<p\hfill \end{array}\right\},`$ (30)
where the normalization parameter $`C_\mathrm{e}(\beta _{\mathrm{th}},\alpha ,p_1,p_2)`$ is determined by $`𝑑pf_{\mathrm{e},\mathrm{th}\&\mathrm{cr}}(p;\beta _{\mathrm{th}},\alpha ,p_1,p_2)=1`$.
The integral in Eq. (16) has to be evaluated numerically for the thermal electron distribution. But the analytical solution for the power-law case can be expressed as
$`P_{\mathrm{cr}}(t;\alpha ,p_1,p_2)={\displaystyle \frac{\alpha 1}{p_1^{1\alpha }p_2^{1\alpha }}}{\displaystyle \frac{3}{16}}(1+t)\times `$
$`[\mathrm{B}_{\frac{1}{1+p^2}}({\displaystyle \frac{1+\alpha }{2}},{\displaystyle \frac{\alpha }{2}})`$
$`\mathrm{B}_{\frac{1}{1+p^2}}({\displaystyle \frac{3+\alpha }{2}},{\displaystyle \frac{2+\alpha }{2}}){\displaystyle \frac{7+3\alpha }{3+\alpha }}`$
$`\mathrm{B}_{\frac{1}{1+p^2}}({\displaystyle \frac{5+\alpha }{2}},{\displaystyle \frac{4+\alpha }{2}}){\displaystyle \frac{12+3\alpha }{5+\alpha }}`$ (31)
$`+p^{5\alpha }\{({\displaystyle \frac{3}{5+\alpha }}+{\displaystyle \frac{2p^2}{3+\alpha }})(2\mathrm{arcsinh}\left(p\right)|\mathrm{ln}\left(t\right)\left|\right)`$
$`+\left|{\displaystyle \frac{1t}{1+t}}\right|({\displaystyle \frac{1+t^2}{2\left(5+\alpha \right)t}}+{\displaystyle \frac{5}{5+\alpha }}+{\displaystyle \frac{4p^2}{3+\alpha }}+{\displaystyle \frac{2p^4}{1+\alpha }})\}]_{\stackrel{~}{p}_1\left(t\right)}^{\stackrel{~}{p}_2\left(t\right)},`$
with $`\stackrel{~}{p}_1(t)=\mathrm{max}(p_1,\sqrt{t}/2)`$ and $`\stackrel{~}{p}_2(t)=\mathrm{max}(p_2,\sqrt{t}/2)`$. Here $`\mathrm{B}_x(a,b)`$ denotes the incomplete beta-function (Abramowitz & Stegun 1965), and
$$\left[f(p)\right]_{p_1}^{p_2}=f(p_2)f(p_1).$$
(32)
$`P_{\mathrm{cr}}(t;\alpha ,p_1,p_2)=0`$ if $`|\mathrm{ln}(t)|>2\mathrm{arcsinh}(p_2)`$. This function is plotted in Fig. 3 and 4. The resulting change in a black-body photon spectrum due to Comptonization by CR-electrons is shown in Fig. 1 as the curves labeled with $`ji`$. Note that for a given frequency the influence of the gain in up-scattered photons of lower initial frequencies, $`j_{\mathrm{cr}}(x)`$, on the spectral distortions are negligible compared to $`i(x)`$, which describes the loss of photons from the CMB of intially this frequency, as long as $`x<10`$ and the lower cutoff the electron spectrum $`p_1>3`$. Therefore the rSZ effect can be well approximated as a pure absorption process of photons in the CMB-spectral range ($`x<10`$) for sufficiently relativistic electrons, justifying Eq. (1).
### 2.3 Relativistic SZ-Effect in Emission
Above the crossover frequency ($`x_{\mathrm{cross}}3.83`$) the arriving up-scattered photons outnumber the disappearing ones giving an SZ increment. This has a different spectral shape in the relativistic case compared with the classical one. The higher the frequency, the stronger is the rSZ effect compared to the thSZ effect. This provides for a promising opportunity to detect rSZ effect via this characteristic signature. Fig. 57 illustrate this. A direct application of the emission signature of the rSZ effect can be found in Sect. 4.2.
## 3 Radio Galaxies
### 3.1 Remnant Radio Plasma
The presence of radio plasma produces all three SZ effects discussed in this paper: The expanding lobes of radio galaxies (RGs) and radio-loud quasars compress and may shock the ambient IGM, thereby increasing its thSZ effect. During the inflation phase of the radio lobes the ambient medium is pushed aside, leading to a possible kSZ effect. And finally the relativistic electron content of the radio plasma produces the rSZ effect.
The inflation phase of a RG is short compared to cosmological times, thus we do not consider the transient kinematic effect. Furthermore, since gas on both sides of the radio lobes is accelerated in different directions, giving contributions of opposite sign but similar strength which should therefore roughly cancel out. In Sect. 3.3 we argue that the combined kSZ effect of an isotropic ensemble of flows is practically indistinguishable from a thSZ effect with the same energy content, further justifying our simplified treatment.
Radio lobes become unobservable due to adiabatic, IC- and synchrotron energy losses of the electrons on a time-scale shorter than the expansion time of the radio lobes which leads to pressure equilibrium with the surrounding medium (e.g. Kaiser et al. 2000). However, the energy transfered to the IGM by the expansion of the lobes, and also the relativistic electron population remain and leave their fingerprints in the CMB spectrum via Comptonization.
The radio plasma and later radio ghosts will expand or contract until they reach pressure equilibrium with the surrounding medium. The pressure of the radio ghost is given by that of the confined relativistic particles and the magnetic fields, assumed to be in rough energy equipartition. Thus the magnetic field energy densities should be of the order of the thermal energy density of the environment. Subsonic turbulence in this environment, which has an energy density below the thermal energy density, is therefore not strong enough to overcome the magnetic elastic forces of the radio ghost if the magnetic field in the ghost is contiguous on large scales. In the case of a magnetic field disjoined on scales smaller than the physical size of the ghost and/or the turbulence being sonic or super-sonic, which is e.g. expected in giant merger events of cluster of galaxies, it can ‘shred’ the ghost into smaller pieces. The size of such pieces will be comparable to the eddy size of the turbulence. This means, since a typical turbulent spectrum has less energy density on smaller scales, that there is a length-scale below which the turbulence is not able to overcome magnetic forces. Turbulent erosion of radio ghosts should stop at this length-scale, leaving small-scale patches of still unmixed old radio plasma. Similarly, fluid instabilities of Kelvin-Helmholtz or Rayleigh-Taylor type at the surface of the ghost will not lead to a complete mixing of the radio plasma with the surrounding matter.
The relativistic electrons suffer from adiabatic losses during the inflation phase of the radio plasma, afterwards only from synchrotron and IC- losses. Coulomb- and bremsstrahlung-losses are negligible in a tenuous relativistic plasma. The cooling of an ultra-relativistic electron with momentum $`p`$ is governed by
$$\frac{dp}{dt}=a_\mathrm{C}+a_\mathrm{b}p+a_\mathrm{s}p^2.$$
(33)
(Kardashev 1962), where we ignored adiabatic losses or gains due to volume changes. The coefficient $`a_\mathrm{C}`$, $`a_\mathrm{b}`$, and $`a_\mathrm{s}`$ for Coulomb, bremsstrahlung, and synchrotron/IC losses are (Rephaeli 1979, Blumenthal & Gould 1970)
$`a_\mathrm{C}`$ $`=`$ $`{\displaystyle \frac{3}{2}}\sigma _\mathrm{T}cn_\mathrm{e}\left(\mathrm{ln}{\displaystyle \frac{m_\mathrm{e}c^2p^{1/2}}{\mathrm{}\omega _\mathrm{p}}}+0.22\right),`$ (34)
$`a_\mathrm{b}`$ $`=`$ $`{\displaystyle \frac{3\alpha }{\pi }}\sigma _\mathrm{T}cn_\mathrm{e}\left(\mathrm{ln}2p{\displaystyle \frac{1}{3}}\right),`$ (35)
$`a_\mathrm{s}`$ $`=`$ $`{\displaystyle \frac{4}{3}}{\displaystyle \frac{\sigma _Tc}{m_\mathrm{e}c^2}}\left(\epsilon _\mathrm{B}+\epsilon _{\mathrm{cmb}}\right),`$ (36)
where the electron density of the background gas is $`n_\mathrm{e}`$, $`\epsilon _\mathrm{B}=B^2/(8\pi )`$ is the magnetic field and $`\epsilon _{\mathrm{cmb}}`$ the CMB photon energy density. The plasma frequency is $`\omega _\mathrm{p}=\sqrt{4\pi e^2n_\mathrm{e}/m_\mathrm{e}}`$ and $`\alpha `$ is the fine-structure constant. In order to judge the importance of the different terms in Eq. (33) we define the cooling times
$$t_\mathrm{C}=\frac{p}{a_C}=22\mathrm{Gyr}\left(1+0.01\mathrm{ln}(p_1/n_5)\right)^1p_1/n_5,$$
(37)
$$t_\mathrm{b}=\frac{1}{a_b}=8.610^3\mathrm{Gyr}\left(1+0.38\mathrm{ln}(p_1)\right)^1/n_5,$$
(38)
$$t_\mathrm{s}=\frac{1}{a_sp}=9.8\mathrm{Gyr}/(\epsilon _{11}p_1),$$
(39)
where $`p_1=p/10`$, $`n_5=n_\mathrm{e}/(10^5\mathrm{cm}^3)`$, and $`\epsilon _{11}=(\epsilon _\mathrm{B}+\epsilon _{\mathrm{cmb}})/(10^{11}\mathrm{erg}\mathrm{cm}^3)`$. We note, that in the following estimates the radio plasma is assumed to contain relativistic electrons only above $`p=10`$. From these timescales it becomes obvious that as long as the electron density within the radio plasma stays below $`10^5\mathrm{cm}^3`$ only synchrotron and IC losses have to be taken into account. If the radio plasma does not mix on a microscopic scale with the ambient medium this can be expected. In the following we assume that this is indeed the case. We note that in the case of a cold denser gas component in the radio plasma the Coulomb cooling heats the cold electron and this heat contributes to the thSZ-effect. The evolution of the relativistic electron distribution function in this case can be calculated analytically for constant conditions in the plasma ($`B,n_\mathrm{e},\epsilon _{\mathrm{cmb}}=const`$) (Enßlin et al. 1999). If Coulomb- and bremsstrahlung-cooling can be neglected Eq. (33) is solved by
$$p(t)=\left[\frac{1}{p_0}+\frac{4}{3}\frac{\sigma _Tc}{m_\mathrm{e}c^2}\left(\frac{B^2}{8\pi }+\epsilon _{\mathrm{cmb}}\right)t\right]^1,$$
(40)
where the brackets ($`\mathrm{}`$) indicate time-averaged quantities, and $`p_0`$ is the momentum at $`t=0`$. In order to demonstrate that even under disadvantageous conditions electrons stay relativistic we insert $`p_0=10`$, $`B^2=(20\mu \mathrm{G})^2`$, and $`\epsilon _{\mathrm{cmb}}=11\epsilon _{\mathrm{cmb},\mathrm{today}}`$ corresponding to an time-average over the CMB energy density between today and $`z=2`$ in an EdS Universe. The final electron momentum after $`t=10`$ Gyr cooling is still $`p>3`$. Thus, practically all relativistic electrons in old radio plasma stay relativistic for cosmological times. Their energies change, but this is not important for the rSZ effect in the spectral region of the CMB: the rSZ decrement depends practically only on the number of relativistic electrons, and not on their energy. This is because all the up-scattered photons have energies far above the CMB range, so that their actual energy does not influence measurements within this range ($`x<10`$). However, a search for up-scattered photons is a promising way to detect the rSZ-effect of very low energy relativistic electrons.
### 3.2 Detectability of Radio Ghosts
Beside the possibility to see a rSZ signature, discussed in this article, there are a few other possibilities to detect radio ghosts.
Radio ghosts are practically unobservable as long as their electron population remains at low energies. But if the population is re-accelerated the ghost becomes radio luminous again. This can happen when the ghost is dragged into a large-scale shock wave, e.g. in a merger event of clusters of galaxies or at the accretion shock where the matter is falling onto a cluster. The emission region is expected to be irregularly shaped, and should exhibit linear polarization due to the compression of the magnetic fields in the shock. Such regions are indeed observed in the periphery of a few clusters of galaxies and are called ‘cluster radio relics’ (for reviews, see Jaffe 1992 and Feretti & Giovannini 1996). Their properties, such as degree and direction of polarization, surface luminosity, peripheral position etc., can be understood within such a scenario (Enßlin et al. 1998a). The observed rarity of the cluster radio relic phenomena can be understood in this context: the presence of a shock wave, which should be quite frequently found in and near clusters of galaxies, is not sufficient to produce a cluster radio relic, a radio ghost has to be present at the same location, too.
Another way to detect the presence of radio ghosts is via their ability to magnetically scatter ultra-high energy (UHE) cosmic rays (CR) (Enßlin 1999). Without such scattering the distribution of sky arrival directions of UHE CR should trace their source distribution within $`50`$ Mpc, due to the limited distance protons above $`310^{19}`$ eV can travel without strong photo-pion energy losses (Greisen 1966, Zatsepin & Kuzmin 1966). If the distribution of UHE CR sources follows that of the matter in the local Universe a non-uniform arrival direction distribution is expected, contrary to the observations. Medina Tanco & Enßlin (in preparation) show that under optimistic assumptions about their distribution radio ghosts can sufficiently isotropize the UHE CR arrival directions. Since the UHE CR scattering angle decreases with particle energy, this scenario can in principle be tested, as soon as sufficient CR data is available.
### 3.3 The Energy Budget of Radio Ghosts
The energy a radio galaxy release into a radio ghost and its environment during its lifetime is
$$E_{\mathrm{gh}}=\epsilon _{\mathrm{gh}}V_{\mathrm{gh}}+PV_{\mathrm{gh}}$$
(41)
where $`V_{\mathrm{gh}}`$ is the volume occupied by the ghost, $`P`$ the pressure inside and outside the ghost and
$$\epsilon _{\mathrm{gh}}=\epsilon _\mathrm{e}+\epsilon _\mathrm{p}+\epsilon _B=(1+k_p+k_B)\epsilon _\mathrm{e}$$
(42)
is the energy density in the remnant radio plasma, split into its leptonic ($`\epsilon _\mathrm{e}`$), hadronic ($`\epsilon _\mathrm{p}`$) and magnetic ($`\epsilon _B`$) part. Unfortunately, the ratios between these energy densities are still undetermined. The parameters $`k_p=\epsilon _p/\epsilon _e`$ and $`k_B=\epsilon _B/\epsilon _e`$ take account of this.
The radio plasma is expected to follow a relativistic equation of state ($`P=\frac{1}{3}\epsilon _{\mathrm{gh}}`$; note that isotropically oriented magnetic fields also follow a relativistic equation of state). The fraction of $`E_{\mathrm{gh}}`$ stored as volume work in the ambient medium is therefore
$$\mathrm{\Delta }E_{\mathrm{th}}=PV_{\mathrm{gh}}=\frac{1}{4}E_{\mathrm{gh}}.$$
(43)
This thermal energy produces a thSZ effect in addition to that expected from the heat in the IGM due to structure formation. It seems to be difficult to follow this energy in an expanding Universe with ongoing structure formation. But it can be argued that this energy is released into gravitationally bound structures (clusters and filaments of galaxies). Filaments expand due to the Hubble flow, but the matter within a given filament flows into the next galaxy cluster and gets compressed thereby. This means, that on a longer time-scale the extra heat from RG is expected to be confined in clusters of galaxies and does not suffer from adiabatic losses.
In fact, part of the energy injected into the environment might also be stored gravitationally whenever gas is pushed to a higher gravitational potential or transformed to kinetic energy of gas flows. Gravitational energy is converted to kinetic energy, whenever the accelerated gas flows into the next potential well. From Eqs. (1) and (6) one would conclude that the kSZ effect is zero for an average of isotropically oriented flows. But Eq. (1) is only an approximation. A better description of the average kSZ effect for an isotropic ensemble of flows can be gained by the following argument: the individual electron velocities of all flows measured in the CMB-rest frame can be combined into a single momentum space distribution function, which is of course isotropic. This function can be used in order to calculate the Comptonization with the formalism described in Sect. 2.2, since the latter does not depend on the position of an electron along the line of sight in the optically thin limit. The kinetic energies are non-relativistic, so that the spectral changes are well described by the function $`g(x)`$ times the total kinetic energy of the electrons (see Fig. 6). Thus the Comptonization is independent of the exact spectral shape of the electron spectrum or the distribution of flow velocities. Therefore we assume that all of the energy, which the inflating radio lobes give to their environment, contributes to the thSZ effect. A similar argumentation for gaussian, isotropic, random velocity fields was given in Zeldovich et al. (1972).
We approximate the relativistic electron spectrum by a single power-law (Eq. (26)), so that the energy density $`\epsilon _\mathrm{e}`$ and the number density $`n_{\mathrm{e},\mathrm{cr}}`$ of the relativistic electron of a ghost are related via the average kinetic energy of the electrons ($`\gamma _\mathrm{e}1m_\mathrm{e}c^2`$):
$$\epsilon _\mathrm{e}=n_{\mathrm{e},\mathrm{cr}}\gamma _\mathrm{e}1m_\mathrm{e}c^2\frac{\alpha 1}{\alpha 2}n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2p_1.$$
(44)
The approximation assumes the particles to be ultra-relativistic ($`1p_1<p_2`$) and the distribution to be dominated by the lower end ($`p_1p_2`$ and $`\alpha >2`$). In the case that the distribution has a trans- or non-relativistic part the exact formulae for energy density, pressure and adiabatic index can be found in Appendix A.
If the ultra-relativistic approximation is valid and the low energy electrons dominate the spectrum, we can write:
$$n_{\mathrm{e},\mathrm{cr}}=\frac{\alpha 2}{\alpha 1}\frac{\epsilon _{\mathrm{gh}}}{(1+k_\mathrm{p}+k_B)m_\mathrm{e}c^2p_1}.$$
(45)
All the quantities on the rhs have to be inserted for the same stage of the radio plasma evolution in order to give a consistent result. The best observational constraints for these quantities are of course given for the moment of injection, when the radio plasma is visible.
### 3.4 CMB-Distortions
We estimate now the Comptonization due to energy from radio galaxies stored in different forms into the IGM. We have to calculate the integrals in Eq. (4) and (7) in a cosmological context. We need to know the distribution of injected heat and relativistic electrons. Since both quantities (injected thermal heat in the IGM and number of relativistic electrons in the radio plasma) are accumulative in the sense, that radio galaxies produce them, but cooling mechanisms are insufficient to remove them. For the thermal energy this is clear, since except for cooling flows in the very center of some cluster of galaxies the cooling time exceeds the Hubble time. For the number density of relativistic electrons this follows from the fact that the dominant cooling mechanism in a tenuous relativistic plasma are synchrotron- and IC-cooling, which become inefficient for low energy electrons (see Sec 3.1).
The integrands in Eq. (4) and (7) are proportional to the amount of remnant radio plasma at a given epoch (e.g. measured by the number density of relativistic electrons $`n_{\mathrm{e},\mathrm{rel}}`$). Since, as we argued, radio plasma is conserved, its electron density is given by
$$n_{\mathrm{e},\mathrm{rel}}(t)=_0^t𝑑t^{}\dot{n}_{\mathrm{e},\mathrm{rel}}(t^{}),$$
(46)
where $`\dot{n}_{\mathrm{e},\mathrm{rel}}(t)`$ is the source electron density of radio plasma at time $`t`$ in comoving coordinates. The optical depth of relativistic plasma can then be written as
$$\tau _{\mathrm{rel}}=\sigma _\mathrm{T}c_0^{z_{\mathrm{max}}}𝑑z\frac{dt}{dz}\dot{n}_{\mathrm{e},\mathrm{rel}}(z)_0^z𝑑z^{}\frac{dt}{dz^{}}(1+z^{})^3,$$
(47)
where $`t(z)`$ is the time between today and redshift $`z`$, and $`z_{\mathrm{max}}`$ is the maximal redshift of injection. We use $`z_{\mathrm{max}}=4`$, which is sufficiently low so that cooling effects can be neglected. Since $`\dot{n}_{\mathrm{e},\mathrm{rel}}(z)`$, the source density of relativistic electrons, is only poorly constrained observationally, we use $`Q_{\mathrm{jet}}(t)`$, the source density of energy from RG instead. For convenience, we define
$$\omega _{\mathrm{gh}}=\frac{\sigma _\mathrm{T}c}{m_\mathrm{e}c^2}_0^{z_{\mathrm{max}}}𝑑z\frac{dt}{dz}\dot{Q}_{\mathrm{jet}}(z)_0^z𝑑z^{}\frac{dt}{dz^{}}(1+z^{})^3,$$
(48)
and the total amount of energy release through radio plasma:
$$\overline{ϵ}_{\mathrm{gh}}=_0^{z_{\mathrm{max}}}𝑑z\frac{dt}{dz}\dot{Q}_{\mathrm{jet}}(z).$$
(49)
For an ultra-relativistic and steep electron population inside ghosts the approximation Eq. (45) gives
$$\tau _{\mathrm{gh}}=\frac{3}{4}\frac{\alpha 2}{\alpha 1}\frac{\omega _{\mathrm{gh}}}{(1+k_\mathrm{p}+k_B)p_1}.$$
(50)
A similar expression gives the additional thSZ effect due to IGM heating by expanding radio lobes:
$$y_{\mathrm{gh}}=\frac{\omega _{\mathrm{gh}}}{12}.$$
(51)
This follows from Eq. (43) if we assume that the IGM protons get half of the heat energy.
The optical depth of radio ghosts can be written as
$$\tau _{\mathrm{gh}}=\frac{\alpha 2}{\alpha 1}\frac{9y_{\mathrm{gh}}/p_1}{(1+k_\mathrm{p}+k_B)}=\frac{y_{\mathrm{gh}}}{10}\frac{3}{1+k_\mathrm{p}+k_B}\frac{10}{p_1}.$$
(52)
Here, we inserted $`\alpha =2.5`$. If we assume protons, electrons and magnetic fields to have all the same energy density, and insert a speculative lower electron cutoff of $`p_1=10`$, this equation suggests that $`y_{\mathrm{gh}}\tau _{\mathrm{gh}}`$. This means that the thSZ decrement due to the thermal IGM heating by inflating radio lobes always dominates over the rSZ decrement of the relativistic electron population in these radio lobes since $`|g(x)||i(x)|`$ in the Rayleight-Jeans regime (see Fig. 1 and Eq. (1) for details). The thSZ increment in the Wien-regime decreases exponential and therefore much faster than the rSZ increment of a power-law electron population. Therefore the rSZ effect of radio plasma dominates over its induced thSZ effect at sufficiently high frequencies as long as the electron population stays relativistic.
The total jet power of radio galaxies per comoving volume can be derived from the radio luminosity function of RG with the additional assumption that there is a unique relation between radio luminosity and jet power. This is not strictly correct, since it is known that for the largest part of the observed lifetime of RGs they should exhibit some luminosity evolution, even in models with constant jet power (for the case of powerful RG of type FRII see e.g. Begelman & Cioffi (1989), Falle (1991), Kaiser & Alexander (1997), Kaiser et al. (1997), Daly (1999)). But this evolution is ‘modest’ (roughly an order of magnitude) and its influence should produce some scatter in an empirical derived relation, but not a large systematic effect. Similarly we neglect the influence of the density of the gas in the radio source environments on its radio luminosity. Again we expect this only to increase the scatter about the empirical function we will use in the following to relate radio luminosity and jet power.
Enßlin et al. (1997) have fit a power-law to the jet power – radio luminosity relation derived by Rawlings & Saunders (1991) for a sample of radio galaxies, which includes both, the most powerful radio galaxies of type FRII and also the somewhat less luminous FRI objects. Their results are based on minimum energy arguments and age estimates. The real energy of radio plasma can easily be much higher than the minimal energy estimate by some factor $`f_{\mathrm{power}}>1`$ due to the presence of relativistic protons, low energy electrons or deviations from equipartition between particle and field energy densities. A rough estimate of $`f_{\mathrm{power}}`$ can be derived from observations of radio lobes embedded in the intra-cluster medium (ICM) of clusters of galaxies. These observations show a discrepancy of the thermal ICM-pressure to the pressure in the radio plasma following minimal energy arguments of a factor of $`510`$, even if projection effects are taken into account (e.g. Feretti et al. 1992). Since also a filling factor smaller than unity of the radio plasma in the radio lobes can mimic a higher energy density we chose $`f_{\mathrm{power}}=3`$ in order to be conservative.
The jet power – radio luminosity correlation at $`\nu =2.7`$ GHz is
$$q_{\mathrm{jet}}(L_\nu )=a_\nu (L_\nu /(\mathrm{Watt}\mathrm{Hz}^1h_{50}^2))^{b_\nu }f_{\mathrm{power}}$$
(53)
for which $`b_\nu =0.82\pm 0.07`$, $`\mathrm{log}_{10}(a_\nu /\mathrm{erg}\mathrm{s}^1h_{50}^2)=45.2826.22b_\nu \pm 0.18`$ (Enßlin et al. 1997). This relation allows to integrate the observed radio luminosity function $`n_\nu (L_\nu ,t)`$ in order to get the total jet power
$$Q_{\mathrm{jet}}(z)=_{L_{\mathrm{min}}}^{L_{\mathrm{max}}}𝑑L_\nu n_\nu (L_\nu ,z)q_{\mathrm{jet}}(L_\nu ).$$
(54)
We use $`L_{\mathrm{min}}=10^{23}\mathrm{Watt}\mathrm{Hz}^1`$, which is just above the region of the RLF which is dominated by starburst galaxies, and $`L_{\mathrm{max}}=10^{29}\mathrm{Watt}\mathrm{Hz}^1`$. We adopt different radio-luminosity functions parameterized by Dunlop & Peacock (1990) and integrate Eq. (54) in an EdS-cosmology up to redshift $`z_{\mathrm{max}}=4`$. Strictly, the expressions for $`n_\nu (L_\nu ,t)`$ used here only apply to those members of the radio source population with steep radio spectra (spectral index $`0.7`$ or less). However, the space density of the flat spectrum population, which is thought to consist mainly of FRI-type objects, is at least a factor 10 smaller at any given redshift (Dunlop & Peacock 1990). Therefore we can safely neglect their contribution to the overall radio luminosity function. We further calculate the ghost distribution function for $`b_\nu =0.7,0.82,1`$ in order to show the dependence on the uncertainties.
Results of the integration of the different RLFs and jet power-radio luminosity correlations are given in Tab. 1. The energy input in form of radio plasma is roughly $`\overline{ϵ}_{\mathrm{gh}}=310^{66}\mathrm{erg}\mathrm{Gpc}^3(f_{\mathrm{power}}/3)`$.
Note that the X-ray background, which is believed to be dominated by AGN emission, corresponds to an injection energy of $`10^{6768}\mathrm{erg}\mathrm{Gpc}^3`$ in comoving coordinates (Soltan 1982; Chokshi & Turner 1992; Fabian & Iwasawa 1999; Fabian 1999). Either the X-ray energy losses of AGNs exceed the radio plasma release, or if X-ray power is comparable to the jetpower then the latter is strongly underestimated here and $`f_{\mathrm{power}}30`$ would be more realistic. The resulting Comptonization would be comparable to the present-day upper limit of $`y<1.510^5`$ (Fixsen et al. 1996). The extra-galactic radio background has an energy density of $`510^{63}\mathrm{erg}\mathrm{Gpc}^3`$ (Longair & Sunyaev 1971), which is several orders of magnitude below the energy density in the X-ray background and the expected ghost energy density. This indicates that radio emission is a very inefficient mechanism to extract energy from radio plasma.
It is interesting to note that the X-ray background predicts a mass density in AGN black holes of $`(1.4\mathrm{}2.2)10^{14}M_{}\mathrm{Gpc}^3`$ for a mass-to-light-conversion-efficiency of AGNs of 0.1. This is consistent with the central black hole mass density derived from the observed black-hole-to-galactic-bulge-mass ratio and the observed masses of galaxies (Faber et al. 1997).
If the total energy of the radio lobes would be thermalized, then no rSZ-effect would result, but the contribution to the thSZ effect would increase by a factor of 4, shifting it closer to the present upper limit.
### 3.5 Discussion
The expected thermal Comptonization due to heating of the environment of RG is of the order of $`y=1.410^6(f_{\mathrm{power}}/3)`$, if it is assumed that radio ghosts do not thermalize their internal energy, otherwise up to a factor of 4 higher. $`f_{\mathrm{power}}`$ is the ratio of true energy over equipartition energy in radio lobes.
There are two attempts in the literature to estimate the thSZ effect due to energy release from RGs.
Enßlin et al. (1998b) give $`y1.010^5\eta _{\mathrm{jet}/\mathrm{X}\mathrm{ray}}`$. For this estimate they assume that the jet power is completely thermalized, and that the energy budget is the same as the X-ray energy release ($`\eta _{\mathrm{jet}/\mathrm{X}\mathrm{ray}}1`$, thus 10 times higher than what we estimate here). It was also assumed that the thermal energy suffers from adiabatic cooling due to the Hubble flow, but since radio galaxies are located in filaments and clusters of galaxies only 1-dimensional instead of 3-dimentional expansion was taken into account.
Yamada et al. (1999) use a model for heating and cooling of the environment of radio lobes. They estimate the jet power from the black-hole-to-galaxy mass ratio ($`f_{\mathrm{bh}}=0.002`$) using a Press-Schechter description of galaxy (and black hole) growth and assuming a constant fraction ($`f_{\mathrm{rg}}=0.01)`$ of galaxies above $`10^{12}M_{}`$ to be active RGs. Since the total energy release in their description should be higher by an order of magnitude, compared to what we estimate here, it is not surprising that they get $`y=5.710^5(f_{\mathrm{bh}}/0.002)(f_\mathrm{r}/0.01)`$, which slightly violates the present day upper limit.
The present day limit on the Comptonization parameter gives already an important constraint on the energy budget of RGs, which will be improved in the near future. If it would be possible to detect the spectral signature of a rSZ effect, a further major step in the investigation of radio plasma could be achieved. We do not give detailed prediction for the rSZ increment, since different to the rSZ decrement at lower frequencies it strongly depends on the unknown shape of the electron distribution. But observations of the rSZ effect in emission could allow a spectral examination of the low-energy end of the relativistic electron content of the Universe.
## 4 Clusters of Galaxies
### 4.1 Embedded Radio Plasma
Clusters of galaxies are known to be strong sources of the thSZ effect. Since the thermal gas causing the thSZ distortions can also be observed with X-ray satellites via its bremsstrahlung emission it is possible to compare the angular diameter with the true line-of-sight diameter of the cluster. This gives directly the Hubble parameter, $`H_0`$, assuming spherical symmetry at least in a statistical sense. It is therefore of principal interest to estimate the possible influence of relativistic plasma on the measured Comptonization.
Radio plasma was mostly produced by outflows from AGN in virialized cosmological structures such as filaments and clusters of galaxies. Subsequent flows into deeper gravitational potentials should have transported a significant fraction of all radio plasma into clusters of galaxies. If the radio plasma is still unmixed with the thermal medium, it resembles a cavity in the X-ray emitting ICM gas. We assume in the following that a (for simplicity constant) fraction $`\mathrm{\Phi }_{\mathrm{gh}}`$ of the cluster volume is occupied by radio ghosts. The X-ray luminosity of this cluster is thus given by
$$L_X=a_X𝑑Vn_\mathrm{e}^2kT_\mathrm{e}^{\frac{1}{2}}(1\mathrm{\Phi }_{\mathrm{gh}})l_{\mathrm{cl}}^3n_{\mathrm{e},0}^2(1\mathrm{\Phi }_{\mathrm{gh}}).$$
(55)
$`n_\mathrm{e}`$ is the cluster electron density, $`n_{\mathrm{e},0}`$ its central value, and $`l_{\mathrm{cl}}`$ the characteristic length scale of the cluster (e.g the core radius). Since the cluster temperature $`kT_\mathrm{e}`$ can be determined spectroscopically we assume it to be known. The thSZ effect
$$y_{\mathrm{cl}}=\frac{\sigma _T}{m_\mathrm{e}c^2}𝑑ln_\mathrm{e}kT_\mathrm{e}(1\mathrm{\Phi }_{\mathrm{gh}})l_{\mathrm{cl}}n_{\mathrm{e},0}(1\mathrm{\Phi }_{\mathrm{gh}})$$
(56)
in combination with Eq. (55) allows to measure the typical scale of the cluster:
$$l_{\mathrm{cl}}\frac{L_X}{y_{\mathrm{cl}}^2}(1\mathrm{\Phi }_{\mathrm{gh}})$$
(57)
Since the angular diameter of a cluster should (at least in a statistical average) be identical to $`l_{\mathrm{cl}}`$ the Hubble constant can be derived:
$$H_0l_{\mathrm{cl}}^1\frac{y_{\mathrm{cl}}^2}{L_X(1\mathrm{\Phi }_{\mathrm{gh}})}.$$
(58)
There are two ways the presence of radio ghosts can affect the determination of $`H_0`$, which shift it in opposite directions: First, if a significant fraction of the volume is filled with relativistic plasma, then the true $`H_0`$ is greater than the one derived under the assumption that $`\mathrm{\Phi }_{\mathrm{gh}}=0`$. Second, if the measured $`y_{\mathrm{obs}}`$ contains a significant contamination due to absorption by the rSZ effect the true $`y_{\mathrm{cl}}`$, and therefore also $`H_0`$, is lower.
The second effect is negligible for the following reason: The optical depth of the relativistic electrons of the ghosts can be written as
$$\tau _{\mathrm{gh}}=\frac{6\varphi _{\mathrm{gh}}y_{\mathrm{cl}}}{(1+k_B+k_p)\gamma _\mathrm{e}1(1\mathrm{\Phi }_{\mathrm{gh}})},$$
(59)
if we assume pressure equilibrium between ghosts and thermal plasma, and use Eqs. (44), (56), and 71. Optimistic values for the unknown parameters of the old radio plasma ($`1+k_B+k_p3`$, $`\gamma _\mathrm{e}110`$) give $`\tau _{\mathrm{gh}}0.2\mathrm{\Phi }_{\mathrm{gh}}y_{\mathrm{cl}}`$. The rSZ decrement is roughly $`\tau _{\mathrm{gh}}i(x)`$, or smaller. The observed, rSZ-contaminated $`y`$-parameter is therefore
$$y_{\mathrm{obs}}=y_{\mathrm{cl}}+\frac{i(x)}{g(x)}\tau _{\mathrm{gh}}y_{\mathrm{cl}}(1+0.06\varphi _{\mathrm{gh}}).$$
(60)
In the last approximation we assumed that the measurement is taken around $`x=h\nu /kT_{\mathrm{cmb}}2`$, where the thSZ absorption is at its maximum and $`i(x)/g(x)0.3`$ (see Fig. 1). Inserting $`y_{\mathrm{obs}}`$ instead of $`y_{\mathrm{cl}}`$ into Eq. (58) gives an systematic error of $`0.1\mathrm{\Phi }_{\mathrm{gh}}`$, which is an order of magnitude smaller than the error due to the factor $`1\mathrm{\Phi }_{\mathrm{gh}}`$ in the denominator of that equation. The error due to the rSZ effect could be larger, if the measurement is done closer to the crossover frequency at $`x=3.83`$. However, this value of $`x`$ is inconvenient for detecting the thSZ effect. Also if $`\gamma _\mathrm{e}11`$ the distortion would be significant, but this is extremely speculative.
We conclude that the presence of a significant fraction of the volume of the ICM filled by radio ghosts biases the determination of the Hubble constant to lower values, if this effect is not accounted for.
### 4.2 Non-Thermal Electrons in the ICM
In some clusters of galaxies a relativistic ICM population of electrons is visible due to their synchrotron emission, which produces the radio halos of clusters of galaxies. Also the recently detected extreme ultraviolet (EUV) excess (Lieu et al. 1996) and the high energy X-ray (HEX) excess (Fusco-Femiano et al. 1998, 1999) in the Coma cluster indicate the presence of non-thermal electrons.
Our knowledge about the slope of the electrons spectrum in the ICM is limited (see Enßlin & Biermann 1998 for an attempt to compile the electron spectrum in the Coma cluster). The number density and therefore the optical depth of the higher energy electrons is too small to detect the rSZ decrement produced by them. See Birkinshaw (1999) for a discussion of the rSZ decrement from radio emitting electrons and McKinnon et al. (1990) for an attempt to detect this decrement in radio galaxies. The rSZ effect emission from these electrons may already have been observed in the EUV (Hwang 1997, Enßlin & Biermann 1998, Sarazin & Lieu 1998, Bowyer & Berghöfer 1998) or in the HEX excess (Fusco-Femiano et al. 1998, 1999). The latter would require a high number of electrons in the range of $`35`$ GeV, which would over-produce radio emission if the magnetic fields are as high as Faraday rotation measurements indicate ($`210\mu `$G; Crusius-Waetzel et al. 1990; Kim et al. 1991; Feretti et al. 1995, 1999, Clarke et al. in preparation). Only if fields are as weak as $`0.16\mu `$G consistency between the expected and observed radio flux is established (Fusco-Femiano et al. 1998, 1999). As a possible solution of this apparent contradiction of magnetic field estimates Enßlin et al. (1999) and Sarazin & Kempner (2000) proposed that the HEX excess could alternatively be produced by bremsstrahlung of a non-thermal high energy tail of the thermal electron distribution. Such a tail could exist due to in-situ particle acceleration powered by ICM turbulence (Enßlin et al. 1999; Dogiel 1999; Blasi 2000).
Detection of the IC emission of such a non-thermal tail could confirm the bremsstrahlung-scenario. In the following we estimate the expected spectral distortions for the Coma cluster of galaxies.
The electron density of the cluster is assumed to follow a beta-profile:
$$n_\mathrm{e}(r)=n_{\mathrm{e},0}\left[1+(r/r_{\mathrm{cl}})^2\right]^{\frac{3}{2}\beta _{\mathrm{cl}}},$$
(61)
where $`r_{\mathrm{cl}}=400\mathrm{kpc}\mathrm{h}_{50}^1`$ is the core radius, $`\beta _{\mathrm{cl}}=0.8`$ is the beta-parameter, and $`n_{\mathrm{e},0}=2.8910^3\mathrm{cm}^3`$ is the central thermal electron density (Briel et al. 1992). We ignore the possible complication due to embedded old radio plasma discussed in Sect. 4.1 for simplicity.
The optical depth of a line of sight passing the cluster center at a distance $`R`$ is then
$$\tau _{\mathrm{cl}}(R)=\mathrm{B}(\frac{1}{2},\frac{3\beta _{\mathrm{cl}}1}{2})\frac{n_{\mathrm{e},0}r_{\mathrm{cl}}\sigma _T}{\left[1+(R/r_{\mathrm{cl}})^2\right]^{\frac{3}{2}\beta _{\mathrm{cl}}\frac{1}{2}}},$$
(62)
giving a central optical depth of $`\tau _{\mathrm{cl}}(0)=5.9510^3`$. The CMB-distortions are
$$\delta I_\nu =i_0\delta i(x)=i_0\tau _{\mathrm{cl}}(j(x)i(x))).$$
(63)
The electron spectrum is given for a thermal spectrum by Eq. (24) and for a modified thermal spectrum by Eq. (26). The non-thermal electrons are treated as an additional population, increasing slightly the total electron number, but leaving the number densities of electrons with $`p<p_1`$ unchanged. We estimate the distortions for a pure thermal spectrum (model 0) and for three modified thermal spectra, which are able to explain the HEX excess of Coma by bremsstrahlung (see Fig. 6 in Enßlin et al. 1999): model 1: $`\alpha =5.3`$, $`p_1=0.45`$, $`E_{\mathrm{kin},1}=50\mathrm{keV}`$; model 2: $`\alpha =2.9`$, $`p_1=0.51`$, $`E_{\mathrm{kin},1}=62.5\mathrm{keV}`$; model 3: $`\alpha =1.6`$, $`p_1=0.58`$, $`E_{\mathrm{kin},1}=80\mathrm{keV}`$; $`p_2=10`$ in all three models. The electron spectra are shown in Fig. 8 and the resulting CMB distortions at the center of the Coma cluster in Fig. 9. The differences between the thermal and the modified thermal models are considerable and change as a function of frequency, so that future multichannel CMB telescopes such as high-altitude ground based observatories (e.g. at the ALMA site or in Antarctica), balloon-experiments, or marginally even the Planck satellite (see Appendix B) can discriminate between these models. A detection of this non-thermal SZ effect would prove in-situ acceleration processes acting in the ICM of Coma. It would confirm the bremsstrahlung origin of the HEX excess and therefore solve the discrepancy between Faraday and IC based magnetic field estimates in Coma.
Since the SZ effect does not depend on distance (as long as the instrument beam resolves the cluster core) many other clusters can be investigated for the presence of a non-thermal SZ effects with future CMB experiments.
## 5 Conclusion
We investigated the transrelativistic Thomson scattering of photons on a isotropic electron distribution in the optically thin limit using the photon redistribution kernel (Eq. (21)). We derived for the first time an analytic formula for the scattering by a population of electrons with a power law momentum distribution, Eq. (26). We demonstrated that relativistic electron populations produce a decrement in the cosmic microwave background, similar to that an absorber with an optical depth of the Thomson depth of the relativistic electrons would produce. Our formalism can be generalized to the optical finite case using the techniques described in Rephaeli (1995a,b) or Molnar & Birkinshaw (1999).
We applied this theory to radio galaxies and clusters of galaxies:
Although a single radio galaxy produces a negligible SZ decrement, the combined effect of several radio galaxies and their remnants might produce a detectable signal. In order to show this we estimated the total cosmological release in jet power of radio galaxies using the observed radio luminosity function converted by an empirical jet power-radio luminosity relation. It is roughly $`310^{66}\mathrm{erg}\mathrm{Gpc}^3(f_{\mathrm{power}}/3)`$, where $`f_{\mathrm{power}}`$ gives the poorly constrained ratio between true energy content of radio lobes and the minimum energy estimate. If completely thermalized this energy would lead to a Comptonization parameter of $`y610^6`$, close to the present day observational limit of $`y<1.510^5`$. We argued that roughly $`3/4`$ of the released energy remains as a relativistic plasma, which rapidly becomes unobservable after the activity of the galaxy stopped. Thus the Comptonization due to heating by radio plasma is expected to be $`y1.510^6`$.
Patches of old plasma, called ‘radio ghosts’, are expected to survive turbulent erosion unmixed with the ambient medium, so that they are able to retain a low energy relativistic electron population. The optical depth for Thomson scattering by radio ghosts is $`\tau _{\mathrm{gh}}10^7`$ for very optimistic assumptions about the low energy cutoff of the fresh electrons during injection. Otherwise $`\tau _{\mathrm{gh}}`$ is much lower .
If radio ghosts occupy a significant volume in clusters of galaxies, they affect the determination of the Hubble constant via SZ and X-ray measurements. The geometric effect of the cavities formed by ghosts in the intra-cluster gas overwhelms the SZ decrement expected due to up-scattering of CMB photons by the ghosts’ relativistic electron populations. SZ based $`H_0`$ determinations have to be corrected to higher values due to this effect.
Finally, we demonstrated that future CMB telescopes such as the Planck satellite seemed to be useful tools to measure supra-thermal electron populations in clusters of galaxies. These are expected in the case of turbulent in-situ particle acceleration, and supported by the recently detected high energy X-ray excess in the Coma and Abell 2556 cluster. Such a detection would be crucial for our understanding of particle acceleration processes in clusters.
###### Acknowledgements.
We thank Rashid Sunyaev, Matthias Bartelmann, and Mark Birkinshaw, the referee, for comments on the manuscript.
## Appendix A Energy Density and Pressure of a Trans-Relativistic Power-Law Electron Distribution
If an electron population, with number density $`n_{\mathrm{e},\mathrm{cr}}`$, has a momentum spectrum which is a single power-law (Eq. (26)) the kinetic energy density $`\epsilon _\mathrm{e}`$ is
$`\epsilon _\mathrm{e}`$ $`=`$ $`n_{\mathrm{e},\mathrm{cr}}\gamma _\mathrm{e}1m_\mathrm{e}c^2`$ (64)
$`=`$ $`n_{\mathrm{e},\mathrm{cr}}{\displaystyle _0^{\mathrm{}}}𝑑pf_\mathrm{e}(p)(\sqrt{1+p^2}1)m_\mathrm{e}c^2`$ (66)
$`=`$ $`{\displaystyle \frac{n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2}{[p^{1\alpha }]_{p_2}^{p_1}}}[{\displaystyle \frac{1}{2}}\mathrm{B}_{\frac{1}{1+p^2}}({\displaystyle \frac{\alpha 2}{2}},{\displaystyle \frac{3\alpha }{2}})`$
$`+p^{1\alpha }(\sqrt{1+p^2}1)]_{p_2}^{p_1}`$
$``$ $`{\displaystyle \frac{\alpha 1}{\alpha 2}}{\displaystyle \frac{[p^{2\alpha }]_{p_2}^{p_1}}{[p^{1\alpha }]_{p_2}^{p_1}}}n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2`$ (67)
$``$ $`{\displaystyle \frac{\alpha 1}{\alpha 2}}n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2p_1`$ (68)
Ḣere, we have used the short-hand notation defined in Eq. (32). The first approximation assumes the particles to be ultra-relativistic ($`1p_1<p_2`$), and the second that the distribution is dominated by the lower end ($`p_1p_2`$ and $`\alpha >2`$). For an ultra-relativistic population of electrons (and similar for protons) the pressure is $`P_\mathrm{e}=\frac{1}{3}\epsilon _\mathrm{e}`$, but this is not correct for a transrelativistic population. There
$`P_\mathrm{e}`$ $`=`$ $`n_{\mathrm{e},\mathrm{cr}}{\displaystyle _0^{\mathrm{}}}𝑑pf_\mathrm{e}(p){\displaystyle \frac{1}{3}}pv(p)m_\mathrm{e}c`$ (69)
$`=`$ $`{\displaystyle \frac{n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2(\alpha 1)}{6[p^{1\alpha }]_{p_2}^{p_1}}}\left[\mathrm{B}_{\frac{1}{1+p^2}}({\displaystyle \frac{\alpha 2}{2}},{\displaystyle \frac{3\alpha }{2}})\right]_{p_2}^{p_1}`$ (70)
$``$ $`{\displaystyle \frac{1}{3}}{\displaystyle \frac{\alpha 1}{\alpha 2}}{\displaystyle \frac{[p^{2\alpha }]_{p_2}^{p_1}}{[p^{1\alpha }]_{p_2}^{p_1}}}n_{\mathrm{e},\mathrm{cr}}m_\mathrm{e}c^2{\displaystyle \frac{1}{3}}\epsilon _\mathrm{e}`$ (71)
The ultra-relativistic approximation was again applied, assuming $`1p_1<p_2`$. If this is not justified, correct results can easily be obtained from Eq. (66) and (70).
## Appendix B Detectability of the Cluster-rSZ-Effect with Multifrequency Experiments
In order to demonstrate that the predicted non-thermal CMB distortions in the clusters like the Coma cluster (Sect. 4.2) are detectable with multifrequency experiments we estimate the expected sensitivity of the Planck satelitte. The spectral distortions at the location of the observed cluster are produced by temperature fluctuations of the CMB on the scale of the cluster ($`\delta T/T10^6\mathrm{}10^5`$), the thSZ and kSZ effect of the ICM gas, and the rSZ effect of the non-thermal electron distribution. The resulting distortion are therefore
$$\delta i(x)=\underset{j=1}{\overset{3}{}}\tau _jf_j(x),$$
(72)
with
$`\tau _1=\tau _{\mathrm{th}},f_1(x)=j_{\mathrm{th}}(x)i(x),`$
$`\tau _2=\tau _{\mathrm{rel}},f_2(x)=j_{\mathrm{rel}}(x)i(x),`$
$`\tau _3=\delta T/T+\overline{\beta }_{\mathrm{gas}},f_3(x)=h(x).`$
Here, $`j_{\mathrm{rel}}(x)`$ are the spectral distorions produced by the electrons in the non-thermal tail with optical thickness $`\tau _{\mathrm{rel}}`$ and spectrum given by $`f_{\mathrm{e},\mathrm{th}\&\mathrm{cr}}(p;\beta _{\mathrm{th}},\alpha ,p_1,p_2)f_{\mathrm{e},\mathrm{th}}(p;\beta _{\mathrm{th}})`$ (with proper normalization).
The observing frequencies and sensitivities of the channels ($`k`$) of the Planck satelitte are taken from Puget (1998) and given in Tab. 2. We assume the beams to be Gaussian, and therefore to have the areas $`A_k=1.13\mathrm{FWHM}_k^2`$. The core region of the Coma cluster is $`A_{\mathrm{Coma}}\pi r_{\mathrm{core}}^2300\mathrm{arcmin}^2`$. Thus, for each frequency $`N_k=A_{\mathrm{Coma}}/A_k`$ independent beams probe the CMB. Using a $`\chi ^2`$ analysis allows to dissentangle the different spectral distortions if sufficient multichannel sensitivity is given. The accuracy in the determination of $`\tau _j`$ is given by $`\mathrm{\Delta }\tau _j=\sqrt{C_{jj}}`$, where $`C=A^1`$ and
$$A_{jl}=\underset{k}{}N_k\frac{f_j(x_k)f_l(x_k)}{(\mathrm{\Delta }i(x_k))^2}$$
(73)
(Press et al. 1992). Assuming $`\tau _{\mathrm{th}}=5.9510^3`$ and $`\tau _{\mathrm{rel}}`$ and $`j_{\mathrm{rel}}`$ according to the models 1,2, and 3 (see Sect. 4.2) we find that $`\tau _{\mathrm{rel}}/\mathrm{\Delta }\tau _{\mathrm{rel}}=`$ 0.61, 1.15, 0.96 correspondingly. This demonstrates that the Planck mission is expected to give a marginal (1-sigma) detection of such a supra-thermal electron population in the Coma cluster. Since some other clusters are also expected to contain supra-thermal electrons – e.g. Abell 2256 revealed very recently a similar HEX excess (Fusco-Femiano 2000) – a combined signal from several (HEX excess selected, or merger and post-merger) clusters might be detectable with statistical significance by Planck. This demonstrates the ability of future sensitive multichannel CMB telescopes – e.g. dedicated baloon experiments – to detect supra-thermal electron populations in clusters of galaxies. |
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