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import torch | |
import torch.nn.functional as F | |
from diffusers import AutoencoderKL, LMSDiscreteScheduler, UNet2DConditionModel | |
# For video display: | |
from PIL import Image | |
from torchvision import transforms as tfms | |
from tqdm.auto import tqdm | |
from transformers import CLIPTextModel, CLIPTokenizer, logging, AutoProcessor, CLIPVisionModel | |
import os, glob | |
from pathlib import Path | |
import gradio as gr | |
torch.manual_seed(1) | |
def pil_to_latent(input_im): | |
# Single image -> single latent in a batch (so size 1, 4, 64, 64) | |
with torch.no_grad(): | |
latent = vae.encode(tfms.ToTensor()(input_im).unsqueeze(0).to(torch_device)*2-1) # Note scaling | |
return 0.18215 * latent.latent_dist.sample() | |
def latents_to_pil(latents): | |
# bath of latents -> list of images | |
latents = (1 / 0.18215) * latents | |
with torch.no_grad(): | |
image = vae.decode(latents).sample | |
image = (image / 2 + 0.5).clamp(0, 1) | |
image = image.detach().cpu().permute(0, 2, 3, 1).numpy() | |
images = (image * 255).round().astype("uint8") | |
pil_images = [Image.fromarray(image) for image in images] | |
return pil_images | |
# Prep Scheduler | |
def set_timesteps(scheduler, num_inference_steps): | |
scheduler.set_timesteps(num_inference_steps) | |
scheduler.timesteps = scheduler.timesteps.to(torch.float32) # minor fix to ensure MPS compatibility, fixed in diffusers PR 3925 | |
def get_output_embeds(input_embeddings): | |
# CLIP's text model uses causal mask, so we prepare it here: | |
bsz, seq_len = input_embeddings.shape[:2] | |
causal_attention_mask = text_encoder.text_model._build_causal_attention_mask(bsz, seq_len, dtype=input_embeddings.dtype) | |
# Getting the output embeddings involves calling the model with passing output_hidden_states=True | |
# so that it doesn't just return the pooled final predictions: | |
encoder_outputs = text_encoder.text_model.encoder( | |
inputs_embeds=input_embeddings, | |
attention_mask=None, # We aren't using an attention mask so that can be None | |
causal_attention_mask=causal_attention_mask.to(torch_device), | |
output_attentions=None, | |
output_hidden_states=True, # We want the output embs not the final output | |
return_dict=None, | |
) | |
# We're interested in the output hidden state only | |
output = encoder_outputs[0] | |
# There is a final layer norm we need to pass these through | |
output = text_encoder.text_model.final_layer_norm(output) | |
# And now they're ready! | |
return output | |
# Defined a latent loss that is purely based on one image instead of multiple images | |
# This is used and working | |
# Need to look at if I can have the latents of multiple images merged together to give a thought | |
# Rather than a image - So that the thought is more on the design rather than pure image | |
def latent_loss(latent, conditioning_image): | |
# How far are the image embeds from lossembeds: | |
# image = Image.open(conditioning_image) | |
r_image = conditioning_image.resize((512,512)) | |
r_latent = pil_to_latent(r_image) | |
error = F.mse_loss(0.5*latent,0.5*r_latent) | |
return error | |
#Generating an image with these modified embeddings | |
def generate_with_embs(text_input, text_embeddings, conditioning_image): | |
height = 512 # default height of Stable Diffusion | |
width = 512 # default width of Stable Diffusion | |
num_inference_steps = 10 # Number of denoising steps | |
guidance_scale = 7.5 # Scale for classifier-free guidance | |
generator = torch.manual_seed(torch.seed()) # Seed generator to create the inital latent noise | |
batch_size = 1 | |
loss_scale = 100 #@param | |
imageCondSteps = 2 | |
max_length = text_input.input_ids.shape[-1] | |
uncond_input = tokenizer( | |
[""] * batch_size, padding="max_length", max_length=max_length, return_tensors="pt" | |
) | |
with torch.no_grad(): | |
uncond_embeddings = text_encoder(uncond_input.input_ids.to(torch_device))[0] | |
text_embeddings = torch.cat([uncond_embeddings, text_embeddings]) | |
# Prep Scheduler | |
set_timesteps(scheduler, num_inference_steps) | |
# Prep latents | |
latents = torch.randn( | |
(batch_size, unet.in_channels, height // 8, width // 8), | |
generator=generator, | |
) | |
latents = latents.to(torch_device) | |
latents = latents * scheduler.init_noise_sigma | |
# Loop | |
for i, t in tqdm(enumerate(scheduler.timesteps), total=len(scheduler.timesteps)): | |
# expand the latents if we are doing classifier-free guidance to avoid doing two forward passes. | |
latent_model_input = torch.cat([latents] * 2) | |
sigma = scheduler.sigmas[i] | |
latent_model_input = scheduler.scale_model_input(latent_model_input, t) | |
# predict the noise residual | |
with torch.no_grad(): | |
noise_pred = unet(latent_model_input, t, encoder_hidden_states=text_embeddings)["sample"] | |
# perform guidance | |
noise_pred_uncond, noise_pred_text = noise_pred.chunk(2) | |
noise_pred = noise_pred_uncond + guidance_scale * (noise_pred_text - noise_pred_uncond) | |
#### ADDITIONAL GUIDANCE ### | |
if conditioning_image: | |
if i%imageCondSteps == 0: | |
# Requires grad on the latents | |
latents = latents.detach().requires_grad_() | |
# Get the predicted x0: | |
latents_x0 = latents - sigma * noise_pred | |
# Decode to image space | |
# denoised_images = vae.decode((1 / 0.18215) * latents_x0).sample / 2 + 0.5 # range (0, 1) | |
# Calculate loss | |
# loss = flag_loss(denoised_images, lossEmbeds) * blue_loss_scale | |
loss = latent_loss(latents_x0, conditioning_image) * loss_scale | |
# loss = blue_loss(denoised_images) * blue_loss_scale | |
# print(i, 'loss item:', loss.item()) | |
# Get gradient | |
cond_grad = torch.autograd.grad(loss, latents)[0] | |
# print("cond_grad:", cond_grad) | |
# Modify the latents based on this gradient | |
latents = latents.detach() - cond_grad * sigma**2 | |
# compute the previous noisy sample x_t -> x_t-1 | |
latents = scheduler.step(noise_pred, t, latents).prev_sample | |
return latents_to_pil(latents)[0] | |
def getImageWithStyle(prompt, style_name, conditioning_image): | |
prompt = prompt + ' in the style of puppy' | |
style_loc = "styles/" + style_name + '/learned_embeds.bin' | |
style_embed = torch.load(style_loc) | |
# print(style_embed) | |
# print(style_embed.keys(), style_embed[style_embed.keys()[0]].shape) | |
new_style = list(style_embed.keys())[0] | |
# Tokenize | |
text_input = tokenizer(prompt, padding="max_length", max_length=tokenizer.model_max_length, truncation=True, return_tensors="pt") | |
# print("text_input:", text_input) | |
input_ids = text_input.input_ids.to(torch_device) | |
# print("Input Ids:", input_ids) | |
# print("Input Ids shape:", input_ids.shape) | |
token_emb_layer = text_encoder.text_model.embeddings.token_embedding | |
# Get token embeddings | |
token_embeddings = token_emb_layer(input_ids) | |
# print("Token Embeddings shape:", token_embeddings.shape) | |
# The new embedding - our special style word | |
replacement_token_embedding = style_embed[new_style].to(torch_device) | |
# Insert this into the token embeddings | |
token_embeddings[0, torch.where(input_ids[0]==6829)] = replacement_token_embedding.to(torch_device) | |
pos_emb_layer = text_encoder.text_model.embeddings.position_embedding | |
position_ids = text_encoder.text_model.embeddings.position_ids | |
# print("position_ids shape:", position_ids.shape) | |
position_embeddings = pos_emb_layer(position_ids) | |
# print("Position Embeddings shape:", token_embeddings.shape) | |
# Combine with pos embs | |
input_embeddings = token_embeddings + position_embeddings | |
# Feed through to get final output embs | |
modified_output_embeddings = get_output_embeds(input_embeddings) | |
# And generate an image with this: | |
image = generate_with_embs(text_input, modified_output_embeddings, conditioning_image) | |
return image | |
# name = "./Outputs/" + filename+".jpg" | |
# image.save(name) | |
# Supress some unnecessary warnings when loading the CLIPTextModel | |
logging.set_verbosity_error() | |
# Set device | |
torch_device = "cuda" if torch.cuda.is_available() else "mps" if torch.backends.mps.is_available() else "cpu" | |
if "mps" == torch_device: os.environ['PYTORCH_ENABLE_MPS_FALLBACK'] = "1" | |
# Load the autoencoder model which will be used to decode the latents into image space. | |
vae = AutoencoderKL.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="vae") | |
# Load the tokenizer and text encoder to tokenize and encode the text. | |
tokenizer = CLIPTokenizer.from_pretrained("openai/clip-vit-large-patch14") | |
text_encoder = CLIPTextModel.from_pretrained("openai/clip-vit-large-patch14") | |
image_encoder = CLIPVisionModel.from_pretrained("openai/clip-vit-large-patch14") | |
image_processor = AutoProcessor.from_pretrained("openai/clip-vit-large-patch14") | |
# The UNet model for generating the latents. | |
unet = UNet2DConditionModel.from_pretrained("CompVis/stable-diffusion-v1-4", subfolder="unet") | |
# The noise scheduler | |
scheduler = LMSDiscreteScheduler(beta_start=0.00085, beta_end=0.012, beta_schedule="scaled_linear", num_train_timesteps=1000) | |
# To the GPU we go! | |
vae = vae.to(torch_device) | |
text_encoder = text_encoder.to(torch_device) | |
unet = unet.to(torch_device) | |
image_encoder = image_encoder.to(torch_device) | |
# Used puppy as a placeholder here since the token is known | |
# Will replace with some other word that is better | |
# prompt = 'A farm' | |
# style_name = 'birb' | |
# conditioning_image_folder = './conditioning_images/' | |
# style_folder = './styles/' | |
# seedlist = [*range(0, 10000, 500)] | |
# print(seedlist) | |
# stylelist = ['birb', ] | |
# i = 0 | |
# for style_path in glob.glob(os.path.join(style_folder, '*')): | |
# seed = seedlist[i] | |
# i = i + 1 | |
# getImageWithStyle(prompt, style_path, None, seed) | |
# for conditioning_image in glob.glob(os.path.join(conditioning_image_folder, '*.jpg')): | |
# print("style_path:", style_path, "conditioning_image:", conditioning_image) | |
# getImageWithStyle(prompt, style_path, conditioning_image, seed) | |
def generateOutput(prompt, style_name, conditioning_image): | |
outputImage = getImageWithStyle(prompt, style_name, conditioning_image) | |
return outputImage | |
title = "Stable Diffusion SD Styles along with image conditioning" | |
description = "Shows the Stable Diffusion usage with SD Styles as well as ways to condition using different loss aspects" | |
examples = [["A farm", 'midjourney', 'conditioning_images/indianflag.jpg'],["A playground", 'lineart', None],["A theme park", 'cute-game-style', 'conditioning_images/vividcolors.jpg'],["A mouse", 'depthmap', 'conditioning_images/autumn.jpg']] | |
style_options = ['birb', 'moebius', 'midjourney', 'cute-game-style', 'depthmap', 'hitokomoru', 'lineart', 'madhubani'] | |
demo = gr.Interface( | |
generateOutput, | |
inputs = [ | |
gr.Textbox(), | |
gr.Dropdown(choices=style_options, label="Choose the Style you want"), | |
gr.Image(width=256, height=256, label="Image to use for Conditioning", type='pil'), | |
], | |
outputs = [ | |
gr.Image(width=512, height=512, label="Output"), | |
], | |
title = title, | |
description = description, | |
examples = examples, | |
cache_examples=False | |
) | |
demo.launch() | |